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warning/0506/math0506140.html | ar5iv | text | # Fermi Markov states
## 1. introduction
The quantum analogues of Markov processes were first constructed in , where the notion of quantum Markov chain on infinite tensor product algebras was introduced.
Nowadays, quantum Markov chains have become a standard computational tool in solid state physics, and several natural applications have emerged in quantum statistical mechanics and quantum information theory. On the other hand, the introduction in of the notion of โproduct stateโ on CAR algebras motivated the analogue construction in , of quantum Markov chains on these algebras as local perturbation of product states. The Fermi extension of product and Markov states is nontrivial because, even if the Fermion algebra is isomorphic to an infinite tensor products of matrix algebras, this embedding does not preserve the natural localization which plays an essential role in the very definition of these states. Product states describe non interacting (free or independent) systems. Markov states describe nearest neighbor interactions. In both cases the notion of โlocalizationโ plays a crucial role. Typically a discrete system is identified to a point in a graph. If this graph is not isomorphic to an interval in $``$ ($`1`$โdimensional case), one speaks of a random field. The crucial role of the localization is at the root of the difficulties to construct nontrivial examples of Markov fields. As the โinteracting degrees of freedomโ localized in a finite volume, increase with the volume, the first step to achieve this goal is to investigate the nonhomogeneous $`1`$โdimensional case, one of the goals of the present paper. Our second goal has to do with the most important difference between tensor and Fermi Markov chains, emerged from the analysis of , which in its turn is related to the difference between quantum Markov chains and quantum Markov states. The origin of this difference lies in the fact that, in the classical case, the simple structure of Markov states is equivalent to a single intrinsic condition: the Markov property. In the quantum case, while the Markov property can be formulated in terms of a localization property of the modular group of the state (see ), there is a class of states which have a Markov like local structure but do not necessarily enjoy the Markov property. This states are called quantum Markov chains (see e.g. , Definition 2.2). This phenomenon is related to the fact that the natural probabilistic extension of the notion of conditional expectation to the quantum case in general is not a projection (cf. ). On the other hand, the results of show that some of the most interesting physical applications involve precisely those Markov chains which are not Markov states. Such Markov chains can be explicitly constructed, but at the moment no intrinsic operator theoretic characterization is known. This distinction also appears in the Fermi case. However, the new phenomenon consists in the fact that, while in the tensor case all Markov states are convex combinations of states which are product states with respect to a new localization canonically associated to the original one (two block factors), this is not true in the Fermi case. The Fermi analogue of the convex combination of two block factors still appear. The last are called in the sequel strongly even Markov states. In addition, there is a completely new class of Fermi Markov states. This new class of non strongly even Markov states that appears in the classification theorem below (cf. Theorem 3.2) is likely to play in the Fermi case, the role played by the entangled states in the tensor product case.
One can define the notion of Fermi entanglement in analogy to the tensor product case. One can expect that the main problem of the entanglement theory, that is to find constructive and easily applicable criteria to discriminate entangled from non entangled states, will be in the Fermi case at least as difficult as in the tensor case. The first step to attack this problem, is to have a full and detailed description of this new class of states. In the paper only the simplest case was considered, that is when there is only one creator and annihilator in each site. The non homogeneous case, discussed in the present paper, includes as a particular case the translation invariant cases described in and its natural translation invariant generalization when there are $`d`$ creators and annihilators localized in each site. This leads to a much larger class of non strongly even Markov states, which can be completely described.
In the present paper, the program outlined below is carried out in the following steps. Section 2 contains the key result on the structure of the even transition expectations associated to Fermi Markov states (Proposition 2.7). It is then possible to provide the full classification of all even Markov quasi conditional expectations. Section 3 is devoted to the study of the most general situation, including also the non homogeneous Fermi Markov states. Even if the structure of the Markov states considered here is more complex than the one in , we are still able to provide their decomposition as direct integrals of minimal ones (Theorem 3.2). Furthermore, the minimal Markov states are the buildingโblocks for the construction of all the Fermi Markov states (Theorem 3.3). Section 4 deals with some general properties of the Markov states, such as the connection with the KMS boundary condition, and the entropy. Section 5 is devoted to the study of the detailed structure of the strongly even Markov states. They can be viewed as the Fermi analogue of the Ising type interactions. We show that a strongly even Markov state $`\phi `$ on the Fermion algebra arises by a lifting of a classical Markov process on the spectrum of a maximal Abelian subalgebra, with respect to the same localization as $`\phi `$. In addition, we establish the equality between the ConnesโNarnhoferโThirring dynamical entropy $`h_\phi (\alpha )`$ with respect the shift and the mean entropy $`s(\phi )`$. Section 6 provides the full list of translation invariant Fermi Markov states for low dimensional single site local algebras. The same method is applicable to higher dimensions. Finally, by using Moriya criterion, we show (cf. Proposition 6.1) that the Fermi Markov states which are not strongly even are indeed entangled, that is they provide a wide class of examples of entangled states on the CAR algebra which can be directly constructed and investigated in detail.
## 2. preliminaries
For the convenience of the reader, we collect some preliminary facts needed in the sequel.
### 2.1. Umegaki conditional expectations
By a (Umegaki) conditional expectation $`E:๐๐
๐`$ we mean a norm one projection of the $`C^{}`$โalgebra $`๐`$ onto the $`C^{}`$โsubalgebra (with the same identity $`1\mathrm{I}`$) $`๐
`$. The case of interest for us is when $`๐`$ is a full matrix algebra. Consider a set $`\{P_i\}`$ of central orthogonal projections of the range $`๐
`$ of $`E`$, summing up to the identity. We have
(2.1)
$$E(x)=\underset{i}{}E(P_ixP_i)P_i.$$
Then $`E`$ is uniquely determined by its values on the reduced algebras $`๐_{P_i}:=P_i๐P_i`$. When the above set $`\{P_i\}`$ consists of minimal projections, we get $`๐_{P_i}=N_i\overline{N}_i`$, and there exist states $`\varphi _i`$ on $`\overline{N}_i`$ such that
(2.2)
$$E(P_i(a\overline{a})P_i)=\varphi _i(\overline{a})P_i(a1\mathrm{I})P_i.$$
The reader is referred to for further details.
### 2.2. Quasi conditional expectations
Consider a triplet $`๐
๐`$ of unital $`C^{}`$โalgebras. A quasi conditional expectation w.r.t. the given triplet, is a completely positive, identity preserving linear map $`E:๐๐
`$ such that
$$E(ca)=cE(a),a๐,c.$$
Notice that, as the quasi conditional expectation $`E`$ is a real map, we have
$$E(ac)=E(a)c,a๐,c.$$
If $`\phi `$ is a normal faithful state on the $`W^{}`$โalgebra $`๐`$, the $`\phi `$โexpectation $`E^\phi :๐๐
`$ by Accardi and Cecchini preserving the restriction of $`\phi `$ to the $`W^{}`$โsubalgebra $`๐
`$, provides an example of quasi conditional expectation. Namely, it is enough to choose for $``$ any unital $`C^{}`$โsubalgebra of $`๐
`$ contained in the $`E^\phi `$โfixed point algebra. $`E^\phi `$ is a conditional expectation if and only if the modular group of $`\phi `$ leaves globally stable the subalgebra $`๐
`$, see .
### 2.3. The CAR algebra
Denote $`[a,b]:=abba`$, $`\{a,b\}:=ab+ba`$ the commutator and anticommutator between elements $`a`$, $`b`$ of an algebra, respectively.
Let $`J`$ be a set. The Canonical Anticommutation Relations (CAR for short) algebra over $`J`$ is the $`C^{}`$โalgebra $`๐_J`$ with the identity $`1\mathrm{I}`$ generated by the set $`\{a_j,a_j^+\}_{jI}`$, and the relations
$$(a_j)^{}=a_j^+,\{a_j^+,a_k\}=\delta _{jk}1\mathrm{I},\{a_j,a_k\}=\{a_j^+,a_k^+\}=0,j,kJ.$$
When there is no matter of confusion, we denote $`๐_J`$ simply as $`๐`$. The parity automorphism $`\mathrm{\Theta }`$, of $`๐`$ acts on the generators as
$$\mathrm{\Theta }(a_j)=a_j,\mathrm{\Theta }(a_j^+)=a_j^+,jJ,$$
and induces on $`๐`$ the $`_2`$โgrading $`๐=๐_+๐_{}`$ where
$$๐_+:=\{a๐|\mathrm{\Theta }(a)=a\},๐_{}:=\{a๐|\mathrm{\Theta }(a)=a\}.$$
Elements in $`๐_+`$ (resp. $`๐_{}`$) are called even (resp. odd).
A map $`T:๐^1๐^2`$ between the CAR algebras $`๐^1`$, $`๐^2`$ with $`_2`$โgradings $`\mathrm{\Theta }_1`$, $`\mathrm{\Theta }_2`$ is said to be even if it is gradingโequivariant:
$$T\mathrm{\Theta }_1=\mathrm{\Theta }_2T.$$
The previous definition applied to states $`\phi ๐ฎ(๐)`$ leads to $`\phi \mathrm{\Theta }=\phi `$, that is $`\phi `$ is even if it is $`\mathrm{\Theta }`$โinvariant.
Let the index set $`J`$ be countable, then the CAR algebra is isomorphic to the $`C^{}`$โinfinite tensor product of $`J`$โcopies of $`๐_2()`$:
(2.3)
$$๐_J\overline{\underset{J}{}๐_2()}^C^{}.$$
For the convenience of the reader, we report the JordanโKleinโWigner transformation establishing the mentioned isomorphism. Fix any enumeration $`j=1,2,\mathrm{}`$ of the set $`J`$. Let $`U_j:=a_ja_j^+a_j^+a_j`$, $`j=1,2,\mathrm{}`$ . Put $`V_0:=1\mathrm{I}`$, $`V_j:={\displaystyle \underset{n=1}{\overset{j}{}}}U_n`$, and denote
$`e_{11}(j):=a_ja_j^+,e_{12}(j):=V_{j1}a_j,`$
(2.4) $`e_{21}(j):=V_{j1}a_j^+,e_{22}(j):=a_j^+a_j.`$
$`\{e_{kl}(j)|k,l=1,2\}_{jI}`$ provides a system of commuting $`2\times 2`$ matrix units realizing the mentioned isomorphism.
Thanks to (2.3), $`๐_J`$ has a unique tracial state $`\tau `$ (at least when $`J`$ is countable) as the extension of the unique tracial state on $`๐_I`$, $`|I|<+\mathrm{}`$. Let $`J_1J`$ be a finite set and $`\phi ๐ฎ(๐)`$. Then there exists a unique positive element $`T`$ such that $`\phi _{๐_{J_1}}=\tau _{๐_{J_1}}(T)`$. The element $`T`$ is called the adjusted matrix of $`\phi _{๐_{J_1}}`$.<sup>1</sup><sup>1</sup>1For the standard applications to quantum statistical mechanics, we also use the density matrix w.r.t. the unnormalized trace, see Section 5. The state $`\phi _{๐_{J_1}}`$ is even (faithful) if and only if its adjusted matrix is even (invertible). The reader is referred to , Section XIV.1 and for further details.
We end the present subsection by recalling the description of product state (cf. ), and the definition of entanglement (cf. , Section 2). Let $`J_1,J_2I`$ with $`J_1J_2=\mathrm{}`$. Fix $`\phi _1๐ฎ(๐_{J_1})`$, $`\phi _2๐ฎ(๐_{J_2})`$. If at least one among them is even, then according to Theorem 11.2 of , the product state extension (called product state for short) $`\phi ๐ฎ(๐_{J_1{\scriptscriptstyle J_2}})`$ is uniquely defined. We write with an abuse of notation, $`\phi =\phi _1\phi _2`$. Suppose that $`J_1,J_2`$ are finite sets. Let $`T_1๐_{J_1}`$, $`T_2๐_{J_2}`$ be the adjusted densities relative to $`\phi _1๐ฎ(๐_{J_1})`$, $`\phi _2๐ฎ(๐_{J_2})`$, respectively. If at least one among $`T_1`$ and $`T_2`$ is even, then $`[T_1,T_2]=0`$ and $`T:=T_1T_2`$ is a well defined positive element of $`๐_{J_1J_2}`$ which is precisely the density matrix of $`\phi =\phi _1\phi _2`$. $`\phi ๐ฎ(๐_{J_1J_2})`$ is even if and only if $`\phi _1`$ and $`\phi _2`$ are both even.
A state $`\phi ๐ฎ(๐_{J_1J_2})`$ is called separable (w.r.t. to the decomposition $`๐_{J_1J_2}=\overline{๐_{J_1}๐_{J_2}}`$) if it is in the closed convex hull of all the product states over $`๐_{J_1J_2}`$. Otherwise it is called entangled.
### 2.4. Preliminaries on Fermi Markov states
Let us start as in , with a totally ordered countable set $`I`$ containing, possibly a smallest element $`j_{}`$ and/or a greatest element $`j_+`$. If $`I`$ contains neither $`j_{}`$, nor $`j_+`$, then $`I`$. If only $`j_+I`$, then $`I_{}`$, and if only $`j_{}I`$, then $`I_+`$. Finally, if both $`j_{}`$ and $`j_+`$ belong to $`I`$, then $`I`$ is a finite set and the analysis becomes easier. If $`I`$ is order isomorphic to $``$, $`_{}`$ or $`_+`$, we put simbolically $`j_{}`$ and/or $`j_+`$ equal to $`\mathrm{}`$ and/or $`+\mathrm{}`$ respectively. In such a way, the objects with indices $`j_{}`$ and $`j_+`$ will be missing in the computations.
Let $`๐_j`$ be the CAR algebra generated by $`d_j`$ creators and annihilators $`\{a_{j,1},a_{j,1}^+,a_{j,2},a_{j,2}^+,\mathrm{},a_{j,d_j},a_{j,d_j}^+\}`$ localized on the site $`jI`$. The numbers of the $`2d_j`$ generators of $`๐_j`$ may depend on $`j`$. We call
(2.5)
$$๐:=\overline{\underset{jI}{}๐_j}^C^{}$$
the Fermion algebra. Let $`J:={\displaystyle \underset{jI}{}}\{1,2,\mathrm{}d_j\}`$ be the disjoint union of the sets $`\{1,2,\mathrm{}d_j\}`$, $`jI`$. Then the Fermion algebra $`๐`$ given in (2.5) is nothing but the CAR algebra over the set $`J`$ previously described.
Now we pass to describe the local structure of the Fermion algebra $`๐`$. For each $`\mathrm{\Lambda }I`$, the local algebra $`๐_\mathrm{\Lambda }๐`$ is defined as $`๐_\mathrm{\Lambda }:=\overline{{\displaystyle \underset{j\mathrm{\Lambda }}{}}๐_j}`$. According to this notation, $`๐_{\{j\}}=๐_j`$ and $`๐_I=๐`$. Then $`\mathrm{\Lambda }I๐_\mathrm{\Lambda }๐`$ describes the local structure of the Fermion algebra. Particular subsets of $`I`$ are
$$[k,n]:=\{lI|kln\},n]:=\{lI|ln\}.$$
We put for $`\mathrm{\Lambda }I`$, $`S_\mathrm{\Lambda }:=S_{๐_\mathrm{\Lambda }}`$, $`S`$ being any map defined on $`๐`$. The reader is referred to , Section 2.6 and , Section 4 for further details.
A state $`\phi ๐`$ is said to be locally faithful if $`\phi _\mathrm{\Lambda }`$ is faithful whenever $`\mathrm{\Lambda }I`$ is finite.
If the number local generators $`d_j`$ depend on $`j`$ we refer to this situation as the nonhomogeneous case. Conversely, when $`I=`$ and $`d_j=d`$, $`j`$, the shift $`jj+1`$ acts in a natural way as an automorphism $`\alpha `$ of $`๐`$. A state $`\phi ๐ฎ(๐)`$ is translation invariant if $`\phi \alpha =\phi `$. If a state is translation invariant, then it is automatically even, see e.g. , Example 5.2.21.
We pass to the definition of Markov states which parallels Definition 4.1 of .
###### Definition 2.1.
An even state $`\phi `$ on $`๐`$ is called a Markov state if, for each $`n<j_+`$, there exists an even quasi conditional expectation $`E_n`$ w.r.t. the triplet $`๐_{n1]}๐_{n]}๐_{n+1]}`$ satisfying
(2.6)
$$\phi _{n]}E_n=\phi _{n+1]},$$
$$E_n(๐_{[n,n+1]})๐_{\{n\}}.$$
Notice that the local structure $`\mathrm{\Lambda }๐_\mathrm{\Lambda }`$, $`\mathrm{\Lambda }`$ finite subset of $`I`$, plays a crucial role in defining the Markov property. In fact, the isomorphism in (2.3) does not preserve neither the grading nor the natural localization.<sup>2</sup><sup>2</sup>2The algebra on the r.h.s. of (2.3) is naturally equipped with the trivial parity automorphism. Thus, its $`_2`$โgrading is trivial. Hence, it does not intertwine the corresponding Markov states.
When the numbers $`d_j`$ of the generators of $`๐_j`$ depend of the site, we call a Markov state $`\phi `$ (or equally well a Markov measure in the Abelian case) a nonhomogeneous Markov state. If $`d_j=d`$ for each $`j`$ we refer to the $`d`$โMarkov property. Thus, homogeneity means $`d`$โMarkov property for some $`d`$. For the applications to quantum statistical mechanics, $`d_j`$ is nothing but the โrange of interactionโ on the chain which might depend on the site, and when $`d=1`$ we are speaking of nearest neighbor interaction. The reader is referred to and the literature cited therein, for the connection between the Markov property and the statistical mechanics, and for further details.
Let $`\phi ๐ฎ(๐)`$ be a locally faithful Markov state. Then the restriction $`e_n:=E_n_{๐_{[n,n+1]}}`$ is a completely positive identity preserving linear map $`e_n:๐_{[n,n+1]}๐_{\{n\}}๐_{[n,n+1]}`$ leaving invariant the faithful state $`\phi _{[n,n+1]}`$. It is a quite standard fact (see e.g. ) that the ergodic average
$$\epsilon _n:=\underset{k}{lim}\frac{1}{k}\underset{h=0}{\overset{k1}{}}(e_n)^h$$
exists and defines a conditional expectation
$$\epsilon _n:๐_{[n,n+1]}(\epsilon _n)๐_{\{n\}}$$
projecting onto the fixed point algebra of $`e_n`$, the last coinciding with the range $`(\epsilon _n)`$ of $`\epsilon _n`$. The sequence $`\{\epsilon _n\}_{n<j_+}`$ of two point conditional expectations is called in the sequel the sequence of transition expectations associated to the locally faithful Markov state $`\phi `$. They uniquely determine, and are determined by the conditional expectations $`_n:๐_{n+1]}๐_{n]}`$, given for $`x๐_{n1]}`$, $`y๐_{[n,n+1]}`$ by
(2.7)
$$_n(xy)=x\epsilon _n(y).$$
In addition, it is quite standard to verify (cf. , Proposition 4.2) that we can freely replace the quasi conditional expectation $`E_n`$ in Definition 2.1 with its ergodic average $`_n`$. For the convenience of the reader we report Proposition 4.3 of .
###### Proposition 2.2.
Let $`f:๐_{[n,n+1]}(f)๐_{\{n\}}`$ be a even conditional expectation. The formula
$$(xy):=xf(y),x๐_{n1]},y๐_{[n,n+1]}$$
uniquely defines a even conditional expectation
$$:๐_{n+1]}๐_{n1]}(f)๐_{n]}.$$
From now on, we deal without further mention with even (quasi) conditional expectations. In addition, all the Markov states we deal with are even, and locally faithful if it is not otherwise specified.
###### Lemma 2.3.
Let $`:๐_{[k,l+1]}()๐_{[k,l]}`$ be a conditional expectation with $`๐_{[k,l1]}()`$. Then $``$ is faithful provided that $`_{๐_{[l,l+1]}}`$ is faithful.
###### Proof.
Let $`\phi _1`$, $`\psi `$ be faithful even states on $`๐_{[k,l1]}`$, $`(_{๐_{[l,l+1]}})`$ respectively. Put $`\phi _2:=\psi _{๐_{[l,l+1]}}`$. The product state $`\phi :=\phi _1\phi _2`$ is a faithful state on $`๐_{[k,l+1]}`$ left invariant by $``$. Fix $`a๐_{[k,l+1]}`$ with $`(a^{}a)=0`$. Then $`\phi (a^{}a)=\phi ((a^{}a))=0`$ which implies that $`a=0`$ as $`\phi `$ is faithful. Namely, $``$ is faithful. โ
We then pass to study the structure of the even conditional expectations
$$\epsilon _n:๐_{[n,n+1]}(\epsilon _n)๐_{\{n\}}.$$
To shorten the notations, it is enough to consider the case when $`n=0`$. After putting $`\epsilon :=\epsilon _0`$, let us start with the finite set $`\{P_j\}`$ of the minimal projections of the centre $`๐ต((\epsilon ))`$ of $`(\epsilon )`$.
###### Lemma 2.4.
The parity automorphism $`\mathrm{\Theta }`$ acts on $`๐ต((\epsilon ))`$, and the orbits of minimal projections consist of one or two elements.
###### Proof.
Let $`P_j`$ be a minimal projection of $`๐ต((\epsilon ))`$. As $`\epsilon `$ is even and $`\mathrm{\Theta }^2=id`$, we have that $`\mathrm{\Theta }(P_j)`$ is a minimal projection of $`๐ต((\epsilon ))`$. This means that either $`\mathrm{\Theta }(P_j)=P_j`$, or $`\mathrm{\Theta }(P_j)`$ is orthogonal to $`P_j`$. The latter means that the orbit of $`\mathrm{\Theta }(P_j)`$ consists of two elements. โ
We showed in that there are interesting examples with $`\mathrm{\Theta }(P_j)P_j`$. Let $`\epsilon `$ be as above. Some useful properties of the pieces $`\epsilon (PxP)P`$, $`P`$ being a even projection of the centre of $`(\epsilon )`$, minimal among the invariant ones, are described below.
###### Lemma 2.5.
Let $`M=M_+M_{}`$ be a $`_2`$โgraded full matrix algebra. If $`xM_{}`$ commutes with $`M_+`$, then $`x=0`$.
###### Proof.
Let the $`_2`$โgrading be implemented by the automorphism $`\mathrm{\Theta }`$. As $`\mathrm{\Theta }_M`$ is inner, there exists an even selfadjoint unitary $`VM`$, uniquely determined up to a sign, implementing $`\mathrm{\Theta }`$ on $`M`$, see , Corollary 8.11. This means that $`M_+=A^{}`$, $`A`$ being the Abelian algebra generated by $`V`$, and the commutant is taken in the full matrix algebra $`M`$. As $`x(M_+)^{}`$, $`xA^{\prime \prime }A`$. As $`x`$ is odd, we have $`VxV=x`$. Collecting together, we obtain $`x=0`$. โ
###### Lemma 2.6.
Let $`M=M_+M_{}`$ be a $`_2`$โgraded full matrix algebra. For every $`\mathrm{\Theta }`$โinvariant full matrix subalgebra $`NM`$, there exists a unique $`\mathrm{\Theta }`$โinvariant full matrix subalgebra $`\overline{N}M`$ such that $`N\overline{N}`$, and
(2.8)
$$x\overline{x}+\sigma (x,\overline{x})\overline{x}x=0,xN_\pm ,\overline{x}\overline{N}_\pm $$
where $`\sigma (x,\overline{x})`$ is 1 if both $`x`$, $`\overline{x}`$ are odd, and $`1`$ in the remaining cases. Moreover, we have $`N\overline{N}=1\mathrm{I}`$.
###### Proof.
Let $`\stackrel{~}{N}:=N^{}M`$ which is a $`\mathrm{\Theta }`$โinvariant full matrix subalgebra of $`M`$ as well. Fix a (even) unitary $`VN`$ uniquely determined up to a sign, implementing $`\mathrm{\Theta }`$ on $`N`$. Define for $`x=x_++x_{}\stackrel{~}{N}`$,
(2.9)
$$\beta (x):=x_++Vx_{}.$$
It is easy to see that $`\beta `$ defines a $``$โalgebra isomorphism between $`\stackrel{~}{N}`$ and $`\beta (\stackrel{~}{N})`$. Thus, the full matrix algebra $`\overline{N}:=\beta (\stackrel{~}{N})`$ is the algebra we are looking for.
For the uniqueness, let $`\overline{R}M`$ be a $`\mathrm{\Theta }`$โinvariant full matrix algebra fulfilling the commutation relations in (2.8) whenever $`xN_\pm `$, $`\overline{x}\overline{R}_\pm `$, such that $`N\overline{R}=M`$. Then it is easy to verify that $`\stackrel{~}{R}:=\left\{x_++Vx_{}\right|x\overline{R}\}`$ is a full matrix subalgebra of $`\stackrel{~}{N}`$. Since $`M=N\stackrel{~}{N}N\stackrel{~}{N}`$, we get that $`\stackrel{~}{R}`$ must coincide with $`\stackrel{~}{N}`$ which implies $`\overline{R}=\overline{N}`$. โ
We call the algebra
(2.10)
$$\overline{N}:=\left(N^{}M\right)_++V\left(N^{}M\right)_{}$$
obtained in Lemma 2.6, the Fermion complement of $`N`$ in $`M`$.
###### Proposition 2.7.
Let $`๐:=_{jI}๐_j`$ be the Fermi $`C^{}`$โalgebra with $`I=\{0,1\}`$.
(i) Let $`P๐_{\{0\}}`$ be a $`\mathrm{\Theta }`$โinvariant projection. Then there is a oneโtoโone correspondence between:
* $`\epsilon :P๐PP๐P`$ an even conditional expectation such that $`(\epsilon )`$ is a full matrix subalgebra of $`P๐_{\{0\}}P`$,
* $`NP๐_{\{0\}}P`$ a $`\mathrm{\Theta }`$โinvariant full matrix subalgebra and $`\mathrm{\Phi }`$ an even state on $`\overline{N}P๐_{\{1\}}P`$.
The correspondence is given for $`xN`$, $`y\overline{N}P๐_{\{1\}}P`$ by
(2.11)
$$\epsilon (xy)=\mathrm{\Phi }(y)x$$
where $`\overline{N}`$ is the Fermion complement of $`N`$ in $`P๐_{\{0\}}P`$ given in (2.10). In particular, $`(\epsilon )=N`$.
(ii) Let $`P_1,P_2๐_{\{0\}}`$ such that $`\mathrm{\Theta }(P_1)=P_2`$, $`P_1P_2=0`$. Then there is a oneโtoโone correspondence between:
* $`\epsilon :(P_1+P_2)๐(P_1+P_2)(P_1+P_2)๐(P_1+P_2)`$ an even conditional expectation such that $`(\epsilon )๐_{\{0\}}`$ and $`๐ต((\epsilon ))=P_1P_2`$,
* $`N_1P_1๐_{\{0\}}P_1`$ full matrix algebra and $`\mathrm{\Phi }`$ a state on $`M_1:=N_1^{}P_1๐P_1`$.
The correspondence is given for $`x_iN_i`$, $`y_iM_i`$, $`i=1,2`$,
(2.12)
$$\epsilon (x_1y_1+x_2y_2)=\mathrm{\Phi }(y_1)x_1+\mathrm{\Phi }(\mathrm{\Theta }(y_2))x_2$$
where $`N_2:=\mathrm{\Theta }(N_1)`$, $`M_2:=\mathrm{\Theta }(M_1)`$. In particular, $`(\epsilon )=N_1\mathrm{\Theta }(N_1)`$.
In addition, if $`z๐_{\{1\}}`$ is even, then
(2.13)
$$\epsilon ((P_1+P_2)z)=\mathrm{\Phi }(P_1zP_1)(P_1+P_2).$$
###### Proof.
(i) Let $`N:=(\epsilon )`$. As $`\epsilon `$ is even, $`N`$ is a $`\mathrm{\Theta }`$โinvariant full matrix algebra of $`P๐_{\{0\}}P`$. Let $`\overline{N}`$ be the Fermion complement of $`N`$ in $`P๐_{\{0\}}P`$, and $`y\overline{N}P๐_{\{1\}}P`$ an odd element. Then $`\epsilon (y)N`$ is odd too, and by the bimodule property of $`\epsilon `$, $`[\epsilon (y),N_+]=0`$. By Lemma 2.5, $`\epsilon (y)=0`$. If $`xN`$, $`y\overline{N}๐_{\{1\}}`$, we have
$$x\epsilon (y)=x\epsilon (y_+)=\epsilon (xy_+)=\epsilon (y_+x)=\epsilon (y_+)x=\epsilon (y)x.$$
This means that $`\epsilon (y)๐ต(N)P`$, that is $`\epsilon (xy)=\mathrm{\Phi }(y)x`$ for a uniquely determined even state $`\mathrm{\Phi }`$ on $`\overline{N}P๐_{\{1\}}P`$.
Fix now an invariant full matrix subalgebra $`N`$ of $`P๐_{\{0\}}P`$. By uniqueness, the Fermion complement of $`N`$ in $`P๐P`$ is all of $`\overline{N}P๐_{\{1\}}P`$. Thus, in order to shorten the notations, we can suppose that $`\overline{N}`$ is the Fermion complement of $`N`$ in $`P๐P`$. Thus,
$$P๐P=N\beta ^1(\overline{N})N\beta ^1(\overline{N}),$$
where $`\beta :\stackrel{~}{N}\overline{N}`$ is the isomorphism given in (2.9). Define $`\epsilon :=E_N^{\mathrm{\Phi }\beta }`$ as the Fubini mapping given in , 9.8.4. Let now $`xN`$, $`y\overline{N}`$. We get
$`\epsilon (xy)=\epsilon (x(y_++y_{}))=\epsilon (xy_+)+\epsilon [(xV)(Vy_{})]`$
$`=`$ $`\mathrm{\Phi }(y_+)x+\mathrm{\Phi }(\beta (Vy_{}))xV=\mathrm{\Phi }(y_+)x+\mathrm{\Phi }(y_{})xV`$
$`=`$ $`\mathrm{\Phi }(y_+)x=\mathrm{\Phi }(y_+)x+\mathrm{\Phi }(y_{})x=\mathrm{\Phi }(y)x`$
as, being $`\mathrm{\Phi }`$ even, it is zero on the odd part of $`\overline{N}`$.
(ii) Take $`N_i:=P_i(\epsilon )P_i`$, $`M_i:=P_i\left((\epsilon )^{}๐\right)P_i`$, $`i=1,2`$. As $`\epsilon `$ is even, we have
$$\mathrm{\Theta }(N_1)=N_2,\mathrm{\Theta }(M_1)=M_2,\mathrm{\Theta }(P_1๐P_1)=P_2๐P_2,$$
and
$`P_1๐P_1+P_2๐P_2=`$ $`N_1{\displaystyle M_1}+N_2{\displaystyle M_2}`$
$``$ $`N_1M_1N_2M_2.`$
As $`\epsilon `$ is uniquely determined by the restriction on the reduced algebras $`๐_{P_i}`$, $`i=1,2`$, according to (2.1) and (2.2), there exist uniquely determined states $`\phi _i`$ on $`M_i`$, such that
$$\epsilon (x_1y_1+x_2y_2)=\phi _1(y_1)x_1+\phi _2(y_2)x_2$$
whenever $`x_iN_i`$, $`y_iM_i`$, $`i=1,2`$. Thus, it is enough to show that $`\phi _2=\phi _1\mathrm{\Theta }`$. We compute
$$\epsilon (\mathrm{\Theta }(x_1y_1+x_2y_2))=\phi _1(\mathrm{\Theta }(y_2))\mathrm{\Theta }(x_2)+\phi _2(\mathrm{\Theta }(y_1))\mathrm{\Theta }(x_1),$$
and
$$\mathrm{\Theta }(\epsilon (x_1y_1+x_2y_2))=\phi _2(y_2)\mathrm{\Theta }(x_2)+\phi _1(y_1)\mathrm{\Theta }(x_1).$$
Thanks to the $`\mathrm{\Theta }`$โequivariance of $`\epsilon `$, we conclude that $`\phi _2=\phi _1\mathrm{\Theta }`$ and vice versa.
Finally, if $`z๐_{\{1\}}`$ is even, then
$$P_izP_iP_izP_i\left((\epsilon )^{}๐_{[0,1]}\right)P_i=M_i,i=1,2.$$
By the first part, we get
$`\epsilon ((P_1+P_2)z)=\mathrm{\Phi }(P_1zP_1)P_1+\mathrm{\Phi }(\mathrm{\Theta }(P_2zP_2))P_2`$
$`=`$ $`\mathrm{\Phi }(P_1zP_1)P_1+\mathrm{\Phi }(P_1zP_1)P_2=\mathrm{\Phi }(P_1zP_1)(P_1+P_2).`$
The previous results relative to the action of the grading automorphism on the centers of the transition expectations allows us to provide the definition of the strongly even and minimal Markov states.
We start by noticing that Definition 2.6 can be slightly generalized by simply requiring that the the local algebras $`๐_j`$ appearing in (2.5) are full matrix $`C^{}`$โalgebras such that the grading automorphism $`\mathrm{\Theta }`$ leaves each algebra $`๐_j`$ globally stable. In this case, the Markov property is still described by the transition expectations $`\epsilon _n`$ previously described, and Lemma 2.4 still works.
###### Definition 2.8.
Let $`\phi ๐ฎ(๐)`$ be a Markov state. It is called strongly even (resp. minimal) if the parity automorphism $`\mathrm{\Theta }`$ acts trivially (resp. transitively) on each $`๐ต((\epsilon _n))`$, $`\epsilon _n`$ being the transition expectations canonically associated to $`\phi `$ through Proposition 2.2.
For some interesting applications (see e.g. Corollary 4.3), it is enough to consider a Markov state as strongly even if $`\mathrm{\Theta }`$ acts trivially on the centers of the transition expectations, infinitely often. Then a Markov state $`\phi `$ will be non strongly even if there exists $`k`$ such that the action of $`\mathrm{\Theta }`$ on $`๐ต((\epsilon _n))`$ is nontrivial for each $`|n|>k`$.
## 3. the structure of general Fermi Markov states
In the present section we investigate the structure of Fermi Markov states. We follow Section 3 of , where we dealt with the quasi local algebra $`๐=\overline{{\displaystyle \underset{jI}{}}๐_{d^j}()}^C^{}`$, equipped with the local structure $`๐_\mathrm{\Lambda }={\displaystyle \underset{j\mathrm{\Lambda }}{}}๐_{d^j}()`$, $`\mathrm{\Lambda }I`$ finite, and trivial $`_2`$โgrading. The forthcoming analysis also represents the extension to the most general Fermion algebra of the results in Section 5 of , where only the homogeneous situation $`๐:=\overline{{\displaystyle \underset{I}{}}๐_2()}^C^{}`$, and only the strongly even Markov states were considered.
The program in cannot be directly implemented in this situation. In fact the parity automorphism $`\mathrm{\Theta }`$ does not act trivially on the centres of $`๐ต((\epsilon _j))`$ in general. Thus, the minimal projections of the centers $`๐ต((\epsilon _j))`$ of the ranges $`(\epsilon _j)`$ does not generate an Abelian algebra. Yet, we are able to decompose non homogeneous Markov states on the Fermion algebras into minimal ones (cf. Definition 2.8).
Let $`\phi `$ be a Markov state, together with the sequence $`\{\epsilon _j\}_{j<j_+}`$ of transition expectations canonically associated to $`\phi `$ as previously explained. We start by considering the centre $`Z_j`$ of $`(\epsilon _j)`$, together with the generating family $`\{P_{\gamma _j}^j\}_{\gamma _j\mathrm{\Gamma }_j}`$ of minimal projections. Define $`\mathrm{\Omega }_j:=\mathrm{\Gamma }_j/`$ where โ$``$โ stands for the equivalence relation induced by $`\mathrm{\Theta }`$ on the spectrum $`\mathrm{\Gamma }_j`$ of $`Z_j`$. Let $`p_j:\mathrm{\Gamma }_j\mathrm{\Omega }_j`$ be the corresponding canonical projection. Put
(3.1)
$$Q_{\omega _j}^j:=\underset{\gamma _j=p_j^1(\{\omega _j\})}{}P_{\gamma _j}^j.$$
For $`j<j_+`$, denote $`C_jZ_j`$ the subalgebra generated by $`\{Q_{\omega _j}^j|\omega _j\mathrm{\Omega }_j\}`$. Notice that $`spec(C_j)=\mathrm{\Omega }_j`$. The $`C_j`$ generate an Abelian subalgebra of $`๐`$ whose spectrum is precisely
(3.2)
$$\mathrm{\Omega }:=\underset{jI}{}\mathrm{\Omega }_j\underset{jI}{}\text{spec}(C_j),$$
where the product in (3.2) stands for the topological product of the finite sets $`\mathrm{\Omega }_j`$. In order to simplify the notations, we define $`\epsilon _{j_+}:=id_{๐_{j_+}}`$. This means $`\mathrm{\Omega }_{j_+}:=\{j_+\}`$, $`Q_{j_+}P_{j_+}:=1\mathrm{I}`$, and finally, for $`N_{j_+}`$, $`\overline{N}_{j_+}`$ given in Proposition 2.7, $`N_{j_+}:=๐_{\{_{j_+}\}}`$, $`\overline{N}_{j_+}:=1\mathrm{I}`$ with an obvious meaning. Put
$$B_j:=\underset{\omega _j\mathrm{\Omega }_j}{}Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j๐_{\{j\}},$$
and
$$๐
:=\overline{\underset{jI}{}B_j}๐.$$
The next step is to construct a conditional expectation of $`๐`$ onto $`๐
`$. Thanks to the fact that the $`Q_{\omega _j}^j`$ are even and thus mutually commuting, we have for each $`x๐_{[k,l]}`$,
$`{\displaystyle \underset{\omega _{k1},\omega _k,\mathrm{},\omega _l,\omega _{l+1}}{}}\left(Q_{\omega _{k1}}^{k1}Q_{\omega _k}^k\mathrm{}Q_{\omega _l}^lQ_{\omega _{l+1}}^{l+1}\right)x\left(Q_{\omega _{k1}}^{k1}Q_{\omega _k}^k\mathrm{}Q_{\omega _l}^lQ_{\omega _{l+1}}^{l+1}\right)`$
$`=`$ $`{\displaystyle \underset{\omega _{k1}}{}}Q_{\omega _{k1}}^{k1}{\displaystyle \underset{\omega _{l+1}}{}}Q_{\omega _{l+1}}^k{\displaystyle \underset{\omega _k,\mathrm{},\omega _l}{}}\left(Q_{\omega _k}^k\mathrm{}Q_{\omega _l}^l\right)x\left(Q_{\omega _k}^k\mathrm{}Q_{\omega _l}^l\right)`$
$`=`$ $`{\displaystyle \underset{\omega _k,\mathrm{},\omega _l}{}}\left(Q_{\omega _k}^k\mathrm{}Q_{\omega _l}^l\right)x\left(Q_{\omega _k}^k\mathrm{}Q_{\omega _l}^l\right).`$
Moreover, if $`x=x_kx_{k+1}\mathrm{}x_l`$, then
$`{\displaystyle \underset{\omega _k,\omega _{k+1},\mathrm{},\omega _l,\omega _l}{}}\left(Q_{\omega _k}^kQ_{\omega _{k+1}}^{k+1}\mathrm{}Q_{\omega _l}^l\right)x\left(Q_{\omega _k}^kQ_{\omega _{k+1}}^{k+1}\mathrm{}Q_{\omega _l}^l\right)`$
$`=`$ $`{\displaystyle \underset{\omega _k,\omega _{k+1},\mathrm{},\omega _l}{}}\left(Q_{\omega _k}^kx_kQ_{\omega _k}^k\right)\left(Q_{\omega _{k+1}}^{k+1}x_{k+1}Q_{\omega _{k+1}}^{k+1}\right)\mathrm{}\left(Q_{\omega _l}^lx_lQ_{\omega _l}^l\right).`$
Thus, on the dense subalgebra $`๐:={\displaystyle \underset{\mathrm{\Lambda }I}{}}๐_\mathrm{\Lambda }`$, $`\mathrm{\Lambda }`$ finite, we get a norm one projection $`E:๐๐
`$, given on the algebraic generators of $`๐`$ by
(3.3)
$$E\left(x_{j_1}\mathrm{}x_{j_n}\right)=\underset{\omega _{j_1},\mathrm{},\omega _{j_n}}{}\left(Q_{\omega _{j_1}}^{j_1}x_{j_1}Q_{\omega _{j_1}}^{j_1}\right)\mathrm{}\left(Q_{\omega _{j_n}}^{j_n}x_{j_n}Q_{\omega _{j_n}}^{j_n}\right)$$
which uniquely extends to a conditional expectation (denoted again by $`E`$ by an abuse of notations) $`E:๐๐
`$ of $`๐`$ onto $`๐
`$. It is also a quite standard fact to see that
$$\phi =\phi E\phi _๐
E.$$
By taking into account the previous considerations we can investigate the structure of Fermi Markov states following the lines in . We recover the following objects canonically associated to the Markov state $`\phi `$ under consideration.
* A classical Markov process on the compact space $`\mathrm{\Omega }`$ given in (3.2), whose law $`\mu `$ is uniquely determined by the sequences of compatible distributions and transition probabilities at the place $`j`$ given respectively by
(3.4) $`\pi _{\omega _j}^j:=\phi (Q_{\omega _j}^j),j<j_+`$
$`\pi _{\omega _j,\omega _{j+1}}^j:={\displaystyle \frac{\phi (\epsilon _j(Q_{\omega _j}^jQ_{\omega _{j+1}}^{j+1}))}{\phi (Q_{\omega _j}^j)}},j<j_+.`$
* For each trajectory $`\omega (\mathrm{},\omega _{j1},\omega _j,\omega _{j+1},\mathrm{})\mathrm{\Omega }`$, the $`C^{}`$โalgebra $`๐
^\omega `$ given by
(3.5)
$$๐
^\omega :=\overline{\underset{jI}{}Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j}^C^{}.$$
Notice that, in the non trivial cases (i.e. when $`I`$ is infinite), $`๐
^\omega `$ cannot be viewed in a canonical way as a subalgebra of $`๐`$. Yet, whenever $`\mathrm{\Lambda }I`$ is finite,
$$๐
_\mathrm{\Lambda }^\omega :=\underset{j\mathrm{\Lambda }}{}Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j$$
is a subalgebra of $`๐_\mathrm{\Lambda }`$ with the identity the projection $`{\displaystyle \underset{j\mathrm{\Lambda }}{}}Q_{\omega _j}^j`$. Namely, $`๐
^\omega `$ is equipped with a canonical localization $`\{๐
_\mathrm{\Lambda }^\omega |\mathrm{\Lambda }\text{finite subset of }I\}`$, and a $`_2`$โgrading implemented by the automorphism $`\mathrm{\Theta }^\omega `$ arising from the restrictions $`\mathrm{\Theta }_{๐_\mathrm{\Lambda }}`$.
* A completely positive identity preserving map $`E^\omega :๐๐
^\omega `$, which is uniquely determined as in (3.3) by
$`x_{j_1}x_{j_2}\mathrm{}x_{j_n}๐`$ $`\left(Q_{\omega _{j_1}}^{j_1}x_{j_1}Q_{\omega _{j_1}}^{j_1}\right)\left(Q_{\omega _{j_2}}^{j_2}x_{j_2}Q_{\omega _{j_2}}^{j_2}\right)`$
(3.6) $`\mathrm{}`$ $`\left(Q_{\omega _{j_n}}^{j_n}x_{j_n}Q_{\omega _{j_n}}^{j_n}\right).`$
The above map satisfies $`E^\omega E=E^\omega `$.
* A sequence $`\{_j^\omega \}_{jI}`$ of maps
$$_j^\omega :๐
_{j+1]}^\omega ๐
_{j]}^\omega $$
given, for $`x๐
_{j1]}^\omega `$, $`y๐
_{[j,j+1]}^\omega `$ by
$$_j^\omega (xy):=\frac{x\epsilon _j(y)}{\pi _{\omega _j,\omega _{j+1}}^j}.$$
###### Proposition 3.1.
The maps $`_j^\omega `$ are even conditional expectations.
###### Proof.
As for each $`kj`$,
$$๐
_{[k,j+1]}^\omega =\left(\underset{l=k}{\overset{j+1}{}}Q_{\omega _l}^l\right)๐_{[k,j+1]}\left(\underset{l=k}{\overset{j+1}{}}Q_{\omega _l}^l\right)๐_{[k,j+1]},$$
and
$$_j^\omega _{๐
_{[k,j+1]}^\omega }=\frac{_j_{๐
_{[k,j+1]}^\omega }}{\pi _{\omega _j,\omega _{j+1}}^j}.$$
Thanks to Proposition 2.2, $`_j^\omega `$ is an even conditional expectation, provided that it is identity preserving. This means that we must check $`_j^\omega (Q_{\omega _j}^jQ_{\omega _{j+1}}^{j+1})=Q_{\omega _j}^j`$. But, we have by (2.13) that $`\epsilon _j(Q_{\omega _j}^jQ_{\omega _{j+1}}^{j+1})=cQ_{\omega _j}^j`$. The proof follows as the number $`c`$ is precisely $`\pi _{\omega _j,\omega _{j+1}}^j`$. โ
* The state $`\psi ^\omega ๐ฎ(๐
^\omega )`$, uniquely determined on localized elements by
(3.7)
$$\psi ^\omega :=\underset{\stackrel{kj_{}}{lj_+}}{lim}\frac{\phi _{Q_{\omega _k}^k๐_{\{k\}}Q_{\omega _k}^k}_k^\omega \mathrm{}_l^\omega }{\pi _{\omega _k}^k}.$$
It is straightforward to check that the state $`\psi ^\omega `$ is a minimal Markov state on $`๐
^\omega `$ w.r.t. the conditional expectations
$$\stackrel{~}{}_j^\omega :=_j^\omega _{j+1}^\omega .$$
In addition, the field
$$\omega \mathrm{\Omega }\psi ^\omega E^\omega ๐ฎ(๐)$$
is $`\sigma (๐^{},๐)`$โmeasurable.
###### Theorem 3.2.
Let $`\phi `$ be a Markov state on the Fermion algebra $`๐`$ w.r.t. the associated sequence $`\{_j\}_{j_{}j<j_+}`$ of conditional expectations given in (2.7). Define the compact set $`\mathrm{\Omega }`$ by (3.2), the probability measure $`\mu `$ on $`\mathrm{\Omega }`$ by (3.4), the quasi local algebra $`๐
^\omega `$ by (3.5), the map $`E^\omega `$ by ((c)), the state $`\psi ^\omega `$ on $`๐
^\omega `$ by (3.7).
Then $`\phi `$ admits the integral decomposition
(3.8)
$$\phi (A)=_\mathrm{\Omega }^{}\psi ^\omega E^\omega (A)\mu (d\omega ),A๐.$$
###### Proof.
We outline the proof which is similar to that of Theorem 3.2 of after writing down the corresponding objects relative to the Fermi case. Consider the Abelian $`C^{}`$โsubalgebra $``$ of $`๐
`$ given by
$$:=\overline{\underset{jI}{}C^j}\overline{\underset{jI}{}C^j}^C^{},$$
together its spectrum $`spec()=\mathrm{\Omega }`$. By restricting $`\phi `$ to $``$, we obtain a possibly nonhomogeneous Markov random process on $`\mathrm{\Omega }`$ with law $`\mu `$ described above. Let $`\pi `$ be the GNS representation of $`๐
`$ relative to $`\phi _๐
`$. Then $`L^{\mathrm{}}(\mathrm{\Omega },\mu )\pi ()^{\prime \prime }\pi (๐
)^{}\pi (๐
)^{\prime \prime }`$. Thus, we have for $`\pi `$ the direct integral decomposition $`\pi ={\displaystyle _\mathrm{\Omega }^{}}\pi _\omega \mu (d\omega )`$, where $`\omega \pi _\omega `$ is a measurable field of representations of $`๐
`$. This leads to the direct integral decomposition of $`\phi _๐
`$, and then the decomposition of $`\phi \phi _๐
E`$ as $`\phi ={\displaystyle _\mathrm{\Omega }}\phi _\omega \mu (d\omega )`$, see e.g. , Section IV. It is then straightforward to see that
$$\phi _\omega (A)=\psi _\omega (E_\omega (A)),$$
almost everywhere on $`\mathrm{\Omega }`$ for each $`A๐`$. โ
The constructive part of Proposition 2.7 allows us to provide the following reconstruction theorem for the class of Fermi Markov states considered in the sequel. It parallels the analogous one (cf. , Theorem 3.3).
Let $`๐`$ be a Fermion algebra. Take for every $`j<j_+`$, a $`\mathrm{\Theta }`$โinvariant commutative subalgebra $`Z_j`$ of $`๐_{\{j\}}`$ with spectrum $`\mathrm{\Gamma }_j`$ and generators $`\{P_{\gamma _j}^j\}_{\gamma _j\mathrm{\Gamma }_j}`$. Put $`Z_{j_+}:=1\mathrm{I}`$. Let โ$``$โ be the equivalence relation on the $`\mathrm{\Gamma }_j`$ induced by the action of $`\mathrm{\Theta }`$, and $`p_j`$ the corresponding projection map. Set $`\mathrm{\Omega }_j:=\mathrm{\Gamma }_j/`$, and define $`Q_{\omega _j}^j`$ as in (3.1). Choose a full matrix subalgebra $`N_{\gamma _j}^jP_{\gamma _j}^j๐_{\{j\}}P_{\gamma _j}^j`$ which is $`\mathrm{\Theta }`$โinvariant whenever $`P_{\gamma _j}^j`$ is a fixed point of $`\mathrm{\Theta }`$.<sup>3</sup><sup>3</sup>3Notice that $`\overline{N}_{\gamma _j}^j`$ given in Proposition 2.7 is also left globally invariant under the parity. Form for $`j<j_+`$, the two point even, faithful conditional expectations
$`\epsilon _{\omega _j,\omega _{j+1}}^j:Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j{\displaystyle Q_{\omega _{j+1}}^{j+1}๐_{\{j+1\}}Q_{\omega _{j+1}}^{j+1}}`$
$`{\displaystyle \underset{\gamma _j=p_j^1(\{\omega _j\})}{}}N_{\gamma _j}^jQ_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j,`$
$`\epsilon _{\omega _{j_+1},j_+}^{j_+1}:Q_{\omega _{j_+1}}^{j_+1}๐_{\{j_+1\}}Q_{\omega _{j_+1}}^{j_+1}{\displaystyle ๐_{\{j_+\}}}`$
$`{\displaystyle \underset{\gamma _{j_+1}=p_j^1(\{\omega _{j_+1}\})}{}}N_{\gamma _{j_+1}}^{j_+1}Q_{\omega _{j_+1}}^{j_+1}๐_{\{j_+1\}}Q_{\omega _{j_+1}}^j,`$
according to Proposition 2.7, by taking for the states in (2.11), (2.12), faithful ones. Define $`๐
^\omega `$, $`E^\omega `$ as in (3.5), ((c)) respectively. For the trajectory $`\omega =(\mathrm{},\omega _{j1},\omega _j,\omega _{j+1},\mathrm{})`$, and $`j<j_+`$, define the map $`_j^\omega `$ as
$$_j^\omega (xy):=x\epsilon _{\omega _j,\omega _{j+1}}^j(y),$$
which is an even faithful conditional expectation according to Proposition 2.2, and Lemma 2.3. Take, for $`j<j_+`$, a compatible sequence of even faithful states $`\phi _j^\omega `$ on $`Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j`$.<sup>4</sup><sup>4</sup>4It can be shown by a standard compactness property (cf. , Proposition 5.1), that the set of sequences of even compatible states $`\phi _j^\omega `$, that is such that $`\phi _{j+1}^\omega =\phi _j^\omega _j^\omega _{๐
_{\{j+1\}}^\omega }`$, is nonvoid. Form the state $`\psi ^\omega ๐ฎ(๐
_\omega )`$ as in (3.7) by taking as initial distributions the $`\phi _j^\omega `$. Finally, fix a Markov process on the product space $`\mathrm{\Omega }:={\displaystyle \underset{jI}{}}\mathrm{\Omega }_j`$ with law $`\mu `$ determined, for $`\omega _j\mathrm{\Omega }_j`$, $`\omega _{j+1}\mathrm{\Omega }_{j+1}`$, by the marginal distributions $`\pi _{\omega _j}^j>0`$, and transition probabilities $`\pi _{\omega _j,\omega _{j+1}}^j>0`$.
###### Theorem 3.3.
In the above notations, the state $`\phi `$ on $`๐`$ given by
$$\phi :=_\mathrm{\Omega }\psi _\omega E_\omega \mu (d\omega )$$
is a Markov state w.r.t. the sequence $`\{_j\}_{j<j_+}`$ of conditional expectations uniquely determined (with the convention $`๐_{\{j_{}1\}}=1\mathrm{I}`$) for $`a๐_{j1]}`$, $`x๐_{\{j\}}`$, $`y๐_{\{j+1\}}`$ by
$`_j(axy)=a{\displaystyle \underset{\omega _j,\omega _{j+1},\omega _{j+2}}{}}\pi _{\omega _j}^j\pi _{\omega _j,\omega _{j+1}}^j\pi _{\omega _{j+1},\omega _{j+2}}^{j+1}`$
$`\times `$ $`_j^\omega \left(Q_{\omega _j}^jx_{j+1}^\omega \left(Q_{\omega _{j+1}}^{j+1}yQ_{\omega _{j+1}}^{j+1}Q_{\omega _{j+2}}^{j+2}\right)Q_{\omega _j}^j\right),jj_+2,`$
$`_j(axy)=a{\displaystyle \underset{\omega _j}{}}\pi _{\omega _j}^j_j^\omega \left(Q_{\omega _j}^jxQ_{\omega _j}^jy\right),j=j_+1.`$
###### Proof.
A straightforward computation shows that, for all generators of the form $`x_k\mathrm{}x_l๐_{[k,l]}`$, $`\phi `$ satisfies (2.6), for the sequence of conditional expectations constructed as above (cf. Theorem 4.1 of ). The proof follows as the state $`\phi `$ is locally faithful, by taking into account Lemma 2.3. โ
## 4. general properties of Fermi Markov states
Let $`\phi ๐ฎ(๐)`$, and $`D_{[k,l]}`$ be the adjusted density matrix of the restriction $`\phi _{[k,l]}`$. For $`k<n<j_+`$, define the unitary $`w_{k,n}(t)๐_{[k,n+1]}`$ as
$$w_{k,n}(t):=D_{[k,n+1]}^{it}D_{[k,n]}^{it},t.$$
The unitaries $`\{w_{k,n}(t)\}_t`$ give rise to a cocycle called transition cocycle when $`\phi `$ is a Markov state (cf ). Denote $`S()`$ the von Neumann entropy (see e.g. ).
The following theorem collects some properties of the Fermi Markov states, which parallel the analogous ones relative to Markov states on tensor product algebras (cf. ). For the natural applications of the properties described below to the variational principle in quantum statistical mechanics, the reader is referred to .
###### Theorem 4.1.
Let $`\phi ๐ฎ(๐)`$ be a locally faithful even state. Then the following assertions are equivalent.
* $`\phi ๐ฎ(๐)`$ be a Markov state;
* for each $`t`$ and $`k<n<j_+`$, $`w_{k,n}(t)๐_{[n,n+1],+}`$.
Moreover, if $`I=`$, $`๐_{\{n\}}=๐_{2^d}()`$ for each $`n`$, and $`\phi `$ is translation invariant, the previous assertions are also equivalent to
* $`S(\phi _{[0,n+1]})S(\phi _{[0,n]})=S(\phi _{[0,1]})S(\phi _{\{0\}})`$, $`n1`$.
###### Proof.
(i)$``$(ii) Thanks to Lemma 4.1 of , if $`\phi `$ is a Markov state, then there exists an unitary $`u_t๐_{[k,n1]}^{}{\displaystyle ๐_{[k,n+1]}}`$ such that, for each $`x๐_{[k,n1]}`$,
$$w_{k,n}(t)xw_{k,n}(t)^{}\sigma _t^{\phi _{[k,n+1]}}\left(\sigma _t^{\phi _{[k,n]}}(x)\right)=u_txu_t^{}x,$$
$`\sigma _t^\phi `$ denoting the modular group of a faithful state $`\phi `$ on a von Neumann algebra, see e.g. . As $`w_{k,n}(t)`$ is even, we have
$$w_{k,n}(t)๐_{[k,n1]}^{}๐_{[k,n+1]}๐_+=๐_{[n,n+1],+},$$
see Lemma 11.1 and Theorem 4.17 of .
(ii)$``$(i) The AccardiโCecchini $`\phi `$โexpectation $`E_{k,n}`$ of $`\phi _{[k,n+1]}`$ w.r.t. the inclusion $`๐_{[k,n]}๐_{[k,n+1]}`$ (cf. ) is written as
$$E_{k,n}(x)=_{[k,n]}^0(w_{k,n}(i/2)^{}xw_{k,n}(i/2))$$
where $`w_{k,n}(i/2)`$ is the analytic continuation at $`i/2`$ of $`w_{k,n}(t)`$, and $`_{[k,n]}^0`$ is the conditional expectation of $`๐_{[k,n+1]}`$ onto $`๐_{[k,n]}`$ preserving the normalized trace. If the $`w_{k,n}(t)`$ satisfy all the properties listed above, the AccardiโCecchini expectation $`E_{k,l}`$ is a $`\phi _{[k,n+1]}`$โpreserving quasi conditional expectation w.r.t. the triplet $`๐_{[k,n1]}๐_{[k,n]}๐_{[k,n+1]}`$. By taking for each fixed $`n`$ the pointwise limit
$$\epsilon _n:=\underset{kj_{}}{lim}(\underset{L}{lim}\frac{1}{L}\underset{l=0}{\overset{L1}{}}\left(E_{k,n}_{๐_{[n,n+1]}})^l\right),$$
we obtain by $`_n(xy):=x\epsilon _n(y)`$, $`x๐_{[n1]}`$, $`y๐_{[n,n+1]}`$, a conditional expectation (cf. Proposition 2.2) fulfilling all the properties listed in Definition 2.1.
(ii)$``$(iii) We have
$$w_{0,n}(t)=[D\phi _{[0,n+1]}:D(\phi _{[0,n+1]}_{[0,n]}^0)]_t,$$
the last being the ConnesโRadonโNikodym cocycle of $`\phi _{[0,n+1]}`$ w.r.t. $`\phi _{[0,n+1]}_{[0,n]}^0`$ (cf. ). The assertion follows from the fact that (iii) is equivalent to the fact that $`๐_{[n,n+1]}`$ is a sufficient subalgebra for both the mentioned states. It turns out to be equivalent to (ii) by translation invariance, see Proposition 11.5 and Proposition 9.3 of . โ
###### Corollary 4.2.
Suppose that $`j_{}I`$. If $`\phi ๐ฎ(๐)`$ is a Markov state, then its support in $`๐^{}`$ is central. In addition, $`\phi `$ is faithful.
###### Proof.
By Theorem 4.1, the pointwise norm limit
$$\underset{nj_+}{lim}D_{[j_{},n]}^{it}xD_{[j_{},n]}^{it}$$
exists as it is asymptotically constant in $`n`$, on localized elements. Thus, it defines a one parameter group of automorphisms $`t\sigma _t`$ of $`๐`$ which admits, by construction, $`\phi `$ as a KMS state. This means that $`\pi _\phi (๐)^{}\xi _\phi `$ is dense in $`_\phi `$, $`(\pi _\phi ,_\phi ,\xi _\phi )`$ being the GNS triplet of $`\phi `$. Furthermore, $`\phi `$ is faithful as $`๐`$ is a simple $`C^{}`$โalgebra, see , Proposition 2.6.17. โ
###### Corollary 4.3.
Suppose that, for each $`nI`$, there exists a $`k(n)I`$ with $`k(n)n`$, such that $`\mathrm{\Theta }`$ acts trivially on $`๐ต((\epsilon _{k(n)}))`$, $`\epsilon _j`$ being the transition expectations associated to the Markov state $`\phi `$. Then the assertions in Corollary 4.2 hold true as well.
###### Proof.
By regrouping the local algebras, we can suppose that there exists a $`j_0I`$ such that, for $`j<j_0`$, $`\mathrm{\Theta }`$ acts trivially on $`๐ต((\epsilon _j))`$. Consider for $`k<j_0`$, $`l>j_0`$ the local algebras
$$๐_{[k,l]}:=_k^c๐_{[k+1,l]},$$
with $`_k^c`$ given in (5.5). The last assertion follows as in Corollary 4.2, by looking at the transition cocycles of $`\phi `$ relative to the new localization $`\{๐_{[k,l]}\}_{k<j_0<j}`$. โ
Let $`\phi `$ be a translation invariant locally faithful state on the Fermion algebra $`๐๐_{}`$. The mean entropy $`s(\phi )`$ of $`\phi `$ (see e.g. ) is defined as
$$s(\phi ):=\underset{n}{lim}\frac{1}{n+1}S(\phi _{[0,n]}),$$
$`S(\phi _{[0,n]})`$ being the von Neumann entropy of $`\phi _{[0,n]}`$.
###### Corollary 4.4.
We have for the translation invariant Markov state $`\phi `$
$$s(\phi )=S(\phi _{[0,1]})S(\phi _{\{0\}}).$$
###### Proof.
It immediately follows by (iii) of Theorem 4.1. โ
## 5. strongly even Markov states
In the present section we investigate the structure of strongly even Markov states (cf. Definition 2.8). By taking into account the structure of the local densities (or equally well the local Hamiltonians by passing to the logaritm) described in (5.4), the strongly even Markov states can be viewed as the Fermi analogue of the Ising type interactions. In addition, they enjoy a kind of local entanglement effect, see Section 4 of for further details.
Notice that the forthcoming analysis extends to the situation when there exists a subsequence $`\{n_j\}I`$ such that $`\mathrm{\Theta }`$ acts trivially on all the $`๐ต((\epsilon _{n_j}))`$.
We start with the following lemma which is known to the experts.
###### Lemma 5.1.
Let $`_n๐
_n`$, $`n`$, be an increasing sequence of inclusions of unital $`C^{}`$โsubalgebras of $`๐
:=\overline{{\displaystyle \underset{n}{}}๐
_n}`$ satisfying $`(_k)^{}๐
_n=_n`$, $`kn`$. Then $`:=\overline{{\displaystyle \underset{n}{}}_n}`$ is a maximal Abelian $`C^{}`$โsubalgebra of $`๐
`$.
###### Proof.
We have for the commutant $`^{}`$ in the ambient algebra $`๐
`$,
$`^{}=`$ $`\overline{{\displaystyle \underset{n}{}}\left(^{}๐
_n\right)}=\overline{{\displaystyle \underset{n}{}}\left(\left({\displaystyle \underset{k}{}}(_k)^{}\right){\displaystyle ๐
_n}\right)}`$
$`=`$ $`\overline{{\displaystyle \underset{n}{}}\left({\displaystyle \underset{k}{}}\left((_k)^{}๐
_n\right)\right)}=\overline{{\displaystyle \underset{n}{}}_n}=.`$
Let $`\omega =(\mathrm{},\omega _{j1},\omega _j,\omega _{j+1},\mathrm{})\mathrm{\Omega }`$ be a trajectory. Thanks to part (i) of Proposition 2.7,
(5.1) $`๐
^\omega \overline{\left({\displaystyle \underset{j<j_+}{}}Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j\right){\displaystyle ๐_{\{j_+\}}}}^C^{}`$
(5.2) $`=`$ $`\overline{N_{\omega _j_{}}^j_{}{\displaystyle \left(\underset{j<j_+1}{}(\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1})\right)(\overline{N}_{\omega _{j_+1}}^{j_+1}๐_{\{j_+\}})}}^C^{},`$
$`N_{\omega _j}^j`$ $`\overline{N}_{\omega _j}^j`$ providing the (Fermi) decompositions of the $`Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j`$ described by (2.10) in Proposition 2.7. This decomposition is quite similar to the analogous one described in Theorem 3.2 of , and generalize the situation treated in Section 5 of .
###### Lemma 5.2.
Any maximal Abelian subalgebra of $`\left(\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1}\right)_+`$ is maximal Abelian in $`\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1}`$ as well.
###### Proof.
Let $`V\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1}`$ be any selfadjoint unitary implementing $`\mathrm{\Theta }`$. Then
$`\left(\overline{N}_{\omega _j}^j{\displaystyle N_{\omega _{j+1}}^{j+1}}\right)_+\left(E_1E_1\right)^{}`$
$`=`$ $`E_1\left(\overline{N}_{\omega _j}^j{\displaystyle N_{\omega _{j+1}}^{j+1}}\right)E_1E_1\left(\overline{N}_{\omega _j}^j{\displaystyle N_{\omega _{j+1}}^{j+1}}\right)E_1,`$
$`V=E_1E_1`$ being the resolution of $`V`$. โ
###### Lemma 5.3.
The unnormalized trace of
$$R:=N_{\omega _k}^k\overline{N}_{\omega _k}^k\mathrm{}N_{\omega _l}^l\overline{N}_{\omega _l}^l$$
is the product of the unnormalized traces of the $`N_{\omega _j}^j`$ and $`\overline{N}_{\omega _j}^j`$, $`kjl`$.
###### Proof.
Put $`R=\left(Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j\right){\displaystyle \mathrm{}\left(Q_{\omega _{j+1}}^{j+1}๐_{\{j+1\}}Q_{\omega _{j+1}}^{j+1}\right)}`$. By the product property of $`Tr_{๐_{[k,l]}}`$, we get
$$Tr{}_{R}{}^{}=\underset{j=k}{\overset{l}{}}Tr{}_{Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j}{}^{}.$$
Thus, we reduce the situation to the algebra $`N_{\omega _j}^j\overline{N}_{\omega _j}^jQ_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j`$. Notice that
$$Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j=N_{\omega _j}^j\stackrel{~}{N}_{\omega _j}^jN_{\omega _j}^j\stackrel{~}{N}_{\omega _j}^j,$$
$`N_{\omega _j}^j`$, $`\stackrel{~}{N}_{\omega _j}^j`$ are both globally stable under the action of $`\mathrm{\Theta }`$, $`\overline{N}_{\omega _j}^j=\stackrel{~}{N}_{\omega _j,+}^j+V\stackrel{~}{N}_{\omega _j,}^j`$, $`V`$ being any unitary of $`N_{\omega _j}^j`$ implementing $`\mathrm{\Theta }`$ on itself, see Proposition 2.7. As the traces are invariant under any automorphism, we get
$`Tr{}_{Q_{\omega _j}^j๐_{\{j\}}Q_{\omega _j}^j}{}^{}=Tr{}_{N_{\omega _j}^j}{}^{}Tr_{\stackrel{~}{N}_{\omega _j}^j}`$
$`=`$ $`\left(Tr{}_{N_{\omega _j,+}^j}{}^{}{\displaystyle \frac{id+\mathrm{\Theta }}{2}}\right)\left(Tr{}_{\stackrel{~}{N}_{\omega _j,+}^j}{}^{}{\displaystyle \frac{id+\mathrm{\Theta }}{2}}\right)`$
$`=`$ $`\left(Tr{}_{N_{\omega _j,+}^j}{}^{}{\displaystyle \frac{id+\mathrm{\Theta }}{2}}\right)\left(Tr{}_{\overline{N}_{\omega _j,+}^j}{}^{}{\displaystyle \frac{id+\mathrm{\Theta }}{2}}\right)`$
$`=`$ $`Tr{}_{N_{\omega _j}^j}{}^{}Tr{}_{\overline{N}_{\omega _j}^j}{}^{}.`$
Let the initial distributions $`\eta _{\omega _j_{}}^j_{}๐ฎ(N_{\omega _j_{}}^j_{})`$, the states $`\eta _{\omega _j,\omega _{j+1}}^j๐ฎ(\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1})`$ be recovered by $`\phi `$ according to (2.11).<sup>5</sup><sup>5</sup>5If $`j_{}`$ and/or $`j_+`$ do not belong to $`I`$, they do not appear in the formulae, the last having an obvious meaning. In addition, as $`\mathrm{\Omega }_{j_+}\{j_+\}`$, we use the symbology $`\eta _{\omega _j,\omega _{j+1}}^j`$ also for the final distributions $`\eta _{\omega _{j_+1},j_+}^{j_+1}`$. Consider the even densities $`T_{\omega _j}^{(j)}`$, $`\widehat{T}_{\omega _j}^{(j)}`$, $`T_{\omega _j,\omega _{j+1}}^{(j)}`$ localized in $`N_{\omega _j}^j`$, $`\overline{N}_{\omega _j}^j`$, $`\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1}`$, and associated to $`\eta _{\omega _j_{}}^j_{}`$ or $`\eta _{\omega _j,\omega _{j+1}}^j_{N_{\omega _{j+1}}^{j+1}}`$, $`\eta _{\omega _j,\omega _{j+1}}^j_{\overline{N}_{\omega _j}^j}`$, $`\eta _{\omega _j,\omega _{j+1}}^j`$ respectively.
###### Proposition 5.4.
The states $`\eta _{\omega _j_{}}^j_{}`$, $`\eta _{\omega _j,\omega _{j+1}}^j`$ uniquely define a product state on $`๐
^\omega `$, coinciding with $`\psi ^\omega `$ in (3.8), which is symbolically written as
(5.3)
$$\psi ^\omega =\eta _{\omega _j_{}}^j_{}\underset{jj_+1}{}\eta _{\omega _j,\omega _{j+1}}^j.$$
###### Proof.
Consider on $`๐
^\omega `$ the localization
$$๐
^\omega =\overline{\underset{jI}{}๐_{\omega _j}^j}$$
suggested by (5.1). Here, $`๐_{\omega _j_{}}^j_{}:=N_{\omega _j_{}}^j_{}`$, $`๐_{\omega _j}^j:=\overline{N}_{\omega _{j1}}^{j1}N_{\omega _j}^{j+1}`$, $`j_{}<j<j_+`$, and finally $`๐_{\omega _{j_+}}^{j_+}:=\overline{N}_{\omega _{j_+1}}^{j_+1}๐_{\{j_+\}}`$. As the above densities commute each other, for each $`k<l`$, the product of local densities
$$T_{\omega _{k1},\omega _k}^{(k1)}\times \mathrm{}\times T_{\omega _{l1},\omega _l}^{(l1)}$$
is a well defined positive even operator on $`_{kjl}๐_{\omega _j}^j`$ which by Lemma 5.3, is the density of $`\psi ^\omega _{_{kjl}๐_{\omega _j}^j}`$ w.r.t. the unnormalized trace of $`_{kjl}๐_{\omega _j}^j`$. As explained in Section 2.3, this means that $`\psi ^\omega `$ is the product states of $`\eta _{\omega _j_{}}^j_{}`$ with the $`\eta _{\omega _j,\omega _{j+1}}^j`$ as explained in (5.3) (see Theorem 11.2 of for a similar situation). โ
As all the states appearing in (5.3) are even, we can explicitely write the local densities associated to the strongly even Markov state. Namely, consider the RadonโNikodym derivatives (i.e. the densities) $`T_{๐_{[k,l]}}`$ w.r.t. the unnormalized trace of $`๐_{[k,l]}`$,
$$\phi _{[k,l]}=Tr{}_{๐_{[k,l]}}{}^{}(T_{๐_{[k,l]}}).$$
Then $`T_{๐_{[k,l]}}`$ has the nice decomposition
(5.4)
$$T_{๐_{[k,l]}}=\underset{\omega _k,\mathrm{},\omega _l}{}T_{\omega _k}^{(k)}T_{\omega _k,\omega _{k+1}}^{(k)}\times \mathrm{}\times T_{\omega _{l1},\omega _l}^{(l1)}\widehat{T}_{\omega _l}^{(l)}.$$
By Corollary 4.2, any strongly even Markov state is a KMS state for the one parameter group of automomorphisms $`\sigma _t`$ given, for $`x๐`$, by
$$\sigma _t(x):=\underset{\stackrel{kj_{}}{lj_+}}{lim}T_{๐_{[k,l]}}^{it}xT_{๐_{[k,l]}}^{it}.$$
In addition, each strongly even Markov state is faithful.
We now show that each strongly even Markov state is a lifting of a classical Markov processes. This result parallels the analogous one relative to the tensor product algebra, obtained first in for some particular cases, and then in for the general situation. Such property was called diagonalizability in . After adapting the situation relative to the tensor product case to the strongly even Fermi Markov states, we can follow the same line of the proof of Theorem 3.2 of .
We start by defining increasing subalgebras of the Fermion algebra $`๐`$ equipped with a natural local structure inherited from that of the original algebra. Let $`_j:=(\epsilon _j)`$, with relative commutant
(5.5)
$$_j^c:=_j^{}๐_{\{j\}}.$$
Define
(5.6) $`๐_{\{k\}}:=๐ต(_k),๐_{[k,k+1]}:=_k^c{\displaystyle _{k+1}},`$
$`๐_{[k,l]}:=_k^c{\displaystyle ๐_{[k+1,l1]}_l},k<l+1.`$
Thanks to Lemma 5.2, for each $`kj<l`$ and $`\omega _j\mathrm{\Omega }_j`$, we can choose a even maximal Abelian subalgebra $`D_{\omega _j,\omega _{j+1}}^j`$ of $`\overline{N}_{\omega _j}^jN_{\omega _{j+1}}^{j+1}`$ containing $`T_{\omega _j,\omega _{j+1}}^{(j)}`$. Put
$`๐_{\{k\}}:=๐_{\{k\}}๐ต(_k),`$
(5.7) $`๐_{[k,l]}:={\displaystyle \underset{\omega _k,\mathrm{},\omega _l}{}}\left(D_{\omega _k,\omega _{k+1}}^k{\displaystyle \mathrm{}D_{\omega _{l1},\omega _l}^{l1}}\right),k<l,`$
$`๐:=\overline{\left({\displaystyle \underset{[k,l]I}{}}๐_{[k,l]}\right)}.`$
###### Theorem 5.5.
Let $`\phi ๐ฎ(๐)`$ be a strongly even Markov state. Then there exists an even maximal Abelian $`C^{}`$โsubalgebra $`๐๐`$, and a conditional expectation $`๐:๐๐`$ such that $`\phi =\phi _๐๐`$. In addition, the measure $`\mu `$ on $`spec(๐)`$ associated to $`\phi _๐`$ is a Markov measure w.r.t. the natural localization of $`๐`$ given in (5).
###### Proof.
Let $`[m_k,n_k]`$ be an increasing sequence of intervals such that $`[m_k,n_k]I`$. Then
$$๐=\overline{\left(\underset{\stackrel{}{[m_k,n_k]I}}{lim}๐_{[m_k,n_k]}\right)}^C^{}.$$
As $`๐_{[m,n]}`$ is an even maximal Abelian subalgebra of $`๐_{[m,n]}`$, the increasing sequence $`๐_{[m_k,n_k]}๐_{[m_k,n_k]}`$ satisfies the hypotheses of Lemma 5.1. Thus, $`๐`$ is a even maximal Abelian $`C^{}`$โsubalgebra of $`๐`$. According to (5.4), we have
$$T_{๐_{[m,n]}}=\underset{\omega _m,\mathrm{},\omega _n}{}T_{\omega _m,\omega _{m+1}}^{(m)}\times \mathrm{}\times T_{\omega _{n1},\omega _n}^{(n1)},$$
that is, $`\{T_{๐_{[m,n]}}\}_{m<n}๐`$. Let $`E_{m,n}^0:๐_{[m,n]}๐_{[m,n]}`$ be the canonical conditional expectation of $`๐_{[m,n]}`$ onto the maximal abelian subalgebra $`๐_{[m,n]}`$ (cf. , Footnote 4). We have
$`\phi _{๐_{[m,n]}}`$ $`Tr{}_{๐_{[m,n]}}{}^{}(T_{๐_{[m,n]}})=Tr{}_{๐_{[m,n]}}{}^{}\left(E_{m,n}^0(T_{๐_{[m,n]}})\right)`$
(5.8) $`=`$ $`Tr{}_{๐_{[m,n]}}{}^{}\left(T_{๐_{[m,n]}}E_{m,n}^0()\right)\phi _{๐_{[m,n]}}E_{m,n}^0.`$
As the sequence $`\{E_{m,n}^0\}_{m<n}`$ is projective, the direct limit $`\underset{\stackrel{}{[m,n]I}}{lim}E_{m,n}^0`$ uniquely defines a conditional expectation $`๐:๐๐`$ fulfilling by (5), $`\phi =\phi _๐๐`$. The measure $`\mu `$ on $`spec(๐)`$ associated to $`\phi _๐`$ is a Markov measure w.r.t. the natural localization of $`๐`$ previously described. This follows as in Section 6 of , after noticing that $`D_{\omega _m,\omega _{m+1}}^m\mathrm{}D_{\omega _{n1},\omega _n}^{n1}`$ in (5) generates a tensor product, and the restriction $`\phi _{D_{\omega _m,\omega _{m+1}}^m{\scriptscriptstyle \mathrm{}D_{\omega _{n1},\omega _n}^{n1}}}`$ defines a product measure on $`spec\left(D_{\omega _m,\omega _{m+1}}^m\right)\times \mathrm{}\times spec\left(D_{\omega _n,\omega _{n+1}}^n\right)`$. โ
Now we pass to the dynamical entropy $`h_\phi (\alpha )`$ w.r.t. the right shift $`\alpha `$ for translation invariant strongly even Markov states. The reader is referred to for the definition and technical details on the dynamical entropy.
The definition of the dynamical entropy $`h_\phi (\alpha )`$ is based on the multiple subalgebra entropy $`H_\phi (N_1,\mathrm{},N_k)`$, with $`N_1,\mathrm{},N_kM`$. We start by pointing out that, if the subalgebras $`N_1,\mathrm{},N_k`$ are the range of $`\phi `$โpreserving conditional expectations and are contained in different factors of a tensor product algebra, then
(5.9)
$$H_\phi (N_1,\mathrm{},N_k)=S(\phi _N),$$
with $`N:=N_1\mathrm{}N_k`$.<sup>6</sup><sup>6</sup>6Fix a faithful trace on $`M`$. Let $`T_1,\mathrm{},T_k`$, $`T`$ be the corresponding densities of $`N_1,\mathrm{},N_k`$, $`M`$ respectively. Choose maximal Abelian subalgebras $`A_j`$ of $`N_j`$ containing $`T_j`$, $`j=1,\mathrm{},k`$. As the $`N_j`$ are expected, we have for $`aA_j`$,
$$T^{it}aT^{it}=T_j^{it}aT_j^{it}=a,$$
that is $`A_jM_\phi `$, $`M_\phi `$ being the centralizer of the faithful state $`\phi `$. As the $`A_j`$ are contained in different factors of a tensor product, $`A_1\mathrm{}A_k`$ is maximal Abelian in $`N`$. Thus, (5.9) follows by Corollary VIII.8 of .
###### Theorem 5.6.
Let $`\phi ๐ฎ(๐)`$ be a translation invariant strongly even Markov state. Then $`h_\phi (\alpha )=s(\phi )`$.
###### Proof.
The proof follows the same lines of the tensor product case. We keep into account some boundary effects which cannot be neglected in proving the result. Fix $`n`$, and consider $`๐_{[0,n+1]}`$ given in (5.6). We have $`๐_{[1,n]}๐_{[0,n]}๐_{[0,n+1]}`$, and $`๐_{[0,n]}`$ is expected. We compute,
$`H(k):=`$ $`H_\phi (๐_{[0,n]},\alpha (๐_{[0,n]}),\mathrm{},\alpha ^{k(n+2)}(๐_{[0,n]}))`$
$``$ $`H_\phi (๐_{[0,n]},\alpha ^{n+2}(๐_{[0,n]}),\mathrm{},\alpha ^{k(n+2)}(๐_{[0,n]}))`$
$``$ $`H_\phi (๐_{[0,n],+},\alpha ^{n+2}(๐_{[0,n],+}),\mathrm{},\alpha ^{k(n+2)}(๐_{[0,n],+})),`$
Now, $`๐_{[0,n],+},\alpha ^{n+2}(๐_{[0,n],+},),\mathrm{},\alpha ^{k(n+2)}(๐_{[0,n],+})`$ are all expected, and generate a tensor product. Then
$$H(k)S(\phi _{M_k})=S(\phi _{M_k},\tau _{M_k})+k\mathrm{ln}d.$$
Here, $`M_k:=๐_{[0,n],+}\alpha ^{n+2}(๐_{[0,n],+})\mathrm{}\alpha ^{k(n+2)}(๐_{[0,n],+})`$, $`d`$ is the tracial dimension of $`๐_{[0,n],+}`$, $`\tau `$ the normalized trace on $`๐`$, and finally $`S(,)`$ the relative entropy (see e.g. ). As $`๐_{[1,m],+}๐_{[0,m],+}`$, and the tracial dimension of $`๐_{[1,m],+}`$ coincides with that of $`๐_{[1,m]}`$ (cf. Lemma 5.2), we obtain by the monotonicity of the relative entropy,
$`H(k)`$ $`S(\phi _{๐_{[1,(n+2)(k+1)]}},\tau _{๐_{[1,(n+2)(k+1)]}})+k\mathrm{ln}d`$
$`=`$ $`S(\phi _{๐_{[1,(n+2)(k+1)]}})+[kn(n+2)(k+1)]\mathrm{ln}l,`$
$`l`$ being the tracial dimension of $`๐_{\{0\}}`$. Finally, we get
$`h_\phi (\alpha )\underset{k}{lim}{\displaystyle \frac{H(k)}{(n+2)k}}\underset{k}{lim}[{\displaystyle \frac{k+1}{k}}s(\phi )`$
$`+`$ $`{\displaystyle \frac{kn(n+2)(k+1)}{(n+2)k}}\mathrm{ln}l]=s(\phi ){\displaystyle \frac{2\mathrm{ln}l}{n+2}}.`$
Since $`h_\phi (\alpha )s(\phi )`$ and $`n`$ is arbitrary, the assertion follows. โ
## 6. examples of translation invariant Fermi Markov states
In the present section we exhibit some examples of Fermi Markov states. We restrict the matter to the translation invariant situation. The non homogeneous cases can be analogously treated. The present construction furnishes the direct application of Theorem 3.3, or equally well Proposition 2.7. Thanks to the translation invariance, it is enough to construct a two point even transition expectation $`\epsilon :๐_{[0,1]}๐_{\{0\}}`$, and compute the stationary even distributions by solving $`\rho =\rho \epsilon \alpha _{๐_{\{0\}}}`$, $`\rho `$ running into the even states of $`๐_{\{0\}}`$. A translation invariant Markov state $`\phi `$ is then recovered by the marginals
(6.1)
$$\phi (x_k\mathrm{}x_l)=\rho (\epsilon _k(x_k\epsilon _{k+1}(x_{k+1}\mathrm{}\epsilon _{l1}(x_{l1}\epsilon _l(x_l))\mathrm{}))).$$
### 6.1. Case 1:
$`๐_{\{n\}}๐_2()`$, $`๐ต((\epsilon ))^2`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=1`$.
We start with the pivotal example in Subsection 6.4 of by showing that it provides examples of Fermi Markov states which are entangled. Define, for a fixed $`\chi `$ in the unit circle $`๐`$,
$$q_\chi :=\frac{1}{2}\left(1\mathrm{I}+\chi a_0+\overline{\chi }a_0^+\right).$$
Choose a faithful state $`\eta ๐ฎ(q_\chi ๐_{[0,1]}q_\chi )`$. Put
(6.2)
$$\epsilon (x)=\eta (q_\chi xq_\chi )q_\chi +\eta (q_\chi \mathrm{\Theta }(x)q_\chi )q_\chi ,x๐_{[0,1]}.$$
With $`\tau `$ the normalized trace on $`๐_2()`$, $`\epsilon _n:=\epsilon \alpha ^n`$, and $`x_k๐_{\{k\}},\mathrm{},x_l๐_{\{l\}}`$, the marginals (6.1) with $`\rho =\tau `$, uniquely determine a translation invariant locally faithful Markov state $`\phi `$ on the Fermion algebra $`๐:=๐_{}`$ satisfying the required properties. Thanks to shift invariance, it suffices to consider $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$.
Let $`\xi _\chi `$, $`\xi _\chi ^{}`$ be the (uniquely determined up to a phase) eigenvectors of $`q_\chi `$, $`q_\chi =q_\chi ^{}`$ acting on $`^2`$, corresponding to the eigenvalues 1, respectively. Put
$$V:=,\xi _\chi \xi _\chi ^{}.$$
As $`V๐_2()=๐_{\{0\}}๐_{[0,1]}`$, $`V`$ is also in $`๐_{[0,1]}`$. Put $`\delta :=\eta (V(\chi a_1+\overline{\chi }a_1^+)q_\chi ))`$. We have
$$\phi (xy)=x_+\xi _\chi ,\xi _\chi \alpha ^1(y_+)\xi _\chi ,\xi _\chi +\delta x_{}\xi _\chi ,\xi _\chi ^{}\alpha ^1(y_{})\xi _\chi ,\xi _\chi .$$
Now we show that there exists a faithful state $`\eta `$ as above, such that $`\eta (X)0`$, where
(6.3)
$$X:=V\alpha (q_{\chi ,})q_\chi \frac{1}{2}V(\chi a_1+\overline{\chi }a_1^+)q_\chi )0.$$
Pick a functional which is different from zero on $`X`$, hence a state $`\eta _0`$ which is nonnull on $`X`$. Let $`pq_\chi ๐_{[0,1]}q_\chi `$ be the support of $`\eta _0`$. Choose a state $`\eta _1`$ with support $`q_\chi p`$. Then $`\eta :=\beta \eta _0+(1\beta )\eta _1`$ is a faithful state on $`q_\chi ๐_{[0,1]}q_\chi `$ which is nonnull on $`X`$ for an appropriate choice of $`\beta [0,1]`$.<sup>7</sup><sup>7</sup>7The last claim easily follows as $`\eta (X)=0`$ means $`\eta _0(X)\eta _1(X)`$, and $`\beta ={\displaystyle \frac{\eta _1(X)}{\eta _1(X)\eta _0(X)}}`$. We then have the following
###### Proposition 6.1.
Let $`\mathrm{\Lambda }_1,\mathrm{\Lambda }_2`$ such that $`\mathrm{\Lambda }_1\mathrm{\Lambda }_2=\mathrm{}`$, $`\mathrm{\Lambda }_1\mathrm{\Lambda }_2=`$. Suppose that $`\eta (X)0`$, where $`\eta `$ is the state in (6.2) and $`X`$ is given in (6.3). Then the state $`\phi `$ described above is entangled w.r.t. the decomposition $`๐=\overline{๐_{\mathrm{\Lambda }_1}๐_{\mathrm{\Lambda }_2}}`$.
###### Proof.
Let $`๐_{\{n\}}๐_{\mathrm{\Lambda }_1}`$, $`๐_{\{n+1\}}๐_{\mathrm{\Lambda }_2}`$ for some $`n`$ (which is always the case after a possible renumbering of $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$). Under the above assumption, $`\phi (x_{}y_{})`$ cannot be identically zero for each $`x๐_{\mathrm{\Lambda }_1}`$, $`y๐_{\mathrm{\Lambda }_2}`$ due to the shift invariance. The proof now follows by applying the Moriya criterion established in Proposition 1 of . โ
By extending the previous computations to more general cases, it is then possible to construct many examples of entangled translation invariant Fermi Markov states for the situation when $`๐_{\{0\}}=๐_{2^d}`$, $`d>1`$. We are going to describe a sample of pivotal examples.
We now consider the successive step $`๐_{\{k\}}๐_4()`$. We exhibit examples for each possible structure of the Abelian algebra $`๐ต((\epsilon ))`$, and for the action of $`\mathrm{\Theta }`$ on it. Let $`\{a_i,a_i^+|i=1,2\}`$ be the creators and annihilators generating $`๐_{\{0\}}`$. Consider the system $`\{e_{kl}(j)|j,k,l=1,2\}`$ of commuting $`2\times 2`$ matrix units obtained via the JordanโKleinโWigner transformation (2.3). Putting $`e_{(i,j)(k,l)}:=e_{ik}(1)e_{jl}(2)`$, we obtain a system of matrix units for $`๐_{\{0\}}`$ which realizes the isomorphism $`๐_{\{0\}}๐_2()๐_2()`$.
### 6.2. Case 2:
$`๐ต((\epsilon ))^4`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=4`$.
Choose $`\{e_{(i,j)(i,j)}|i,j=1,2\}`$ as the generators of $`๐ต((\epsilon ))`$. In this situation, there exist even states $`\phi _{ij}`$, $`i,j=1,2`$ on $`๐_{\{1\}}`$ such that for $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$,
$$\epsilon (xy)=\underset{i,j=1}{\overset{2}{}}Tr(xe_{(i,j)(i,j)})\phi _{ij}(y)e_{(i,j)(i,j)}.$$
This is nothing but the example in Subsection 6.2 of . Thus, $`\phi `$ is strongly clustering w.r.t. the shift on the chain, and the von Neumann algebra $`\pi _\phi (๐)^{\prime \prime }`$ generated by the GNS representation $`\pi _\phi `$ of $`\phi `$ is a type $`\mathrm{III}_\lambda `$ factor for some $`\lambda (0,1]`$, see .
### 6.3. Case 3:
$`๐ต((\epsilon ))^4`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=3`$.
For a fixed $`\chi `$ in the unit circle $`๐`$, define
$$Q_\chi :=\frac{1}{2}\left(1\mathrm{I}+\chi a_2+\overline{\chi }a_2^+\right).$$
Choose $`\{e_{(1,j)(1,j)},e_{22}(1)Q_{\pm \chi }|j=1,2\}`$ as the generators of $`๐ต((\epsilon ))`$. In this situation, there exist even states $`\phi _j`$, $`j=1,2`$ on $`๐_{\{1\}}`$, and a state $`\phi `$ on $`e_{22}(1)Q_\chi ๐_{[0,1]}e_{22}(1)Q_\chi `$ such that, for $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$,
$`\epsilon (xy)=`$ $`{\displaystyle \underset{j=1}{\overset{2}{}}}Tr(xe_{(1,j)(1,j)})\phi _j(y)e_{(1,j)(1,j)}`$
$`+`$ $`\phi (e_{22}(1)Q_\chi xye_{22}(1)Q_\chi )e_{22}(1)Q_\chi `$
$`+`$ $`\phi (e_{22}(1)Q_\chi \mathrm{\Theta }(xy)e_{22}(1)Q_\chi )e_{22}(1)Q_\chi .`$
### 6.4. Case 4:
$`๐ต((\epsilon ))^4`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=2`$.
First choose $`\{e_{ii}(1)Q_{\pm \chi }|i=1,2\}`$ as the generators of $`๐ต((\epsilon ))`$. In this situation, there exist states $`\phi _i`$, on $`e_{ii}(1)Q_\chi ๐_{[0,1]}e_{ii}(1)Q_\chi `$, $`i=1,2`$ such that, for $`x๐_{[0,1]}`$,
$`\epsilon (x)=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}(\phi _i(e_{ii}(1)Q_\chi xe_{ii}(1)Q_\chi )e_{ii}(1)Q_\chi `$
$`+`$ $`\phi _i(e_{ii}(1)Q_\chi \mathrm{\Theta }(x)e_{ii}(1)Q_\chi )e_{ii}(1)Q_\chi ).`$
Next, for fixed $`(\chi ,\eta )๐^2`$, define with $`V:=a_1^+a_1a_1a_1^+`$,
$$P_{\chi ,\eta }:=\frac{1}{4}\left(1\mathrm{I}+\chi a_1+\overline{\chi }a_1^+\right)\left(1\mathrm{I}+\eta Va_2+\overline{\eta }Va_2^+\right).$$
Choose $`\{P_{\pm \chi ,\pm \eta }\}`$ as the generators of $`๐ต((\epsilon ))`$. In this situation, there exist states $`\phi _\pm `$ on $`P_{\pm \chi ,\eta }๐_{[0,1]}P_{\pm \chi ,\eta }`$ respectively, such that for $`x๐_{[0,1]}`$,
$`\epsilon (x)=`$ $`\phi _+(P_{\chi ,\eta }xP_{\chi ,\eta })P_{\chi ,\eta }+\phi _+(P_{\chi ,\eta }\mathrm{\Theta }(x)P_{\chi ,\eta })P_{\chi ,\eta }`$
$`+`$ $`\phi _{}(P_{\chi ,\eta }xP_{\chi ,\eta })P_{\chi ,\eta }+\phi _{}(P_{\chi ,\eta }\mathrm{\Theta }(x)P_{\chi ,\eta })P_{\chi ,\eta }.`$
### 6.5. Case 5:
$`๐ต((\epsilon ))^3`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=3`$.
First choose $`\{e_{11}(1)e_{jj}(2),e_{22}(1)|j=1,2\}`$ as the generators of $`๐ต((\epsilon ))`$. We have two possibilities. Namely, there exist even states $`\phi _j`$, on $`๐_{\{1\}}`$, $`i=1,2`$, and an even state $`\phi `$ either on $`๐_{\{1\}}`$, or on $`(e_{22}(1)๐_{\{0\}}e_{22}(1))๐_{\{1\}}`$ such that, for $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$,
$$\epsilon (xy)=\underset{j=1}{\overset{2}{}}Tr(xe_{11}(1)e_{jj}(2))\phi _j(y)e_{11}(1)e_{jj}(2)+\phi (y)e_{22}(1)xe_{22}(1),$$
respectively
$$\epsilon (xy)=\underset{j=1}{\overset{2}{}}Tr(xe_{11}(1)e_{jj}(2))\phi _j(y)e_{11}(1)e_{jj}(2)+\phi (e_{22}(1)xe_{22}(1)y)e_{22}(1).$$
Next, put $`P:=e_{(1,2)(1,2)}+e_{(2,1)(2,1)}`$ and choose $`\{e_{(i,i)(i,i)},P|i=1,2\}`$ as the generators of $`๐ต((\epsilon ))`$. Again, we have two possibilities. Namely, there exist even states $`\phi _j`$, on $`๐_{\{1\}}`$, $`i=1,2`$, and an even state $`\phi `$ either on $`๐_{\{1\}}`$, or on $`(P๐_{\{0\}}P)๐_{\{1\}}`$ such that, for $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$,
$$\epsilon (xy)=\underset{j=1}{\overset{2}{}}Tr(xe_{11}(1)e_{jj}(2))\phi _j(y)e_{11}(1)e_{jj}(2)+\phi (y)PxP,$$
respectively
$$\epsilon (xy)=\underset{j=1}{\overset{2}{}}Tr(xe_{11}(1)e_{jj}(2))\phi _j(y)e_{11}(1)e_{jj}(2)+\phi (PxPy)P.$$
Notice that the last possibilities correspond to nontrivial cases with $`(\epsilon )๐_+`$.
### 6.6. Case 6:
$`๐ต((\epsilon ))^3`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=2`$.
For $`\chi ๐`$, choose $`\{e_{11}(1)Q_{\pm \chi },e_{22}(1)\}`$ as the generators of $`๐ต((\epsilon ))`$. We have two possibilities. Namely, choose a state $`\phi `$ on $`e_{11}(1)Q_\chi ๐_{[0,1]}e_{11}(1)Q_\chi `$, and a even state $`\psi `$ either on $`๐_{\{1\}}`$, or on $`(e_{22}(1)๐_{\{0\}}e_{22}(1))๐_{\{1\}}`$ such that, for $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$,
$`\epsilon (xy)=`$ $`\phi (e_{11}(1)Q_\chi xye_{11}(1)Q_\chi )e_{11}(1)Q_\chi `$
$`+`$ $`\phi (e_{11}(1)Q_\chi \mathrm{\Theta }(xy)e_{11}(1)Q_\chi )e_{11}(1)Q_\chi `$
$`+`$ $`\psi (y)e_{22}(1)xe_{22}(1),`$
respectively
$`\epsilon (xy)=`$ $`\phi (e_{11}(1)Q_\chi xye_{11}(1)Q_\chi )e_{11}(1)Q_\chi `$
$`+`$ $`\phi (e_{11}(1)Q_\chi \mathrm{\Theta }(xy)e_{11}(1)Q_\chi )e_{11}(1)Q_\chi `$
$`+`$ $`\psi (e_{22}(1)xe_{22}(1)y)e_{22}(1).`$
### 6.7. Case 7:
$`๐ต((\epsilon ))^2`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=2`$.
We treat only the following cases, the remaining ones follow analogously. Choose $`p=e_{(1,1)(1,1)}`$, $`p^{}`$ as the generators of $`๐ต((\epsilon ))`$. We have two possibilities. Namely, there exists a even state $`\phi `$ on $`๐_{\{1\}}`$, and a even state $`\psi `$ either on $`๐_{\{1\}}`$, or on $`(p^{}๐_{\{0\}}p^{})๐_{\{1\}}`$ such that, for $`x๐_{\{0\}}`$, $`y๐_{\{1\}}`$,
$$\epsilon (xy)=Tr(xp)\phi (y)p+\psi (y)p^{}xp^{},$$
respectively
$$\epsilon (xy)=Tr(xp)\phi (y)p+\psi (p^{}xp^{}y)p^{}.$$
### 6.8. Case 8:
$`๐ต((\epsilon ))^2`$, $`\mathrm{\#}`$ of orbits of $`\mathrm{\Theta }_{๐ต((\epsilon ))}=1`$.
Choose $`\{Q_{\pm \chi }\}`$ as the generators of $`๐ต((\epsilon ))`$. We have two possibilities. First
$$\epsilon (x)=\phi (Q_\chi xQ_\chi )Q_\chi +\phi (Q_\chi \mathrm{\Theta }(x)Q_\chi )Q_\chi ,x๐_{[0,1]},$$
$`\phi `$ being a state on $`Q_\chi ๐_{[0,1]}Q_\chi `$. Second, let $`๐
๐_{[0,1]}`$ be the tensor completion of $`๐_{\{0\}}`$ in $`๐_{[0,1]}`$.<sup>8</sup><sup>8</sup>8According to (2.3), the subalgebra $`๐
`$ is obtained by constructing a systems $`\{e_{kl}(j),f_{kl}(j)|j,k,l=1,2\}`$ of four mutually commuting $`2\times 2`$ matrix units for $`๐_{[0,1]}`$. Notice that $`๐
`$ is localized in the whole $`๐_{[0,1]}`$, and is $`\mathrm{\Theta }`$โinvariant. Then there exists a state $`\phi `$ on $`๐
`$ such that for $`x๐_{\{0\}}`$, $`y๐
`$,
$$\epsilon (xy)=\phi (y)Q_\chi xQ_\chi +\phi (\mathrm{\Theta }(y))Q_\chi xQ_\chi ).$$
### 6.9. Case 9:
$`๐ต((\epsilon ))`$.
We treat only one possibility, the two remaining ones generating one step product states (see e.g. Subsection 6.1 of ). Let $`N`$, $`\overline{N}`$ be the algebra generated by $`a_1,a_1^+`$, $`a_2,a_2^+`$ respectively. Then there exists an even state $`\phi `$ on $`\overline{N}๐_{\{1\}}`$ such that for $`xN`$, $`y\overline{N}๐_{\{1\}}`$,
$$\epsilon (xy)=\phi (y)x.$$
Notice that this example is nothing but that the two block factor treated in Subsection 6.3 of . This is easily seen by passing in , to the two point regrouped algebra.
### 6.10. Case 10:
two examples with $`๐_{\{n\}}๐_{2^3}()`$.
We describe two examples relative to more complicated situations than the previous ones. Let $`\{a_i,a_i^+|i=1,2,3\}`$, $`\{b_i,b_i^+|i=1,2,3\}`$ be the generators of $`๐_{\{0\}}`$, $`๐_{\{1\}}`$ respectively. Let $`\{e_{kl}(j),f_{kl}(j)|j,k,l=1,2\}`$ of commuting $`2\times 2`$ matrix units obtained according to (2.3), and realizing the isomorphism $`๐_{[0,1]}\underset{6\text{โtimes}}{\underset{}{๐_2()\mathrm{}๐_2()}}`$. Put for $`\chi ๐`$,
$$P_\chi :=\frac{1}{2}\left(1\mathrm{I}+\chi a_1+\overline{\chi }a_1^+\right).$$
First define $`N_i`$, $`\overline{N}_i`$ as the algebras generated by $`\{e_{ii}(1)a_2,e_{ii}(1)a_2^+\}`$, $`\{e_{ii}(1)a_3,e_{ii}(1)a_3^+\}`$, $`i=1,2`$ respectively. Choose even states $`\phi _i`$ on $`\overline{N}_i๐_{\{1\}}`$. Then for $`x_iN_i`$, $`y\overline{N}_i๐_{\{1\}}`$,
$$\epsilon \left(\underset{i=1}{\overset{2}{}}x_iy_i\right)=\underset{i=1}{\overset{2}{}}\phi _i(y_i)x_i.$$
Second define $`N_\chi `$, $`M_\chi `$ as the algebras generated by $`\{P_\chi e_{ij}(2)|i,j=1,2\}`$, $`\{P_\chi e_{ij}(3)f_{kl}(n)|i,j,k,l=1,2,n=1,2,3\}`$ respectively. Choose a state $`\phi `$ on $`M_\chi `$. Then for $`x_{\pm \chi }N_{\pm \chi }`$, $`y_{\pm \chi }M_{\pm \chi }`$,
$$\epsilon (x_\chi y_\chi +x_\chi y_\chi )=\phi (y_\chi )x_\chi +\phi (\mathrm{\Theta }(y_\chi ))x_\chi .$$
## acknowledgements
The author is grateful to L. Accardi for suggesting the themes and for several useful discussions. He also acknowledges one of the referees whose suggestions considerably contributed to improve the presentation of the present work. |
warning/0506/astro-ph0506420.html | ar5iv | text | # Spectral Synthesis of SDSS Galaxies
## 1. Introduction
In the proceedings of the meeting on Stellar Populations held in Baltimore in 1986 Searle (1986), discussing integral light spectral synthesis, states that โthis subject has a bad reputation. Too much has been claimed, and too few have been persuadedโ. Indeed, recovering the stellar content of a galaxy from its observed integrated spectrum is not an easy task, as can be deduced from the amount of work devoted to this topic over the past half century. The situation is however much more favorable nowadays. Huge observational and theoretical efforts in the past few years have produced large sets of high quality spectra of stars (e.g., Prugniel & Soubiran 2001; Le Borgne et al. 2003; Bertone et al. 2004; Gonzรกlez-Delgado et al. 2005). These libraries are being implemented in a new generation of evolutionary synthesis models, allowing the prediction of galaxy spectra with an unprecedented level of detail (Vazdekis 1999; Bruzual & Charlot 2003, hereafter BC03; Le Borgne et al. 2004). At the same time, galaxy spectra are now more abundant than ever. The Sloan Digital Sky Survey (SDSS), in particular, is providing a homogeneous data base of hundreds of thousands of galaxy spectra in the 3800โ9200 ร
range, with a resolution of $`\lambda /\mathrm{\Delta }\lambda 1800`$ (York et al. 2000; Stoughton et al. 2002; Abazajian et al. 2003, 2004). This enormous amount of high quality data will undoubtedly be at the heart of tremendous progress in our understanding of galaxy constitution, formation and evolution. Actually, galaxy spectra encode information on the age and metallicity distributions of the constituent stars, which in turn reflect the star-formation and chemical histories of the galaxies. Retrieving this information from observational data in a reliable way is crucial for a deeper understanding of galaxy formation and evolution and, in fact, significant steps in this direction have recently been made (Kauffmann et al. 2003a, hereafter K03; Brinchmann et al. 2004; Tremonti et al. 2004; Heavens et al. 2004; Panter, Heavens & Jimenez 2004).
Here we demonstrate that, besides providing excellent starlight templates to aid emission line studies, spectral synthesis recovers reliable stellar population properties out of galaxy spectra of realistic quality. We show that this simple method provides robust information on the stellar age ($`t_{}`$) and stellar metallicities ($`Z_{}`$) distributions, as well as on the extinction, velocity dispersion and stellar mass. The ability to recover information on $`Z_{}`$ is particularly welcome, given that stellar metallicities are notoriously more difficult to assess than other properties. In order to reach this goal we follow: (1) a priori arguments, based on simulations; (2) comparisons with independent work based on a different method, and (3) an a posteriori empirical analysis of the consistency of results obtained for a large sample of SDSS galaxies. Most of the results presented here are discussed in detail in Cid Fernandes et al. (2005).
In Section 2 we present an overview of our synthesis method and simulations designed to test it and evaluate the uncertainties involved. The discussion is focused on how to use the synthesis to derive robust estimators of physically interesting stellar population properties. Section 3 defines a volume limited sample of SDSS galaxies and presents the results of the synthesis of their spectra, along with measurements of emission lines. We also compare in this section our results to those obtained by other authors. Stellar population and emission line properties are used in Section 4 to investigate whether the synthesis produces astrophysically plausible results. In Section 5 we present preliminary results of an analysis of a magnitude-limited sample of 20000 galaxies also extracted from SDSS. Finally, Section 6 summarizes our main findings.
## 2. Spectral Synthesis
### 2.1. Method
Our synthesis code, which we call STARLIGHT, was first discussed in Cid Fernandes et al. (2004, hereafter CF04). We fit an observed spectrum $`O_\lambda `$ with a combination of $`N_{}`$ Simple Stellar Populations (SSP) from the evolutionary synthesis models of BC03. Extinction is modeled as due to foreground dust, and parametrized by the V-band extinction $`A_V`$. The Galactic extinction law of Cardelli, Clayton & Mathis (1989) with $`R_V=3.1`$ is adopted. Line of sight stellar motions are modeled by a Gaussian distribution $`G`$ centered at velocity $`v_{}`$ and with dispersion $`\sigma _{}`$. With these assumptions the model spectrum is given by
$$M_\lambda =M_{\lambda _0}\left[\underset{j=1}{\overset{N_{}}{}}x_jb_{j,\lambda }r_\lambda \right]G(v_{},\sigma _{})$$
(1)
where $`b_{j,\lambda }`$ is the spectrum of the $`j^{\mathrm{th}}`$ SSP normalized at $`\lambda _0`$, $`r_\lambda 10^{0.4(A_\lambda A_{\lambda _0})}`$ is the reddening term, $`M_{\lambda _0}`$ is the synthetic flux at the normalization wavelength, $`\stackrel{}{x}`$ is the population vector and $``$ denotes the convolution operator. Each component $`x_j`$ ($`j=1\mathrm{}N_{}`$) represents the fractional contribution of the SSP with age $`t_j`$ and metallicity<sup>1</sup><sup>1</sup>1In this paper we follow the convention used in stellar evolution studies, which define stellar metallicities in terms of the fraction of mass in metals. In this system the Sun has $`Z_{}=0.02`$. $`Z_j`$ to the model flux at $`\lambda _0`$. The base components can be equivalently expressed as a mass fractions vector $`\stackrel{}{\mu }`$. In this work we adopt a base with $`N_{}=45`$ SSPs, encompassing 15 ages between $`10^6`$ and $`1.3\times 10^{10}`$ yr and 3 metallicities: $`Z=0.2`$, 1 and 2.5 $`Z_{}`$. Their spectra were computed with the STELIB library (Le Borgne et al. 2003), Padova 1994 tracks, and Chabrier (2003) IMF (see BC03 for details).
The fit is carried out with a simulated annealing plus Metropolis scheme which searches for the minimum $`\chi ^2=_\lambda \left[\left(O_\lambda M_\lambda \right)w_\lambda \right]^2`$, where $`w_\lambda ^1`$ is the error in $`O_\lambda `$. Regions around emission lines, bad pixels or sky residuals are masked out by setting $`w_\lambda =0`$. Pixels which deviate by more than 3 times the rms between $`O_\lambda `$ and an initial estimate of $`M_\lambda `$ are also given zero weight.
Fig. 1 illustrates the spectral fit obtained for a galaxy drawn from the SDSS database. The top-left panel shows the observed spectrum (thin line) and the model (thick), as well as the error spectrum (dashed). The bottom-left panel shows the $`O_\lambda M_\lambda `$ residual spectrum, while the panels in the right summarize the derived star-formation history encoded in the age-binned population vector. This example, along with those in K03 and CF04, demonstrates that this simple method is capable of reproducing real galaxy spectra to an excellent degree of accuracy.
An important application of the synthesis is to measure emission lines from the residual spectrum, as done by K03. Another, of course, is to infer stellar population properties from the fit parameters. Our central goal here is to investigate whether spectral synthesis can also recover reliable stellar population properties. In the remainder of this section we address this issue by means of simulations.
### 2.2. Robust description of the synthesis results
The existence of multiple solutions is an old known problem in stellar population synthesis, and even superb spectral fits as that shown in Fig. 1 do not guarantee that the resulting parameter estimates are trustworthy. There is, consequently, a need to assess the degree to which one can trust the parameters involved in the fit before using them to infer stellar populations properties. The spectral fits involve $`N_{}+3`$ parameters: $`N_{}1`$ of the $`\stackrel{}{x}`$ components (one degree of freedom is removed by the normalization constraint), $`M_{\lambda _0}`$, $`A_V`$ and the two kinematical parameters, $`v_{}`$ and $`\sigma _{}`$. The reliability of parameter estimation is best studied by means of simulations which feed the code with spectra generated with known parameters, add noise, and then examine the correspondence between input and output values.
#### Simulations
We have carried out simulations designed to test the method and investigate which combinations of the parameters provide robust results. Several sets of simulations were performed. Given our interest in modeling SDSS galaxies, here we focus on simulations tailored to match the characteristics of this data set. Test galaxies were built from the average $`\stackrel{}{x}`$, $`A_V`$ and $`\sigma _{}`$ within 65 boxes in the mean stellar age versus mean stellar metallicity plane obtained for the sample described in Section 3. These new simulations confirm previous results reported in CF04 that the individual components of $`\stackrel{}{x}`$ are very uncertain, so we skip a detailed comparison between $`\stackrel{}{x}_{\mathrm{input}}`$ and $`\stackrel{}{x}_{\mathrm{output}}`$ and jump straight to results based on more robust descriptions of the synthesis output.
#### Condensed population vector
A coarse but robust description of the star-formation history of a galaxy may be obtained by binning $`\stackrel{}{x}`$ onto โYoungโ ($`t_j<10^8`$ yr), โIntermediate-ageโ ($`10^8t_j10^9`$ yr), and โOldโ ($`t_j>10^9`$ yr) components ($`x_Y`$, $`x_I`$ and $`x_O`$ respectively). These age-ranges were defined on the basis of the simulations, by seeking which combinations of $`x_j`$โs produce smaller input $``$ output residuals. Our simulations show that these 3 components are very well recovered by the method, with uncertainties smaller than $`\mathrm{\Delta }x_Y=0.05`$, $`\mathrm{\Delta }x_I=0.1`$, and $`\mathrm{\Delta }x_O=0.1`$ for $`S/N10`$.
#### Mass, extinction and velocity dispersion
We have also analyzed the input versus output values of $`A_V`$, $`\sigma _{}`$ and the stellar mass $`M_{}`$. The latter is not an explicit input parameter of the models, but may be computed from $`\stackrel{}{\mu }`$ and the $`M_{}/L_{\lambda _0}`$ ratio of the different populations in the base. The uncertainties in the recovery of these parameters are $`\mathrm{\Delta }A_V<0.05`$ mag, $`\mathrm{\Delta }\mathrm{log}M_{}<0.1`$ dex and $`\mathrm{\Delta }\sigma _{}<12`$ km s<sup>-1</sup> for $`S/N10`$.
#### Mean stellar age
If one had to choose a single parameter to characterize the stellar population mixture of a galaxy, the option would certainly be its mean age. We define two versions of mean stellar age (the logarithm of the age, actually), one weighted by light
$$<\mathrm{log}t_{}>_L=\underset{j=1}{\overset{N_{}}{}}x_j\mathrm{log}t_j$$
(2)
and another weighted by stellar mass
$$<\mathrm{log}t_{}>_M=\underset{j=1}{\overset{N_{}}{}}\mu _j\mathrm{log}t_j$$
(3)
Note that, by construction, both definitions are limited to the 1 Myrโ13 Gyr range spanned by the base. The mass weighted mean age is in principle more physical, but, because of the non-constant $`M/L`$ of stars, it has a much less direct relation with the observed spectrum than $`<\mathrm{log}t_{}>_L`$.
The simulations show that the mean age is a very robust quantity. The rms difference between input and output $`<\mathrm{log}t_{}>_L`$ values is $`0.08`$ dex for $`S/N>10`$, and $`0.14`$ dex for $`<\mathrm{log}t_{}>_M`$. Although the uncertainties of $`<\mathrm{log}t_{}>_L`$ and $`<\mathrm{log}t_{}>_M`$ are comparable in absolute terms, the latter index spans a smaller dynamical range (because of the large $`M/L`$ ratio of old populations), so in practice $`<\mathrm{log}t_{}>_L`$ is the more useful of the two indices.
#### Mean stellar metallicity
Given an option of what to choose as a second parameter to describe a mixed stellar population, the choice would likely be its typical metallicity. Analogously to what we did for ages, we define both light and mass-weighted mean stellar metallicities:
$$<Z_{}>_L=\underset{j=1}{\overset{N_{}}{}}x_jZ_j$$
(4)
and
$$<Z_{}>_M=\underset{j=1}{\overset{N_{}}{}}\mu _jZ_j$$
(5)
both of which are bounded by the 0.2โ2.5 $`Z_{}`$ base limits. Our simulations show that the rms of $`\mathrm{\Delta }\mathrm{log}<Z_{}>_M=\mathrm{log}<Z_{}>_{M,\mathrm{output}}\mathrm{log}<Z_{}>_{M,\mathrm{input}}`$ is of order 0.1 dex. In absolute terms this is comparable to $`\mathrm{\Delta }<\mathrm{log}t_{}>`$, but note that $`<Z_{}>`$ covers a much narrower dynamical range than $`<\mathrm{log}t_{}>`$, so that in practice mean stellar metallicities are more sensitive to errors than mean ages. This is not surprising, given that age is the main driver of variance among SSP spectra, metallicity having a โsecond-orderโ effect (e.g., Schmidt et al. 1991; Ronen, Aragon-Salamanca & Lahav 1999). This is the reason why studies of the stellar populations of galaxies have a much harder time estimating metallicities than ages, to the point that one is often forced to bin-over the $`Z`$ information and deal only with age-related estimates such as $`<\mathrm{log}t_{}>`$ (e.g., Cid Fernandes et al. 2001; Cid Fernandes, Leรฃo & Rodrigues Lacerda 2003; K03).
Notwithstanding these notes, it is clear that uncertainties of $`0.1`$ dex in $`<Z_{}>`$ are actually good news, since they do allow us to recover useful information on an important but hard to measure property. This new tracer of stellar metallicity is best applicable to large samples of galaxies such as the SDSS. The statistics of samples help reducing uncertainties associated with $`<Z_{}>`$ estimates for single objects and allows one to investigate correlations between $`<Z_{}>`$ and other galaxy properties (Sodrรฉ et al. , in preparation; Section 4.1.).
#### Age-Metallicity degeneracy
Our method tends to underestimate $`<Z_{}>`$ for metal-rich systems and vice-versa. This trend is due to the infamous age-metallicity degeneracy. In order to verify to which degree our synthesis is affected by this well known problem (e.g., Renzini & Buzzoni 1986; Worthey 1994; Bressan, Chiosi & Tantalo 1996) we have examined the correlation between the output minus input residuals in $`<\mathrm{log}t_{}>_L`$ and $`\mathrm{log}<Z_{}>_L`$. The age-$`Z`$ degeneracy acts in the sense of confusing old, metal-poor systems with young, metal-rich ones and vice-versa, which should produce anti-correlated residuals. This anti-correlation is also present with the mass-weighted mean stellar metallicity $`<Z_{}>_M`$, but is not as strong as for $`<Z_{}>_L`$. On the other hand, the uncertainty in $`<Z_{}>_M`$ is always larger than for $`<Z_{}>_L`$.
The age-$`Z`$ degeneracy is thus present in our method, introducing correlated residuals in our $`<Z_{}>`$ and $`<\mathrm{log}t_{}>`$ estimates at the level of up to $`0.1`$โ0.2 dex. However, none of the results reported here rely on this level of precision.
#### Base limitations
We have carried out simulations using the $`Z=0.02Z_{}`$ BC03 SSPs to examine some of the limitations of the base addopted here. Galaxies with such low metallicity are not expected to be present in significant numbers in the sample described in Section 3.1., given that it excludes low luminosity systems like HII galaxies and dwarf ellipticals (which are also the least metallic ones by virtue of the mass-metallicity relation). Still, it is interesting to investigate what would happen in this case. When synthesized with our 0.2โ$`2.5Z_{}`$ base, these extremely metal poor galaxies are modeled predominantly with the $`0.2Z_{}`$ components, as intuitively expected. Moreover, the mismatch in metallicity introduces non-negligible biases in other properties, like masses, mean ages and extinction ($`M_{}`$, for instance, is systematically underestimated by 0.3 dex). Similar problems should be encountered when modeling systems with $`Z>2.5Z_{}`$. These results serve as a reminder that our base spans a wide but finite range in stellar metallicity, and that extrapolating these limits has an impact on the derived physical properties. While there is no straightforward a priori diagnostic of which galaxies violate these limits, in general, one should be suspicious of objects with mean $`Z_{}`$ too close to the base limits.
Our simulations demonstrate that we are able of producing reliable estimates of several parameters of astrophysical interest, at least in principle. We must nevertheless emphasize that this conclusion relies entirely on models and on an admittedly simplistic view of galaxies. When applying the synthesis to real galaxy spectra, a series of other effects come into play. For instance, the extinction law appropriate for each galaxy likely differs from the one used here. Similarly, while in the evolutionary tracks adopted here the metal abundances are scaled from the solar values, non-solar abundance mixtures are known to occur in stellar systems (e.g., Trager et al. 2000a,b; see also Section 5), not to mention uncertainties in the SSP models and the always present issue of the IMF. Accounting for all these effects in a consistent way is not currently feasible. We mention these caveats not to dismiss simple models, but to highlight that all parameter uncertainties discussed above are applicable within the scope of the model.
Hence, while the simulations lend confidence to the synthesis method, one might remain skeptical of its actual power. The next sections further address the reliability of the synthesis, this time from a more empirical perspective.
## 3. Analysis of a volume-limited galaxy sample
In this section we apply our synthesis method to a large sample of SDSS galaxies to estimate their stellar population properties. We also present measurements of emission line properties, obtained from the observed minus synthetic spectra. The information provided by the synthesis of so many galaxies allows one to address a long menu of astrophysical issues related to galaxy formation and evolution. Before venturing in the exploration of such issues, however, it is important to validate the results of the synthesis by as many means as possible. Hence, the goal of this study is not so much to explore the physics of galaxies but to provide an empirical test of our synthesis method.
### 3.1. Sample definition
The spectroscopic data used in this work were taken from the SDSS. This survey provides spectra of objects in a large wavelength range (3800โ9200 ร
) with mean spectral resolution $`\lambda /\mathrm{\Delta }\lambda 1800`$, taken with 3 arcsec diameter fibers. The most relevant characteristic of this survey for our study is the enormous amount of good quality, homogeneously obtained spectra. The data analyzed here were extracted from the SDSS main galaxy sample available in the Data Release 2 (DR2; Abazajian et al. 2004). This flux-limited sample consists of galaxies with reddening-corrected Petrosian $`r`$-band magnitudes $`r17.77`$, and Petrosian $`r`$-band half-light surface brightnesses $`\mu _{50}24.5`$ mag arcsec<sup>-2</sup> (Strauss et al. 2002).
From the main sample, we first selected spectra with a redshift confidence $`0.35`$. Following the conclusions of Zaritsky, Zabludoff, & Willick (1995), we have imposed a redshift limit of $`z>0.05`$ (trying to avoid aperture effects and biases; see e.g. Gรณmez et al. 2003) and selected a volume limited sample up to $`z=0.1`$, corresponding to a $`r`$-band absolute magnitude limit of $`M(r)=20.5`$. The absolute magnitudes used here are k-corrected with the help of the code provided by Blanton et al. (2003; kcorrect v3\_2) and assuming the following cosmological parameters: $`H_0`$ = 70 km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_M=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$. We also restricted our sample to objects for which the observed spectra show a $`S/N`$ ratio in $`g`$, $`r`$ and $`i`$ bands greater than 5. These restrictions leave us with a volume limited sample containing $`50362`$ galaxies, which leads to a completeness level of $``$ 98.5 per cent.
### 3.2. Results of the spectral synthesis
All 50362 spectra were brought to the rest-frame (using the redshifts in the SDSS database), sampled from 3650 to 8000 ร
in steps of 1 ร
, corrected for Galactic extinction<sup>2</sup><sup>2</sup>2Unlike in the first data release, the final calibrated spectra from the DR2 are not corrected for foreground Galactic reddening. using the maps given by Schlegel, Finkbeiner & Davis (1998) and the extinction law of Cardelli et al. (1989, with $`R_V=3.1)`$, and normalized by the median flux in the 4010โ4060 ร
region. The $`S/N`$ ratio in this spectral window spans the 5โ30 range, with median value of 14. Besides the masks around the lines listed in Section 2.2., we exclude points with SDSS flag $`2`$, which signals bad pixels, sky residuals and other artifacts. After this pre-processing, the spectra are fed into the STARLIGHT code described in Section 2.1.. On average, the synthesis is performed with $`N_\lambda =3677`$ points, after discounting the ones which are clipped by our $`3`$ sigma threshold (typically 40 points) and the masked ones. The spectral fits are generally very good.
The total stellar masses of the galaxies were obtained from the stellar masses derived from the spectral synthesis (which correspond to the light entering the fibers) by dividing them by $`(1f)`$, where $`f`$ is the fraction of the total galaxy luminosity in the $`z`$-band outside the fiber. This approach, which neglects stellar population and extinction gradients, leads to an increase of typically 0.5 in $`\mathrm{log}M_{}`$. We did not apply any correction to the velocity dispersion estimated by the code given that the spectral resolution of the BC03 models and the data are very similar.
We point out that we did not constrain the extinction $`A_V`$ to be positive. There are several reasons for this choice: (a) some objects may be excessively dereddened by Galactic extinction; (b) some objects may indeed require bluer SSP spectra than those in the base; (c) the observed light may contain a scattered component, which would induce a bluening of the spectra not taken into account by the adopted pure extinction law; (d) constraining $`A_V`$ to have only positive values produces an artificial concentration of solutions at $`A_V=0`$, an unpleasant feature in the $`A_V`$ distribution. Interestingly, most of the objects for which we derive negative $`A_V`$ (typically $`0.1`$ to $`0.3`$ mag) are early-type galaxies. These galaxies are dominated by old populations, and expected to contain little dust. This is consistent with the result of K03, who find negative extinction primarily in galaxies with a large $`D_n(4000)`$. The distribution of $`A_V`$ for these objects, which can be selected on the basis of spectral or morphological properties, is strongly peaked around $`A_V=0`$, so that objects with $`A_V<0`$ can be considered as consistent with having zero extinction. In any case, none of the results reported in this paper is significantly affected by this choice.
### 3.3. Emission line measurements
Besides providing estimates of stellar population properties, the synthesis models allow the measurement of emission lines from the โpure-emissionโ, starlight subtracted spectra $`(O_\lambda M_\lambda )`$. We have measured the lines of \[O ii\]$`\lambda \lambda `$3726,3729, \[O iii\]$`\lambda `$4363, H$`\beta `$, \[O iii\]$`\lambda \lambda `$4959,5007, \[O i\]$`\lambda `$6300, \[N ii\]$`\lambda `$6548, H$`\alpha `$, \[N ii\]$`\lambda `$6584 and \[S ii\]$`\lambda \lambda `$6717,6731. Each line was treated as a Gaussian with three parameters: width, offset (with respect to the rest-frame central wavelength), and flux. Lines from the same ion were assumed to have the same width and offset. We have further imposed \[O iii\]$`\lambda 5007`$/\[O iii\]$`\lambda 4959=2.97`$ and \[N ii\]$`\lambda 6584`$/\[N ii\]$`\lambda 6548=3`$ flux ratio constraints. Finally, we consider a line to have significant emission if its fit presents a $`S/N`$ ratio greater than 3.
In some of the following analysis, galaxies with emission lines are classified according to their position in the \[O iii\]/H$`\beta `$ versus \[N ii\]/H$`\alpha `$ diagram proposed by Baldwin, Phillips & Terlevich (1981) to distinguish normal star-forming galaxies from galaxies containing active galactic nuclei (AGN). We define as normal star-forming galaxies those galaxies that appear in this diagram and are below the curve defined by K03 (see also Brinchmann et al. 2004). Objects above this curve are transition objects and galaxies containing AGN.
### 3.4. Comparisons with the MPA/JHU database
The SDSS database has been explored by several groups, using different approaches and techniques. The MPA/JHU group has recently publicly released catalogues<sup>3</sup><sup>3</sup>3available at http://www.mpa-garching.mpg.de/SDSS/ of derived physical properties for 211894 SDSS galaxies, including 33589 narrow-line AGN (K03, see also Brinchmann et al. 2004). These catalogues are based on the K03 method to infer the star formation histories, dust attenuation and stellar masses of galaxies from the simultaneous analysis of the 4000 ร
break strength, $`D_n(4000)`$, and the Balmer line absorption index H$`\delta _A`$. These two indices are used to constrain the mean stellar ages of galaxies and the fractional stellar mass formed in bursts over the past few Gyr, and a comparison with broad-band photometry then allows to estimate the extinction and stellar masses.
The MPA/JHU catalogues provide very useful benchmarks for similar studies. In this section we summarize the comparison between values of some of the parameters from these catalogues with our own estimates. A more complete discussion is presented in Cid Fernandes et al. (2005).
#### Stellar extinction
The MPA/JHU group estimates the $`z`$-band stellar extinction $`A_z`$ through the difference between model and measured colours, assuming an attenuation curve proportional to $`\lambda ^{0.7}`$. In our case, the extinction $`A_V`$ is derived directly from the spectral fitting, carried out with the Milky Way extinction law (Cardelli et al. 1989, c.f. Section 2.1.). These two independent estimates are very strongly and linearly correlated, with a Spearman rank correlation coefficient $`r_S=0.95`$. However, the values of $`A_z`$ reported by the MPA/JHU group are systematically larger than our values: $`A_z(\mathrm{MPA}/\mathrm{JHU})2.51A_z(\mathrm{This}\mathrm{Work})`$ in the median. This discrepancy, however, is only apparent because the Galactic extinction law is substantially harder than $`\lambda ^{0.7}`$. One thus expects to need less extinction when modeling a given galaxy with the former law than with the latter. Our analysis indicates that there are no substantial differences between the MPA/JHU and our estimates of the stellar extinction other than those implied by the differences in the reddening laws adopted in the two studies.
#### Stellar masses
Our results for the total stellar masses compares very well to the MPA/JHU extinction-corrected stellar masses, with $`r_S=0.89`$. The quantitative agreement is also good, with a median difference of just 0.1 dex. This small offset seems to be due to a subtle technicality. Whereas we adopt the $`M/L`$ ratio of the best $`\chi ^2`$ model, the MPA/JHU group derives $`M/L`$ comparing the observed values of the $`D_n(4000)`$ and $`H\delta _A`$ indices with a library of 32000 models. Each model is then weighted by its likelihood, and a probability distribution for $`M/L`$ is computed. The MPA/JHU mass is the median of this distribution, which is not necessarily the same as the best-$`\chi ^2`$ value.
#### Velocity dispersion
We use the sub-sample of galaxies with active nuclei to compare our measurements of absorption line broadening due to galaxy velocity dispersion and/or rotation, $`\sigma _{}`$, with those of the MPA/JHU group, since they list this quantity only in their AGN Catalogue. The Spearman correlation-coefficient in this case is $`r_S=0.91`$ and the median of the difference between the two estimates is just 9 km s<sup>-1</sup>, indicating an excellent agreement between both studies.
#### Emission lines and nebular metallicities
Brinchmann et al. (2004) also provide, in their Emission line Catalogue, data on emission lines which can be compared with our own measurements. We have compared the fluxes and equivalent widths of H$`\alpha `$, \[N ii\]$`\lambda 6584`$, \[O ii\]$`\lambda 3727`$, H$`\beta `$ and \[O iii\]$`\lambda 5007`$ as measured by our code and that obtained by the MPA/JHU group. We do not find any significant difference between these values; the largest discrepancy ($``$ 5 per cent) was found for the equivalent widths of H$`\alpha `$ and \[N ii\], probably due to different estimates of the continuum level and the associated underlying stellar absorption. Our emission line measurements are also in good agreement with those in Stasiลska et al. (2004), who fit the Balmer lines with emission and absorption components, instead of subtracting a starlight model.
Overall, we conclude that our spectral synthesis method yields estimates of physical parameters in good agreement with those obtained by the MPA/JHU group, considering the important differences in approach and underlying assumptions.
## 4. Empirical relations
Yet another way to assess the validity of physical properties derived through a spectral synthesis analysis is to investigate whether this method yields astrophysically reasonable results. In this section we follow this empirical line of reasoning by comparing some results obtained from our synthesis of SDSS galaxies (which excludes emission lines) with those obtained from a direct analysis of the emission lines. Our aim here is to demonstrate that our synthesis results do indeed make sense.
### 4.1. Stellar and Nebular Metallicities
Our spectral synthesis approach yields estimates of the mean metallicity of the stars in a galaxy, $`<Z_{}>`$. The analysis of emission lines, on the other hand, gives estimates of the present-day abundances in the warm interstellar medium. Although stellar and nebular metallicities are not expected to be equal, it is reasonable to expect that they should roughly scale with each other.
Fig. 2 shows the correlation between mass-weighted stellar metallicities and the nebular oxygen abundance (computed as in Section 3.4.), both in solar units<sup>4</sup><sup>4</sup>4The solar unit adopted for the nebular oxygen abundance is $`12+\mathrm{log}(\mathrm{O}/\mathrm{H})_{}=8.69`$ (Allende Prieto, Lambert & Apslund 2001)., for our sample of normal star-forming galaxies. A correlation is clearly seen, although with large scatter ($`r_S=0.42`$). Galaxies with large stellar metallicities also have large nebular oxygen abundances; galaxies with low stellar abundances tend to have smaller abundances. The observed scatter is qualitatively expected due to variations in enrichment histories among galaxies.
Nebular and stellar metallicities are estimated through completely different and independent methods, so the correlation depicted in Fig. 2 provides an a posteriori empirical validation for the stellar metallicity derived by the spectral synthesis. The possibility to estimate stellar metallicities for so many galaxies is one of the major virtues of spectral synthesis, as it opens an important window to study the chemical evolution of galaxies and of the universe as a whole (Panter et al. 2004; Sodrรฉ et al. in prep.).
### 4.2. Stellar and Nebular extinctions
The stellar extinction in the V-band is one of the products of our STARLIGHT code. A more traditional and completely independent method to evaluate the extinction consists of comparing the observed H$`\alpha `$/H$`\beta `$ Balmer decrement to the theoretical value, the โBalmer extinctionโ (e.g., Stasiลska et al. 2004).
Fig. 3 presents a comparison between the stellar $`A_V`$ and $`A_V^{\mathrm{Balmer}}`$. These two extinctions are determined in completely independent ways, and yet, our results show that they are closely linked, with $`r_S=0.61`$. A linear bisector fitting yields $`A_V^{\mathrm{Balmer}}=0.24+1.81A_V`$. Note that the angular coefficient in this relation indicates that nebular photons are roughly twice as extincted as the starlight. This โdifferential extinctionโ is in very good qualitative and quantitative agreement with empirical studies (Fanelli et al. 1988; Calzetti, Kinney & Storchi-Bergmann 1994; Gordon et al. 1997; Mas-Hesse & Kunth 1999).
### 4.3. Relations with mean stellar age
The equivalent width (EW) of H$`\alpha `$ is related to the ratio of present to past star formation rate of a galaxy (e.g., Kennicutt 1998). It is thus expected to be smaller for older galaxies. Fig. 4a shows the relation between EW(H$`\alpha `$) and the mean light-weighted stellar age obtained by our spectral synthesis. The anti-correlation, which has $`r_S=0.78`$, is evident. It is worth stressing that these two quantities are obtained independently, since the spectral synthesis does not include emission lines.
Another quantity that is considered a good age indicator, even for galaxies without emission lines, is the 4000 ร
break, D4000. We measured this index following Bruzual (1983), who define D4000 as the ratio between the average value of $`F_\nu `$ in the 4050โ4250 and 3750โ3950 ร
bands. The relation between $`<\mathrm{log}t_{}>_L`$ and D4000 is shown in Fig. 4b. Note that the concentration of points at the high age end reflects the upper age limit of the base adopted here, 13 Gyr (c.f. Section 2.2). The correlation is very strong ($`r_S=0.94`$), showing that indeed D4000 can be used to estimate empirically mean light-weighted galaxy ages, despite its metallicity dependence for very old stellar populations (older than 1 Gyr, as shown by K03).
### 4.4. $`M_{}`$$`\sigma _{}`$ relation
Fig. 5 shows the relation between stellar mass and $`\sigma _{}`$, obtained from the synthesis. The relation is quite good, with $`r_S=0.79`$. The solid line displayed in the figure is $`\mathrm{log}M_{}=6.44+2.04\mathrm{log}\sigma _{}`$ for $`M_{}`$ in $`M_{}`$ and $`\sigma _{}`$ in km s<sup>-1</sup>, obtained with a bisector fitting. The figure also shows as a dashed line a fit assuming $`M_{}\sigma _{}^4`$, expected from the virial theorem under the (unrealistic) assumption of constant mass surface density. In both cases we have excluded from the fit galaxies with $`\sigma _{}<35`$ km s<sup>-1</sup>, which corresponds to less than half the spectral resolution of both data and models.
This is another relation that is expected a priori if we have in mind the Faber-Jackson relation for ellipticals and the Tully-Fisher relation for spirals. For early-type galaxies, $`\sigma _{}`$ is a measure of the central velocity dispersion, which is directly linked to the gravitational potential depth, and, through the virial theorem, to galactic mass. For late-type systems, $`\sigma _{}`$ has contributions of isotropic motions in the bulges, as well as of the rotation of the disks, and is also expected to relate with galactic mass. Another aspect that it is interesting to point out in Fig. 5 is that the dispersion in the $`M_{}`$$`\sigma _{}`$ relation decreases as we go from low-luminosity, rotation-dominated systems, for which the values of $`\sigma _{}`$ depend on galaxy inclination and bulge-to-disk ratio, to high-luminosity, mostly early-type systems, which obey a much more regular (and steeper) relation between $`\sigma _{}`$ and $`M_{}`$.
This relation, between a quantity that is not directly linked to the synthesis, $`\sigma _{}`$, and another one that is a product of our synthesis, $`M_{}`$, is yet another indication that the results of our STARLIGHT code do make sense.
## 5. Recent results
We have started an analysis of a magnitude-limited sample with 20000 galaxies extracted from DR2, aiming to probe galaxies with luminosities smaller than those discussed in the previous sections. We present below a summary of some of the results obtained so far.
#### The bimodality of galaxy populations
With the recent advance of redshift surveys, the study of galaxy populations has quantitatively revealed the existence of a bimodal distribution in some fundamental galaxy properties, found both in photometric and spectroscopic data: star-forming and passive galaxies have quite distinct properties. Perhaps the most representative bimodal distribution is that found in galaxy colours. Since photometric parameters are easily measured, their bimodal behavior has been studied for galaxies in the local universe by using both the SDSS and the Two Degree Field Galaxy Redshift Survey (2dFGRS) data (Strateva et al. 2001; Hogg et al. 2002; Blanton et al. 2003, Wild et al. 2004). A bimodality has been found in other galaxy parameters, like mass (Kauffmann et al. 2003b) and star formation properties (Madgwick et al. 2002; Wild et al. 2004; Brinchmann et al. 2004). We show in Fig. 6 that the bimodal character of galaxy populations is also clearly present in some parameters derived from our spectral synthesis (Mateus et al., in prep.). These results suggest that galaxies have evolved through two major paths.
#### The down-sizing of galaxy populations
Another interesting result emerging from our synthesis of SDSS galaxies is the confirmation that there is a correlation between mass and galaxy age, in the sense that most of the stars in massive galaxies were formed long time ago, whereas galaxies with a large fraction of young or intermediate-age stars tend to be less massive. This phenomenon, known as the down-sizing of galaxy populations (Cowie et al. 1996; Juneau et al. 2004; Kodama et al. 2004), is clearly seem in our results with this new sample and, actually, is an essential piece of information on galaxy formation and evolution.
#### The problem of the alpha-enhancement
A major problem of spectral synthesis with a spectral base using BC03 spectra is that while their evolutionary tracks are scaled from the solar values, non-solar abundance patterns may occur in galaxies, mainly in early-types (e.g., Trager et al. 2000a,b). In practice, as discussed by BC03, lines associated to $`\alpha `$-elements are not fitted as well as other lines. This problem is indeed present in our results, as shown in Fig. 8 for a subsample of 2488 bona fide ellipticals with high S/N spectra. This figure indicates that the residuals in the absorption lines of some $`\alpha `$-elements correlate with the depth of the gravitational potential well of the galaxy (Cid Fernandes et al., in prep.), suggesting that our spectral synthesis results of early-type galaxies may be biased due to inadequacy of the spectral base adopted here to synthesize this type of galaxy.
## 6. Summary
We have developed and tested a method to fit galaxy spectra with a combination of spectra of individual simple stellar populations generated with state-of-the art evolutionary synthesis models. The main goal of this investigation was to examine the reliability of physical properties derived in this way. This goal was pursued by three different means: simulations, comparison with independent studies, and analysis of empirical results.
Our simulations show that the individual SSP strengths, encoded in the population vector $`\stackrel{}{x}`$, are subjected to large uncertainties, but robust results can be obtained by compressing $`\stackrel{}{x}`$ into coarser but useful indices. In particular, physically motivated indices such as mean stellar ages and metallicities are found to be well recovered by spectral synthesis even for relatively noisy spectra. Stellar masses, velocity dispersion and extinction are also found to be accurately retrieved.
We have applied our STARLIGHT code to a volume limited sample of over 50000 galaxies from the SDSS Data Release 2. The spectral fits are generally very good, and allow accurate measurements of emission lines from the starlight subtracted spectrum. We have compared our results to those obtained by the MPA/JHU group (K03; Brinchmann et al. 2004) with a different method to characterize the stellar populations of SDSS galaxies. The stellar extinctions and masses derived in these two studies are very strongly correlated. Furthermore, differences in the values of $`A_V`$ and $`M_{}`$ are found to be mostly due to the differences in the model ingredients (extinction law). Our estimates of stellar velocity dispersions and emission line properties are also in good agreement with those of the MPA/JHU group.
The confidence in the method is further strengthened by several empirical correlations between synthesis results and independent quantities. We find strong correlations between stellar and nebular metallicites, stellar and nebular extinctions, mean stellar age and the equivalent width of H$`\alpha `$, mean stellar age and the 4000 ร
break, stellar mass and velocity dispersion. These are all astrophysically reasonable results, which reinforce the conclusion that spectral synthesis is capable of producing reliable estimates of physical properties of galaxies. These results validate spectral synthesis as a powerful tool to study the history of galaxies.
We have also presented some preliminary results of an analysis of a magnitude-limited sample containing 20000 galaxies from SDSS which clearly reveal that the bimodality of galaxy populations is present in the parameters obtained from the synthesis. Our results are also consistent with the โdown-sizingโ scenario of galaxy formation and evolution. Finally, we point out one of the major problems of the spectral synthesis of early-type systems: the spectral base adopted here is based on solar-scaled evolutionary tracks which may not be appropriate for this type of galaxy.
## Acknowledgments
We congratulate D. Valls-Gabaud and M. Chaves for the organization of this workshop. Partial support from CNPq, FAPESP and the France-Brazil PICS program are also acknowledged. |
warning/0506/quant-ph0506086.html | ar5iv | text | # Holonomic quantum computation in decoherence-free subspaces
## Abstract
We show how to realize, by means of non-abelian quantum holonomies, a set of universal quantum gates acting on decoherence-free subspaces and subsystems. In this manner we bring together the quantum coherence stabilization virtues of decoherence-free subspaces and the fault-tolerance of all-geometric holonomic control. We discuss the implementation of this scheme in the context of quantum information processing using trapped ions and quantum dots.
Introduction.โ The implementation of quantum information processing (QIP) poses an unprecedented challenge to our capabilities of controlling the dynamics of quantum systems. The challenge is twofold and somewhat contradictory. On the one hand one must (i) maintain as much as possible the isolation of the computing degrees of freedom from the environment, in order to preserve their โquantumnessโ; on the other hand (ii) their dynamical evolution must be enacted with extreme precision in order to avoid errors whose propagation would quickly spoil the whole quantum computational process. To cope with the decoherence problem (i), active strategies such as quantum error correcting codes Steane:99 , as well as passive ones such as error avoiding codes Duan:97PRLZanardi:97cLidar:PRL98 , have been contrived. The latter are based on the symmetry structure of the system-environment interaction, which under certain circumstances allows for the existence of decoherence-free subspaces (DFS), i.e., subspaces of the system Hilbert state-space over which the dynamics is still unitary. DFSs have been experimentally demonstrated in a host of physical systems (e.g., Kielpinski:01 ; Kwiat:00Mohseni:02Ollerenshaw:02Bourennanne04a ). The DFS idea of symmetry-aided protection has been generalized to noiseless subsystems Knill:99aZanardi:99dKempe:00 , experimentally tested in Ref. Viola:01b .
Holonomic quantum computation (HQC) ZanardiRasetti:99 is an all-geometric strategy wherein QIP is realized by means of adiabatic non-abelian quantum holonomies Wilczek:84 . Quantum information is encoded in a degenerate eigenstate of a Hamiltonian depending on a set of controllable parameters e.g., external laser fields. When the latter are adiabatically driven along a suitable closed path, the initial quantum state is transformed by a non-trivial unitary transformation (holonomy) that is geometrical in nature. Following the ion traps HQC-implementation proposal Duan-Science:01 , several other schemes, based on a variety of physical setups, were proposed Recati:02Faoro:03Li:04Bernevig:05 . One expects the geometrical nature of quantum holonomies to endow HQC with inherent robustness against certain errors. This alleged fault-tolerance has only recently been subjected to serious scrutiny Solinas:04Fuentes-Guridi:05 ; the resulting, still developing picture, is that while stability against decoherence must be further assessed (indeed the adiabatic theorem was only recently generalized to open quantum systems SarandyLidar:04 ), HQC seems to exhibit a strong robustness against stochastic errors in the control process generating the required adiabatic loops Zhu:04 . From this point of view HQC seems to be promising with respect to the general challenge (ii).
In this work we describe a QIP scheme which combines DFSs and HQC HQC-DFS-note ; carollo . More specifically, we show how to perform universal quantum computation within a two-qubit DFS for collective dephasing by using non-abelian holonomies only. The discussion is then extended to consider general collective decoherence as well. The appeal of such a strategy should, in view of the above, be evident: try to bring together the best of two worlds, namely the resilience of the DFS approach against environment-induced decoherence and the operational robustness of HQC. Moreover, we formulate our results using rather generic Hamiltonians, so that the scheme proposed in this work appears to be a suitable candidate for experimental demonstration in a variety of systems, including trapped ions and quantum dots.
Dark-states in a Decoherence-Free subspace.โ Let us start by considering a four-qubit system with state-space $`_4(๐^2)^{\mathrm{\hspace{0.17em}4}}`$. We denote by $`X_l,Y_l`$ and $`Z_l`$ the three Pauli matrices acting on the $`l`$th qubit; for any pair $`lm`$ of qubits we define the operators $`R_{lm}^x:=\frac{1}{2}\left(X_lX_m+Y_lY_m\right)`$, $`R_{lm}^y:=\frac{1}{2}\left(X_lY_mY_lX_m\right)`$, $`R_{lm}^z:=\frac{1}{2}\left(Z_mZ_l\right)`$. These operators have a non-trivial action over the subspace of $`_{lm}๐_l^2๐_m^2`$ spanned by $`|0_l1_m`$ and $`|1_l0_m`$ (where $`|0/|1`$ are the $`+/`$ eigenstates of $`Z`$, respectively): they provide a faithful representation of the su(2) algebra over this subspace, where they act as the Pauli matrices, while they vanish on the orthogonal complement spanned by $`|0_l0_m`$ and $`|1_l1_m`$. Let $`Z:=_{i=1}^4Z_i`$, then $`[R_{lm}^\alpha ,Z]=0`$ $`(l,m)`$. It follows that every eigenspace of $`Z`$ is invariant under the action of the $`R_{lm}^\alpha `$โs. In particular this holds true for the subspace
$$๐:=\mathrm{span}\{|1000,|0100,|0010,|0001\}.$$
(1)
$`๐`$ is a DFS against collective dephasing Duan:97PRLZanardi:97cLidar:PRL98 , i.e., states in $`๐`$ are immune from decoherence induced by system-bath interactions of the form $`ZB`$, where $`B`$ is an arbitrary bath operator. Collective dephasing is known to be a major source of decoherence in ion-trap based QIP Kielpinski:01 .
We assume that the system dynamics is generated by
$$H=\underset{l>m}{}(J_{lm}^xR_{lm}^x+J_{lm}^yR_{lm}^y),$$
(2)
where $`l,m`$ are qubit indices and the $`J_{lm}`$โs are *controllable* coupling constants. These are the parameters that will be driven along controlled adiabatic loops to enact quantum gates via non-abelian holonomies. When $`J_{lm}^y=0`$, $`H`$ is the XY Hamiltonian found in a variety of quantum computing proposals, e.g., the quantum Hall proposal Mozyrsky:01 , quantum dots Imamoglu:99 and atoms in cavities Zheng:00 . It also describes trapped ions subject to the Sรธrensen-Mรธlmer scheme Sorensen:00 . The case $`J_{lm}^x=0`$ is related to the XY model via a unitary transformation.
By way of introduction to our HQC-DFS scheme, note that Hamiltonian (2) has a hidden multi-level structure with interesting properties (isomorphic to the one exploited in Ref. Duan-Science:01 ). Indeed, the Hamiltonian $`H_{lmn}=_{j=m,n}J_{jl}(R_{jl}^x\mathrm{cos}\phi _{jl}+R_{jl}^y\mathrm{sin}\phi _{jl})`$, in the basis $`\{|e:=|100_{lmn},|g_1:=|010_{lmn},|g_2:=|001_{lmn}\}`$, with $`lmn`$, takes the form
$$H_{lmn}=J_{lm}e^{i\phi _{lm}}|eg_1|+J_{ln}e^{i\phi _{ln}}|eg_2|+\mathrm{h}.\mathrm{c}.$$
(3)
This is a so-called Lambda scheme, with $`|e`$ at the top and $`|g_1,|g_2`$ at the bottom. It is well known that for every value of the $`J_{jl}`$โs and $`\phi _{jl}`$โs, Hamiltonian (3) has one *dark state* $`|D(J_{lm},J_{ln})J_{lm}e^{i\phi _{lm}}|g_2J_{ln}e^{i\phi _{ln}}|g_1`$, i.e., a state satisfying $`H_{lmn}|D(J_{lm},J_{ln})=0`$.
The key idea is now as follows: moving in the control parameter space to the point $`J_{lm}/J_{ln}=0`$, one has that the dark state is given by $`|g_1`$; then, if one adiabatically changes the parameters along a closed loop, at the end the state $`|g_1`$ will pick up a non-trivial geometric phase. This simple fact, suitably generalized, is the basic ingredient of our HQC scheme. We further supplement the dark state $`|g_1`$ by another state, that is also annihilated by the system Hamiltonian (and thus does not acquire a dynamical phase either); these states together will form our qubit. In other words, the crucial observation is that one can embed within the DFS $`๐`$ \[Eq. (1)\] a manifold of dark states that, for specific values of the coupling constants in Eq. (2), coincide with logical encoded qubits. Holonomic manipulations are then used to generate a universal set of quantum logic gates.
Let us stress that even if during state manipulation the system described by Eq. (2) leaks out of the logical encoding subspace (the states $`|0_L`$ and $`|1_L`$ defined below), e.g., due to the breakdown of adiabaticity SarandyLidar:04 , this results just in the kind of errors that HQC is robust against Solinas:04Fuentes-Guridi:05 . Moreover, during such leakage the system *never* abandons the DFS $`๐`$. Thus protection against collective dephasing is maintained throughout the whole gating period, which in turns allows one to stretch the gating time in such a way as to fulfill the adiabatic constraint for a longer period of time than would be possible without the DFS. These remarks at least partially counter a standard objection to HQC, that the use of slow gates gives decoherence more time to exert its detrimental effects.
One-qubit holonomic gates.โ We now show how to realize the single-qubit phase gate $`\mathrm{exp}(i\phi Z_L)`$, where $`Z_L|\alpha =(1)^\alpha |\alpha _L`$ (note that $`Z_L=R_{12}^z`$). We encode a logical qubit in $`๐`$ by $`|1_L:=|0010`$ and $`|0_L:=|0001`$, while the remaining pair $`|a_1:=|1000`$ and $`|a_2:=|0100`$ plays the role of ancillae. Let us set all the $`J_{lm}`$ to zero, except $`J_{24}`$ and $`J_{34}`$, such that the Hamiltonian (2) reduces to
$$H_Z=J_{24}(R_{24}^x\mathrm{cos}\phi _{24}+R_{24}^y\mathrm{sin}\phi _{24})+J_{34}R_{34}^x.$$
(4)
We can also write $`H_Z=J_{34}|a_1a_2|+J_{24}e^{i\phi }|a_11|_L+\mathrm{h}.\mathrm{c}`$. This is a Lambda configuration with $`|a_1`$ at the top and $`|a_2,|1_L`$ at the bottom. Therefore, as in our discussion above, $`H_Z`$ has a zero-eigenvalue eigenstate (dark state) given by $`|\mathrm{\Psi }_1=\mathrm{cos}\theta |1_L\mathrm{sin}\theta e^{i\phi }|a_2`$ where $`\theta =\mathrm{tan}^1(J_{24}/J_{34})`$ and $`\phi =\phi _{24}`$. The state $`|0_L`$ is also a zero-eigenvalue eigenstate that does not depend on the parameters $`\theta `$ and $`\phi `$. By adiabatically changing $`\phi `$ in such a way as to have a loop starting from $`\theta =\phi =0`$, the state $`|1_L`$ acquires a Berry phase which is proportional to the solid angle $`\mathrm{\Omega }_Z`$ swept out by the vector $`(\varphi ,\theta )`$ Berry:84 . Therefore after this adiabatic loop one has that: $`|0_L|0_L`$ and $`|1_Le^{i\mathrm{\Omega }_Z/2}|1_L`$, which is clearly equivalent to the operation $`\mathrm{exp}(i\mathrm{\Omega }_ZZ_L/2)`$.
In order to obtain a universal set of gates we need to generate at least two non-commuting single-qubit gates. Therefore we next show how to implement $`\mathrm{exp}(i\phi X_L)=\mathrm{exp}(i\phi R_{12}^x)`$. One way is to once again establish an isomorphism between the system governed by Eq. (2), restricted to $`๐`$, and the HQC model of Ref. Duan-Science:01 . Here we provide an independent derivation. We turn on the couplings in such a way as to obtain the Hamiltonian
$$H_X=J_{34}R_{34}^x+J_{24}(\mathrm{cos}\phi \frac{R_{24}^xR_{14}^x}{\sqrt{2}}+\mathrm{sin}\phi \frac{R_{24}^yR_{14}^y}{\sqrt{2}}).$$
(5)
Let $`|\pm _L:=(|1_L\pm |0_L)/\sqrt{2}`$. Then one can readily check that under the action of $`(R_{24}^xR_{14}^x)/\sqrt{2}`$ the state $`|+_L0`$,$`|_L|a_1`$, and $`|a_1|_L`$, etc. Therefore $`H_X=J_{34}|a_1a_2|+J_{24}e^{i\phi }|a_1|_L+\mathrm{h}`$.$`\mathrm{c}`$. This is a Lambda configuration with $`|a_1`$ at the top and $`|a_2,|_L`$ at the bottom, so that $`H_X`$ supports a dark state $`|\mathrm{\Psi }_2=\mathrm{cos}\theta |_L\mathrm{sin}\theta e^{i\phi }|a_2`$ where $`\theta =\mathrm{tan}^1(J_{24}/J_{34})`$. The similarity between $`H_X`$ and $`H_Z`$ is evident. Then, by executing an adiabatic loop in the parameter space in analogy to the $`H_Z`$ case, one obtains the geometric evolution $`|+_L|+_L`$, $`|_Le^{i\mathrm{\Omega }_X/2}|_L`$, where now $`\mathrm{\Omega }_X`$ is the solid angle swept out by the vector $`(\theta ,\phi )`$. Switching back to the computational basis this transformation amounts to the map $`|0_L\mathrm{cos}\frac{\mathrm{\Omega }_X}{4}|0_L+i\mathrm{sin}\frac{\mathrm{\Omega }_X}{4}|1_L`$, and $`|1_L\mathrm{cos}\frac{\mathrm{\Omega }_X}{4}|1_L+i\mathrm{sin}\frac{\mathrm{\Omega }_X}{4}|0_L`$, which is equivalent to $`\mathrm{exp}(i\mathrm{\Omega }_XR_{12}^x/4)`$.
Two-Qubit holonomic gates.โ A crucial, and typically rather demanding part of any QIP implementation proposal, is the realization of an entangling two-qubit gate. We next consider this problem and demonstrate how to solve it following a strategy relying on the same abstract holonomic structure as that discussed above for one-qubit gates. The total state-space is now $`_4^{\mathrm{\hspace{0.17em}2}}`$, while the two-qubit code is spanned by $`|\alpha _L|\beta _L`$,$`(\alpha ,\beta =0,1)`$. The states $`|a_i^{\mathrm{\hspace{0.17em}2}}`$($`i=1,2`$) will function as ancillae. Importantly, once again all the relevant states belong to a DFS against collective dephasing.
Let us suppose that one can engineer the following controllable four-qubit interaction
$`H_4`$ $`=`$ $`J_{24,68}(R_{24}^x\mathrm{cos}\phi +R_{24}^y\mathrm{sin}\phi )(R_{68}^x\mathrm{cos}\phi +R_{68}^y\mathrm{sin}\phi )`$ (6)
$`+`$ $`J_{34,78}R_{34}^xR_{78}^x,`$
which should be recognizable as a straightforward extension of the one-qubit Hamiltonian (4). Below we discuss the implementation of such a Hamiltonian. To explicitly exhibit the dark state structure, we write this Hamiltonian in the form $`H_4=J_{34,78}R_{34}^x(|a_1a_2|)^{\mathrm{\hspace{0.17em}2}}+J_{24,68}e^{i\phi }(|a_11_L|)^{\mathrm{\hspace{0.17em}2}}+\mathrm{h}.\mathrm{c}.`$, from which it is easily seen that there is one dark state, given by $`\mathrm{cos}\theta |1_L^{\mathrm{\hspace{0.17em}2}}\mathrm{sin}\theta e^{i\phi }|a_2^{\mathrm{\hspace{0.17em}2}}`$, where $`\theta =\mathrm{tan}^1(J_{24,68}/J_{34,78})`$. After the same kind of adiabatic cyclic evolution as described in the single-qubit gate case, one obtains in $`๐^{\mathrm{\hspace{0.17em}2}}`$ the controlled phase-shift gate $`CP=\mathrm{diag}(1,1,1,e^{i\mathrm{\Omega }_P/2})`$, where, as usual, $`\mathrm{\Omega }_P`$ is the solid angle swept out by the vector $`(\theta ,\phi )`$.
Extensions and generalizations.โ The scheme described so far for the case of collective dephasing is straightforwardly scalable to an arbitrary number $`N`$ of encoded DF qubits. The total space is now given by $`๐^N(๐^4)^N`$, with $`๐`$ given in Eq. (1). This, of course, is still an eigenspace of the collective spin $`z`$-component, i.e., $`Z=_{l=1}^{4N}Z_l`$. Following the procedure established above, the controllable Hamiltonian used to generate a controlled phase-shift between the $`i`$-th and the $`j`$-th encoded qubits, has the same structure as $`H_4`$ \[Eq. (6)\], where, with obvious notation, $`J_{24;68}J_{4(i1)+2,4(i1)+4;4(j1)+2,4(j1)+4}`$ and $`J_{34;78}J_{4(i1)+3,4(i1)+4;4(j1)+3,4(j1)+4}`$, and similarly for the $`R_{kl}^x,R_{kl}^y`$ operators.
Next, by making use of *noiseless subsystems* (NS) Knill:99aZanardi:99dKempe:00 , we show that our combined HQC-DFS strategy can also be applied against *general* collective decoherence Duan:97PRLZanardi:97cLidar:PRL98 . Our arguments will be existential in nature, with constructive details to be discussed elsewhere. Let $`_J๐^{2J+1}`$ denote the total spin-$`J`$ irreducible representation of $`su(2)`$. The state-space of five qubits, i.e., $`_{1/2}^{\mathrm{\hspace{0.17em}5}},`$ decomposes, with respect to the collective $`su(2)`$-representation (Clebsch-Gordan decomposition) as follows
$$_{1/2}^{\mathrm{\hspace{0.17em}5}}๐^4_{3/2}๐^5_{1/2}๐_{5/2}.$$
(7)
Each of the $`๐^{n_J}`$ factors represents the multiplicity of the total spin-$`J`$ irreducible representation *and corresponds to a NS against collective general decoherence* Knill:99aZanardi:99dKempe:00 . Consider the first term in (7): the multiplicity factor $`๐^4`$ for the $`J=3/2`$ representation provides a four-dimensional NS. It might then encode two noiseless qubits, but, since we wish to perform QIP with holonomies, we will instead use this $`๐^4`$ space as a code for just one noiseless qubit $`|\stackrel{~}{\alpha }_L`$, ($`\alpha =0,1`$) and two ancillary states $`|\stackrel{~}{a}_i`$, ($`i=1,2`$). Suppose now that one is able to enact the controllable Hamiltonian $`H_{\mathrm{NS}}=J^{}|\stackrel{~}{a}_1\stackrel{~}{a}_2|+J_0^{\prime \prime }e^{i\phi }|\stackrel{~}{a}_1\stackrel{~}{0}|_L+J_1^{\prime \prime }e^{i\phi }|\stackrel{~}{a}_1\stackrel{~}{1}|_L+\mathrm{h}.\mathrm{c}.`$, which, when $`J_0^{\prime \prime }=0`$, admits a dark state given by $`|\mathrm{\Psi }_3=\mathrm{cos}\theta |\stackrel{~}{1}_L\mathrm{sin}\theta e^{i\phi }|\stackrel{~}{a}_2`$, where $`\theta =\mathrm{tan}^1(J_1^{\prime \prime }/J^{})`$. By resorting to the same considerations as previously developed for the collective dephasing case, it should be clear that this allows us to enact a phase gate $`Z_L`$ between $`|\stackrel{~}{0}_L`$ and $`|\stackrel{~}{1}_L`$. By choosing $`J_0^{\prime \prime },J_1^{\prime \prime }`$ such that $`H_{\mathrm{NS}}`$ becomes $`J^{}|\stackrel{~}{a}_1\stackrel{~}{a}_2|+J_{}^{\prime \prime }e^{i\phi }|\stackrel{~}{a}_1\stackrel{~}{}|_L+\mathrm{h}.\mathrm{c}.`$ we can enact the $`X_L`$ gate, so that universal single-qubit control by holonomies can be achieved in this case as well.
To realize the required multi-level controllable Hamiltonian, we observe that the $`๐^4`$ space under consideration is a four-dimensional irreducible representation of the permutation group $`๐ฎ_5`$ \[acting over the whole space as $`\widehat{\sigma }_{j=1}^5|\stackrel{~}{\alpha }_j=_{j=1}^5|\stackrel{~}{\alpha }_{\sigma (j)}`$, where$`\sigma ๐ฎ_5`$\]. Following Ref. Zanardi:04 , universal control over this irreducible representation space amounts to the ability to switch on and off a pair of *generic* Hamiltonians in the group algebra of $`๐ฎ_5`$. An important example is provided by Heisenberg *exchange* Hamiltonians, i.e., $`_{l<m}J_{lm}๐_l๐_m`$ \[where $`๐_l=(X_l,Y_l,Z_l)/2`$\], which are naturally available interactions in several spin-based proposals for QIP Burkard:99 ; Kane:98Vrijen:00 . The construction of two-encoded-qubit holonomic gates, and the generalization to arbitrary numbers of such qubits, again follow the same pattern as in the collective dephasing case.
Implementation.โ We note that all required Hamiltonians, $`H_Z`$, $`H_X`$ and $`H_4`$, have a similar form. These Hamiltonians all involve control over both $`\theta `$ and $`\phi `$. However, we are free to choose any loop $`C(\theta ,\phi )`$ in the $`(\theta ,\phi )`$-parameter space, and we can choose a loop which toggles between $`\theta `$ and $`\phi `$. For example, the loop $`C(\theta ,\phi ):`$ $`(0,0)(\frac{\pi }{2},0)(\frac{\pi }{2},\phi _0)(0,\phi _0)(0,0)`$ has this property, where $`\phi _0`$ is the solid angle $`\mathrm{\Omega }`$ swept out by the vector $`(\theta ,\phi )`$ for this specific loop.
Let us briefly address the feasibility of the control of $`H_Z`$, $`H_X`$ and $`H_4`$ in the context of actual quantum computing proposals. In proposals based on electron spins in quantum dots Burkard:99 , the $`H_Z`$ and $`H_X`$ Hamiltonians are available when one takes into account the spin-orbit interaction, for then the effective spin-spin interaction becomes $`H_{ij}(t)=J_{ij}(t)[\beta (t)(X_iY_jY_iX_j)+(1+\gamma (t))(X_iX_j+Y_iY_j)+Z_iZ_j]`$ Kavokin:01 ; Stepanenko:04 . This Hamiltonian has enough degrees of freedom to implement $`H_Z`$ and $`H_X`$, via separate control of $`J`$ and the dimensionless anisotropy parameters $`\beta `$ and $`\gamma `$, assuming the latter can be made $`1`$, e.g., when the inter-dot distance is large Kavokin:01 . The parameters $`\beta `$ and $`\gamma `$ can be controlled via the confining potential or via pulse shaping, as has been discussed in detail in Ref. Stepanenko:04 , and $`\phi =\mathrm{tan}^1\beta /(1+\gamma )`$. The parameter $`\theta `$ is controllable via, e.g., $`J_{24}`$ and $`J_{34}`$. Implementation of $`H_4`$ is undoubtedly more challenging, since four-body interactions are involved. Recent results indicate that such terms arise by simultaneously coupling four quantum dots MizelLidar:04 ; their relative strength can be controlled, e.g., by adjusting the confining potentials WoodworthMizelLidar:05 .
In proposals based on trapped ions, the implementation of $`H_Z`$ and $`H_X`$ is directly possible using the Sรธrensen-Mรธlmer (SM) scheme Sorensen:00 . The SM gate between two ions implements $`H_Z`$ and $`H_X`$ with control over the various terms achieved via the phases of two lasers Kielpinski:02LidarWu:02 . It is also possible to implement $`H_4`$ using the SM scheme, via control over two pairs of ions Kielpinski:02LidarWu:02 . Given that a geometric two ion-qubit phase gate has already been demonstrated Leibfried:03 , trapped ions seem to be particularly favorable for the implementation of our proposed HQC-DFS scheme.
Conclusions.โ We have combined the DFS and HQC techniques, and showed how to realize universal quantum computation over a scalable DFS against collective dephasing by using adiabatic holonomies only. We discussed an extension to the general collective decoherence case, arguing that controllability of exchange Hamiltonians would suffice. Remarkably, the whole computational process is carried out *within* the DFS. The DFS embedding, along with the all-geometrical nature of HQC, promises to give to this scheme a twofold resilience, against decoherence and stochastic control errors. The proposed universal quantum gates are carried out via adiabatic manipulation of Hamiltonians that are controllable in several proposals for QIP implementations. We are therefore hopeful that the theoretical ideas presented here may stimulate corresponding experimental activity.
Acknowledgements.โ P.Z. acknowledges financial support from the European Union FET project TOPQIP (Contract No. IST-2001-39215) and discussions with A. Carollo. D.A.L. acknowledges financial support from the DARPA-QuIST program and the Sloan Foundation. |
warning/0506/math0506565.html | ar5iv | text | # Local Lie algebra determines base manifold 11footnote 1Research supported by the Polish Ministry of Scientific Research and Information Technology under grant No 2 P03A 020 24.
## 1 Introduction
The classical result of Shanks and Pursell \[PS\] states that the Lie algebra $`๐ณ_c(M)`$ of all compactly supported smooth vector fields on a smooth manifold $`M`$ determines the manifold $`M`$, i.e., the Lie algebras $`๐ณ_c(M_1)`$ and $`๐ณ_c(M_2)`$ are isomorphic if and only if $`M_1`$ and $`M_2`$ are diffeomorphic. A similar theorem holds for other complete and transitive Lie algebras of vector fields \[KMO1, KMO2\] and for the Lie algebras of all differential and pseudodifferential operators \[DS, GP\].
There is a huge list of papers in which special geometric situations (hamiltonian, contact, group invariant, foliation preserving, etc., vector fields) are concerned. Let us mention the results of Omori \[O1\] (Ch. X) and \[O2\](\[Ch. XII), or \[Ab, AG, FT, HM, Ry, G5\], for which specific tools were developed in each case. There is however a case when the answer is more or less complete in the whole generality. These are the Lie algebras of vector fields which are modules over the corresponding rings of functions (we shall call them modular). The standard model of a modular Lie algebra of vector fields is the Lie algebra $`๐ณ()`$ of all vector fields tangent to a given (generalized) foliation $``$. If Pursell-Shanks-type results are concerned in this context, let us recall the work of Amemiya \[Am\] and our paper \[G1\], where the developed algebraic approach made it possible to consider analytic cases as well. The method of Shanks and Pursell consists of the description of maximal ideals in the Lie algebra $`๐ณ_c(M)`$ in terms of the points of $`M`$: maximal ideals are of the form $`\stackrel{~}{p}`$ for $`pM,`$ where $`\stackrel{~}{p}`$ consists of vector fields which are flat at $`p.`$ This method, however, fails in analytic cases, since analytic vector fields flat at $`p`$ are zero on the corresponding component of $`M.`$ Therefore in \[Am, G1\] maximal finite-codimensional subalgebras are used instead of ideals. A similar approach is used in \[GG\] for proving that the Lie algebras associated with Lie algebroids determine base manifolds.
The whole story for modular Lie algebras of vector fields has been in a sense finished by the brilliant purely algebraic result of Skryabin \[S1\], where one associates the associative algebra of functions with the Lie algebra of vector fields without any description of the points of the manifold as ideals. This final result implies in particular that, in the case when modular Lie algebras of vector fields contain finite families of vector fields with no common zeros (we say that they are strongly nonโsingular), isomorphisms between them are generated by isomorphisms of corresponding algebras of functions, i.e., by diffeomorphisms of underlying manifolds.
On the other hand, there are many geometrically interesting Lie algebras of vector fields which are not modular, e.g. the Lie algebras of hamiltonian vector fields on a Poisson manifold etc. For such algebras the situation is much more complicated and no analog of Skryabin method is known in these cases. In \[G6\] a Pursell-Shanks-type result for the Lie algebras associated with Jacobi structures on a manifold has been announced. The result suggests that the concept of a Jacobi structure should be developed for sections of an arbitrary line bundle rather than for the algebra of functions, i.e., sections of the trivial line bundle. This is exactly the concept of local Lie algebra in the sense of A. A. Kirillov \[Ki\] which we will call also Jacobi-Kirillov bundle.
In the present note we complete the Lie algebroid result of \[GG\] by proving that the local Lie algebra determines the base manifold up to a diffeomorphism if only the anchor map is nowhere-vanishing (Theorem 7). The methods, however, are more complicated (due to the fact that the Lie algebra of Jacobi-hamiltonian vector fields is not modular) and different from those in \[GG\]. A part of these methods is a modification of what has been sketched in \[G6\]. However, the full generalization of \[G6\] for local Lie algebras on arbitrary line bundles, i.e., the description of isomorphisms of local Lie algebras, is much more delicate and we postpone it to a separate paper. Note also that in our approach we admit different categories of differentiability: smooth, real-analytic, and holomorphic (Stein manifolds).
## 2 Jacobi modules
What we will call Jacobi module is an algebraic counterpart of geometric structures which include Lie algebroids and Jacobi structures (or, more generally, local Lie algebras in the sense of Kirillov \[Ki\]). For a short survey one can see \[G7\], where these geometric structures appeared under the name of Lie QD-algebroids.
The concept of a Lie algebroid (or its pure algebraic counterpart โ a Lie pseudoalgebra) is one of the most natural concepts in geometry.
Definition 1. Let $`R`$ be a commutative and unitary ring, and let $`๐`$ be a commutative $`R`$-algebra. A Lie pseudoalgebra over $`R`$ and $`๐`$ is an $`๐`$-module $``$ together with a bracket $`[,]:\times `$ on the module $``$, and an $`๐`$-module morphism $`\alpha :\text{Der}(๐)`$ from $``$ to the $`๐`$-module $`\text{Der}(๐)`$ of derivations of $`๐`$, called the anchor of $``$, such that
(i) the bracket on $``$ is $`R`$-bilinear, alternating, and satisfies the Jacobi identity:
$$[[X,Y],Z]=[X,[Y,Z]][Y,[X,Z]].$$
(ii) For all $`X,Y`$ and all $`f๐`$ we have
$$[X,fY]=f[X,Y]+\alpha (X)(f)Y;$$
(1)
(iii) $`\alpha ([X,Y])=[\alpha (X),\alpha (Y)]_c`$ for all $`X,Y`$, where $`[,]_c`$ is the commutator bracket on $`\mathrm{Der}(๐)`$.
A Lie algebroid on a vector bundle $`E`$ over a base manifold $`M`$ is a Lie pseudoalgebra on the $`(,C^{\mathrm{}}(M))`$-module $`=Sec(E)`$ of smooth sections of $`E`$. Here the anchor map is described by a vector bundle morphism $`\alpha :ETM`$ which induces the bracket homomorphism from $`(,[,])`$ into the Lie algebra $`(๐ณ(M),[,]_{vf})`$ of vector fields on $`M`$. In this case, as in the case of any faithful $`๐`$-module $``$, i.e., when $`fX=0`$ for all $`X`$ implies $`f=0`$, the axiom (iii) is a consequence of (i) and (ii). Of course, we can consider Lie algebroids in the real-analytic or holomorphic (on complex holomorphic bundles over Stein manifolds) category as well.
Lie pseudoalgebras appeared first in a paper by Herz \[He\] but one can find similar concepts under more than a dozen of names in the literature (e.g. Lie modules, $`(R,A)`$-Lie algebras, Lie-Cartan pairs, Lie-Rinehart algebras, differential algebras, etc.). Lie algebroids were introduced by Pradines \[Pr\] as infinitesimal parts of differentiable groupoids. In the same year a book by Nelson \[Ne\] was published, where a general theory of Lie modules together with a big part of the corresponding differential calculus can be found. We also refer to a survey article by Mackenzie \[Ma2\].
Note that Lie algebroids on a singleton base space are Lie algebras. Another canonical example is the tangent bundle $`TM`$ with the canonical bracket $`[,]_{vf}`$ on the space $`๐ณ(M)=Sec(TM)`$ of vector fields.
The property (1) of the bracket in the $`๐`$-module $``$ can be expressed as the fact that $`ad_X=[X,]`$ is a quasi-derivation in $``$, i.e., an $`R`$-linear operator $`D`$ in $``$ such that $`D(fY)=fD(Y)+\widehat{D}(f)Y`$ for any $`f๐`$ and certain derivation $`\widehat{D}`$ of $`๐`$ called the anchor of $`D`$. The concept of quasi-derivation can be traced back to N. Jacobson \[J1, J2\] as a special case of his pseudo-linear endomorphism. It has appeared also in \[Ne\] under the name of a module derivation and used to define linear connections in the algebraic setting. In the geometric setting, for Lie algebroids, it has been studied in \[Ma1\], Ch. III, under the name covariant differential operator.
Starting with the notion of Lie pseudoalgebra we obtain the notion of Jacobi module when we drop the assumption that the anchor map is $`๐`$-linear.
Definition 2. Let $`R`$ be a commutative and unitary ring, and let $`๐`$ be a commutative $`R`$-algebra. A Jacobi module over $`(R,๐)`$ is an $`๐`$-module $``$ together with a bracket $`[,]:\times `$ on the module $``$, and an $`R`$-module morphism $`\alpha :\text{Der}(๐)`$ from $``$ to the $`๐`$-module $`\text{Der}(๐)`$ of derivations of $`๐`$, called the anchor of $``$, such that (i)-(iii) of Definition 1 are satisfied. Again, for faithful $``$, the axiom (iii) follows from (i) and (ii). This concept is in a sense already present in \[He\], although in \[He\] it has been assumed that $`๐`$ is a field. It has been observed in \[He\] that every Jacobi module (over a field) of dimension $`>1`$ is just a Lie pseudoalgebra.
Definition 3. (cf. \[G7\]) A Lie QD-algebroid is a Jacobi module structure on the $`(,C^{\mathrm{}}(M)`$-module $`=\text{Sec}(E)`$ of sections of a vector bundle $`E`$ over a manifold $`M`$.
The case $`\text{rank}(E)=1`$ is special by many reasons and it was originally studied by A. A. Kirillov \[Ki\]. For a trivial bundle, well-known examples are those given by Poisson or, more generally, Jacobi brackets (cf. \[Li\]). In \[Ki\] such structures on line bundles are called local Lie algebras and in \[Mr\]Jacobi bundles. We will refer to them also as to local Lie algebras or Jacobi-Kirillov bundles and to the corresponding brackets as to Jacobi-Kirillov brackets.
Definition 4. A Jacobi-Kirillov bundle (local Lie algebra in the sense of Kirillov) is a Lie QD-algebroid on a vector bundle of rank 1. In other words, a Jacobi-Kirillov bundle is a Jacobi module structure on the $`(,C^{\mathrm{}}(M))`$-module $``$ of sections of a line bundle $`E`$ over a smooth manifold $`M`$. The corresponding bracket on $``$ we call Jacobi-Kirillov bracket and the values of the anchor map $`\alpha :๐ณ(M)`$ we call Jacobi-hamiltonian vector fields.
It is easy to see (cf. \[G7\]) that any Lie QD-algebroid on a vector bundle of rank $`>1`$ must be a Lie algebroid. Of course, we can consider Lie QD-algebroids and Jacobi-Kirillov bundles in real-analytic or in holomorphic category as well.
Since quasi-derivations are particular first-order differential operators in the algebraic sense, it is easy to see that, for a Jacobi module $``$ over $`(R,๐)`$, the anchor map $`\alpha :\text{Der}(A)`$ is also a first-order differential operator, i.e.,
$$\alpha (fgX)=f\alpha (gX)+g\alpha (fX)fg\alpha (X)$$
(2)
for any $`f,g๐`$ and $`X`$. Denoting the Jacobi-hamiltonian vector field $`\alpha (X)`$ shortly by $`\widehat{X}`$, we can write for any $`f,g๐`$ and $`X,Y`$,
$`[gX,fY]`$ $`=`$ $`\widehat{gX}(f)Yf\widehat{Y}(g)X+fg[X,Y]`$
$`=`$ $`g\widehat{X}(f)Y\widehat{fY}(g)X+fg[X,Y],`$
so that for the map $`\mathrm{\Lambda }_X:๐\times ๐๐`$ defined by $`\mathrm{\Lambda }_X(g,f):=\widehat{gX}(f)g\widehat{X}(f)`$ we have
$$\mathrm{\Lambda }_X(g,f)Y=\mathrm{\Lambda }_Y(f,g)X.$$
(4)
The above identity implies clearly that, roughly speaking, $`\text{rank}_๐=1`$ โat points where $`\mathrm{\Lambda }`$ is non-vanishingโ (cf. \[G7\]) and that
$$\mathrm{\Lambda }_X(g,f)X=\mathrm{\Lambda }_X(f,g)X.$$
(5)
The identity (5) does not contain much information about $`\mathrm{\Lambda }_X`$ if there is โmuch torsionโ in the module $``$. But if, for example, there is a torsion-free element in $``$, say $`X_0`$, (this is the case of the module of sections of a vector bundle), then the situation is simpler. In view of (5), $`\mathrm{\Lambda }_{X_0}`$ is skew-symmetric and, in turn, by (4) every $`\mathrm{\Lambda }_X`$ is skew-symmetric. Every $`\mathrm{\Lambda }_X`$ is by definition a derivation with respect to the second argument, so, being skew-symmetric, it is a derivation also with respect to the first argument. Since in view of (2),
$$[gX,fX]=\left(g\widehat{X}(f)f\widehat{X}(g)+\mathrm{\Lambda }_X(g,f)\right)X,$$
and since $`\mathrm{\Lambda }_X`$ and $`\widehat{X}`$ respect the annihilator $`\text{Ann}(X)=\{f๐:fX=0\}`$, we get easily the following.
###### Proposition 1
If $``$ is a Jacobi module over $`(R,๐)`$, then, for every $`X`$, the map $`\mathrm{\Lambda }_X:๐\times ๐๐`$ induces a skew-symmetric bi-derivation of $`๐/\mathrm{Ann}(X)`$, the derivation $`\widehat{X}`$ of $`๐`$ induces a derivation of $`๐/\mathrm{Ann}(X)`$ and the bracket
$$\{\overline{f},\overline{g}\}_X=\mathrm{\Lambda }_X(\overline{f},\overline{g})+\overline{f}\widehat{X}(\overline{g})\overline{g}\widehat{X}(\overline{f}),$$
(6)
where $`\overline{f}`$ denotes the class of $`f๐`$ in $`๐/\mathrm{Ann}(X)`$, is a Jacobi bracket on $`๐/\mathrm{Ann}(X)`$ associated with the Jacobi structure $`(\mathrm{\Lambda }_X,\widehat{X})`$. Moreover, $`๐/\mathrm{Ann}(X)\overline{f}fX`$ is a Lie algebra homomorphism of the bracket $`\{,\}_X`$ into $`[,]`$.
For pure algebraic approaches to Jacobi brackets we refer to \[S2, S3, G4\].
###### Corollary 1
If the $`๐`$-module $``$ is generated by torsion-free elements, then for every $`X`$, the map $`\mathrm{\Lambda }_X:๐\times ๐๐`$ is a skew-symmetric bi-derivation and the bracket bracket
$$\{f,g\}_X=\mathrm{\Lambda }_X(f,g)+f\widehat{X}(g)g\widehat{X}(f),$$
(7)
is a Jacobi bracket on $`๐`$ associated with the Jacobi structure $`(\mathrm{\Lambda }_X,\widehat{X})`$. Moreover, $`๐ffX`$ is a Lie algebra homomorphism of the bracket $`\{,\}_X`$ into $`[,]`$.
For any torsion-free generated Jacobi module, e.g. a module of sections of a vector bundle, we have additional identities as shows the following.
###### Proposition 2
If the $`๐`$-module $``$ is generated by torsion-free elements, then for all $`f_1,\mathrm{},f_m๐`$, $`m2`$, and all $`X,Y`$
$$(a)(m1)[FX,Y]=\underset{i=1}{\overset{m}{}}[F_iX,f_iY][X,FY]$$
and
$$(b)(m2)[FX,Y]=\underset{i=1}{\overset{m1}{}}[F_iX,f_iY]+[F_mY,f_mX],$$
where $`F=_{i=1}^mf_i`$, $`F_k=_{ik}f_i`$.
Proof.- (a) We have (cf. (2))
$`{\displaystyle \underset{i=1}{\overset{m}{}}}[F_iX,f_iY]`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}\left(F_i\widehat{X}(f_i)Yf_i\widehat{Y}(F_i)X+F[X,Y]+\mathrm{\Lambda }_X(F_i,f_i)Y\right)`$
$`=`$ $`\widehat{X}(F)Y(m1)\widehat{Y}(F)X+mF[X,Y]+{\displaystyle \underset{ij}{}}F_{ij}\mathrm{\Lambda }_X(f_j,f_i)Y`$
$`=`$ $`[X,FY]+(m1)[FX,Y]+{\displaystyle \underset{ij}{}}F_{ij}\mathrm{\Lambda }_X(f_j,f_i)Y,`$
where $`F_{ij}=_{ki,j}f_k`$. The calculations are based on the Leibniz rule for derivations:
$$\widehat{X}(\underset{i=1}{\overset{m}{}}f_i)=\underset{i=1}{\overset{m}{}}F_i\widehat{X}(f_i),$$
etc. Since, due to Corollary 1, $`\mathrm{\Lambda }_X`$ is skew-symmetric and $`F_{ij}=F_{ji}`$, we have
$$\underset{ij}{}F_{ij}\mathrm{\Lambda }_X(f_j,f_i)Y=0$$
and (a) follows.
(b) In view of (a), we have
$`(m2)[FX,Y]`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}[F_iX,f_iY][X,FY][FX,Y]`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{m1}{}}}[F_iX,f_iY]+[F_mX,f_mY][X,F_mf_mY][F_mf_mX,Y].`$
But
$$[F_mX,f_mY][X,F_mf_mY][F_mf_mX,Y]=[F_mY,f_mX]$$
is a particular case of (a).
$`\mathrm{}`$
## 3 Useful facts about associative algebras
In what follows, $`๐`$ will be an associative commutative unital algebra over a field $`๐`$ of characteristic 0. Our standard model will be the algebra $`๐(N)`$ of class $`๐`$ functions on a manifold $`N`$ of class $`๐`$, $`๐=C^{\mathrm{}},C^๐,`$. Here $`C^{\mathrm{}}`$ refers to the smooth category with $`๐=`$, $`C^๐`$ โ to the $``$-analytic category with $`๐=`$, and $``$ โ to the holomorphic category of Stein manifolds with $`๐=`$ (cf. \[G1, AG\]). All manifolds are assumed to be paracompact and second countable. It is obvious what is meant by a Lie QD-algebroid or a Jacobi-Kirillov bundle of class $`๐`$. The rings of germs of class $`๐`$ functions at a given point are noetherian in analytic cases that is no longer true in the $`C^{\mathrm{}}`$ case. However, all the algebras $`๐(N)`$ are in a sense noetherian in finite codimension. To explain this, let us start with the following well-known observation.
###### Theorem 1
Every maximal finite-codimensional ideal of $`๐(N)`$ is of the form $`\overline{p}=\{f๐(N):f(p)=0\}`$ for a unique $`pN`$ and $`\overline{p}`$ is finitely generated.
Proof.- The form of such ideals is proven e.g. in \[G1\], Proposition 3.5. In view of embedding theorems for all types of manifolds we consider, there is an embedding $`f=(f_1,\mathrm{},f_n):N๐^n`$, $`f_i๐(N)`$. Then, the ideal $`\overline{p}`$ is generated by $`\{f_if_i(p)1:i=1,\mathrm{},n\}`$. In the smooth case it is obvious, in the analytic cases it can be proven by means of some coherent analytic sheaves and methods parallel to those in \[G2\], Note 2.3.
$`\mathrm{}`$
Remark. Note that in the case of a non-compact $`N`$ there are maximal ideals of $`๐(N)`$ which are not of the form $`\overline{p}`$. They are of course of infinite codimension. It is not known if the above theorem holds also for manifolds which are not second countable (cf. \[G8\]).
For a subset $`B๐`$, by $`\text{Sp}(๐,B)`$ we denote the set of those maximal finite-codimensional ideals of $`๐`$ which contain $`B`$. For example, $`\text{Sp}(๐,\{0\})`$ is just the set of all maximal finite-codimensional ideals which we denote shortly by $`\text{Sp}(๐)`$. Put $`\overline{B}=_{I\text{Sp}(A,B)}I`$. For an ideal $`I๐`$, by $`\sqrt{I}`$ we denote the radical of $`I`$, i.e.,
$$\sqrt{I}=\{f๐:f^nI,\text{for some }n=1,2,\mathrm{}\}.$$
The following easy observations will be used in the sequel.
###### Theorem 2
(a) If $`I`$ is an ideal of codimension $`k`$ in $`๐`$, then $`\sqrt{I}=\overline{I}`$ and $`(\overline{I})^kI`$.
(b) Every finite-codimensional prime ideal in $`๐`$ is maximal.
(c) If a derivation $`D\text{Der}(๐)`$ preserves a finite-codimensional ideal $`I`$ in $`๐`$, then $`X(๐)\overline{I}`$.
(d) If $`I_1,\mathrm{},I_n`$ are finite-codimensional and finitely generated ideals of $`๐`$, then the ideal $`I_1\mathrm{}I_n`$ is finite-codimensional and finitely generated.
Proof.- (a) The descending series of ideals
$$I+\overline{I}I+(\overline{I})^2\mathrm{}$$
stabilizes at $`k`$th step at most, so $`I+(\overline{I})^k=I+(\overline{I})^{k+1}`$. Applying the Nakayamaโs Lemma to the finite-dimensional module $`(I+(\overline{I})^k)/I`$ over the algebra $`A/I`$, we get $`(I+(\overline{I})^k)/I=\{0\}`$, i.e., $`(\overline{I})^kI`$, thus $`\overline{I}\sqrt{I}`$. Since for all $`J\text{Sp}(A,I)`$ we have $`\sqrt{I}\sqrt{J}=J`$, also $`\overline{I}\sqrt{I}`$.
(b) If $`I`$ is prime and finite-codimensional, $`\sqrt{I}=I`$ and $`\sqrt{I}=\overline{I}`$ by (a). But a finite intersection of maximal ideals is prime only if they coincide, so $`\overline{I}=J`$ for a single $`J\text{Sp}(๐)`$.
(c) By Lemma 4.2 of \[G1\], $`D(I)I`$ for a finite-codimensional ideal $`I`$ implies $`D(๐)J`$ for each $`J\text{Sp}(๐,I)`$.
(d) It suffices to prove (d) for $`n=2`$ and to use the induction. Suppose that $`I_1,I_2`$ are finite-codimensional and finitely generated by $`\{u_i\}`$ and $`\{v_j\}`$, respectively. It is easy to see that $`I_1I_2`$ is generated by $`\{u_iv_j\}`$ and that $`I_1I_2`$ is finite-codimensional in $`I_1`$. Indeed, if $`c_1,\mathrm{},c_k๐`$ represent a basis of $`A/I_2`$, then $`\{c_lu_i\}`$ represent a basis of $`I_1/(I_1I_2)`$.
$`\mathrm{}`$
###### Theorem 3
For an associative commutative unital algebra $`๐`$ the following are equivalent:
(a) Every finite-codimensional ideal of $`A`$ is finitely generated.
(b) Every maximal finite-codimensional ideal of $`A`$ is finitely generated.
(c) Every prime finite-codimensional ideal of $`A`$ is finitely generated.
Proof.- (a) $``$ (b) is trivial, (b) $``$ (c) follows from Theorem 2 (b), and (c) $``$ (a) is a version of Cohenโs Theorem for finite-dimensional ideals.
$`\mathrm{}`$
Definition 5. We call an associative commutative unital algebra $`๐`$ noetherian in finite codimension if one of the above (a), (b), (c), thus all, is satisfied.
An immediate consequence of Theorem 1 is the following.
###### Theorem 4
The algebra $`๐=๐(N)`$ is noetherian in finite codimension.
## 4 Spectra of Jacobi modules
Let us fix a Jacobi module $`(,[,])`$ over $`(๐,๐)`$. Throughout this section we will assume that $``$ is finitely generated by torsion-free elements and that $`๐`$ is a noetherian algebra in finite codimension over a field $`๐`$ of characteristic 0. The $`(๐,๐(N))`$-modules of sections of class $`๐`$ vector bundles over $`N`$ can serve as standard examples.
For $`L`$, by $`\widehat{L}`$ denote the image of $`L`$ under the anchor map: $`\widehat{L}=\{\alpha (X):XL\}\text{Der}(๐)`$. The set $`\widehat{}`$ is a Lie subalgebra in $`\text{Der}(๐)`$ with the commutator bracket $`[,]_c`$ and we will refer to $`\widehat{}`$ as to the Lie algebra of โJacobi-hamiltonian vector fieldsโ. The main difference with the โmodularโ case (in particular, with that of Pursell and Shanks \[PS\]) is that $`\widehat{}`$ is no longer, in general, an $`๐`$-module, so we cannot multiply by โfunctionsโ inside $`\widehat{}`$. However, we still can try to translate some properties of the Lie algebra $`(,[,])`$ into the properties of the Lie algebra $`\widehat{}`$ of Jacobi-hamiltonian vector fields by means of the anchor map and to describe some โLie objectsโ in $``$ or $`\widehat{}`$ by means of โassociative objectsโ in $`๐`$.
The spectrum of the Jacobi module $``$, denoted by $`\text{Sp}()`$, will be the set of such maximal finite-codimensional Lie subalgebras in $``$ that not contain finite-codimensional Lie ideals of $``$. In nice geometric situations, $`\text{Sp}()`$ will be interpreted as a set of points of the base manifold at which the anchor map does not vanish. Note that the method developed in \[GG\] for Lie pseudoalgebras fails, since the Lemma 1 therein in no longer true for Jacobi modules. In fact, as easily shows the example of a symplectic Poisson bracket on a compact manifold, $`[,]`$ may include no non-trivial $`๐`$-submodules of $``$. Therefore we will modify the method from \[G3\] where Poisson brackets have been considered.
Let us fix some notation. For a liner subspace $`L`$ in $``$ and for $`J๐`$, denote
* $`๐ฉ_L=\{X:[X,L]L\}`$ โ the Lie normalizer of $`L`$;
* $`U_L=\{X:[X,]L\}`$;
* $`I(L)=\{f๐:X[fXL]\}`$ โ the largest associative ideal $`I`$ in $`๐`$ such that $`IL`$;
* $`_J=\{X:\widehat{X}(A)J\}`$.
It is an easy excercise to prove the following proposition (cf. \[G3\], Theorem 1.6).
###### Proposition 3
If $`L`$ is a Lie subalgebra in $``$, then $`๐ฉ_L`$ is a Lie subalgebra containing $`L`$, the set $`U_L`$ is a Lie ideal in $`๐ฉ_L`$, and $`\widehat{๐ฉ_L}(I(U_L))I(U_L)`$.
Choose now generators $`X_1,\mathrm{},X_n`$ of $``$ over $`๐`$. For a fixed finite-codimensional Lie subalgebra $`L`$ in $``$ put $`U_i=\{f๐:fX_iU_L\}`$ and $`U=_{i=1}^nU_i`$. Since $`U_L`$ is clearly finite-codimensional in $``$, all $`U_i`$ are finite-codimensional in $`๐`$, so is $`U`$.
###### Lemma 1
(a) $`[U^mX_j,X_k]L`$ for all $`j,k=1,\mathrm{},n`$ and $`m3`$.
(b) $`[U^mX_j,U^lX_k]L`$ for all $`j,k=1,\mathrm{},n`$ and $`m,l1`$.
Proof.- (a) Take $`f_1,\mathrm{},f_mU`$. Since $`f_iX_kU_L`$, Proposition 2 (b) implies $`[f_1\mathrm{}f_mX_j,X_k]L`$.
(b) The inclusion is trivial for $`l=1`$, so suppose $`l2`$. Take $`f_1,\mathrm{},f_mU`$, $`f_{m+1}U^l`$ and put $`F=f_1\mathrm{}f_{m+1}`$, $`F_i=_{ri}f_r`$. By Proposition 2 (b)
$$[f_1\mathrm{}f_mX_k,f_{m+1}X_j]=(m1)[FX_j,X_k]\underset{i=1}{\overset{m}{}}[F_iX_j,f_iX_k].$$
Since $`FU^{m+l}`$, according to (a), $`[FX_j,X_k]L`$ and $`[F_iX_j,f_iX_k][,U_L]L`$, so the lemma follows.
$`\mathrm{}`$
###### Theorem 5
The ideal $`I(U_L)`$ is finite-codimensional in $`๐`$ provided $`L`$ is a finite-codimensional Lie subalgebra in $``$.
Proof.- Let $`๐ฐ`$ be the associative subalgebra in $`๐`$ generated by $`U`$. It is finite-codimensional and, in view of Lemma 1 (b), $`[๐ฐX_j,๐ฐX_k]L`$. Being finite-codimensional in $`๐`$, the associative subalgebra $`๐ฐ`$ contains a finite-codimensional ideal $`J`$ of $`๐`$ (cf. \[G3\], Proposition 2.1 b)). Hence $`[JX_j,JX_k]L`$ and, since $`X_i`$ are generators of $``$, $`[J,J]L`$. Note that we do not exclude the extremal case $`๐ฐ=๐=J`$. Applying the identity
$$[f_1f_2X,Y]=[f_2X,f_1Y]+[f_1X,f_2Y][f_1f_2Y,X]$$
for $`f_1,f_2J`$, $`XU_L`$, we see that $`J^2U_LU_L`$. In particular, $`J^2UX_iU_L`$ for all $`i=1,\mathrm{},n`$, so $`J^2UU`$ and hence $`J^2๐ฐU`$ and $`J^3U_L`$. Consequently $`J^3I(U_L)`$. Since $`J`$ is finite-codimensional and finitely generated, $`J^3`$ is finite-codimensional (Theorem 2 (d)), so $`I(U_L)`$ is finite-codimensional.
$`\mathrm{}`$
Denote $`\text{Sp}_{}(๐)`$ the set of these maximal finite-codimensional ideals $`I๐`$ which do not contain $`\widehat{}(๐)`$, i.e., $`_I`$. Geometrically, $`\text{Sp}_{}(๐)`$ can be interpreted as the support of the anchor map. Recall that $`\text{Sp}()`$ is the set of these maximal finite-codimensional Lie subalgebras in $``$ which do not contain finite-codimensional Lie ideals.
###### Theorem 6
The map $`J_J`$ constitutes a bijection of $`\text{Sp}_{}(๐)`$ with $`\text{Sp}()`$. The inverse map is $`L\sqrt{I(L)}`$.
Proof.- Let us take $`J\text{Sp}_{}(๐)`$. In view of (2), $`J^2_J`$ which implies that $`_J`$ is finite-codimensional, as $`J^2`$ is finite-codimensional and $``$ is finitely generated.
We will show that $`_J`$ is maximal. Of course, $`_J`$ and $`_J`$ is of finite codimension, so there is a maximal Lie subalgebra $`L`$ containing $`_J`$. We have
$$J^2_JLJ^2I(L)J\sqrt{I(L)}J=\sqrt{I(L)}.$$
Moreover, $`I(L)`$ is finite-codimensional, and since, due to (1), $`\widehat{L}(I(L))I(L)`$, then, by Theorem 2 (c), $`\widehat{L}(A)J`$, i.e., $`L_J`$ and finally $`L=_J`$.
Finally, suppose $`P`$ is a finite-codimensional Lie ideal of $``$ contained in $`_J`$. Then $`U_P`$ is a Lie ideal in $``$ of finite codimension and, according to Theorem 5, $`I(U_P)`$ is a finite-codimensional ideal in $`๐`$. Since $`\widehat{}(I(U_P))I(U_P)`$, and since $`I(U_P)I(U_L)J`$, we have $`\widehat{}(A)J`$, i.e., $`=_J`$; a contradiction.
Suppose now that $`L\text{Sp}()`$. Observe first that $`๐ฉ_L=L`$, since otherwise $`L`$ would be a Lie ideal, that would, in turn, imply $`U_LL`$ and $`I(U_L)I(L)`$. Since $`U_L`$ is finite-codimensional, Theorem 5 shows that $`I(L)`$ is finite-codimensional. Exactly as above we show that $`\widehat{L}(๐)\sqrt{I(L)}`$, i.e., $`L_J`$, where $`J=\sqrt{I(L)}`$. By Theorem 2 (a), $`J^kI(L)`$ for some $`k`$, so if we had $`_J=`$, then $`J^k`$ would be a finite-codimensional Lie ideal contained in $`L`$. Thus $`_J`$ and there is $`I\text{Sp}(๐,J)`$ with $`_I`$. We know already that in this case $`_I`$ is maximal. Since $`L_J_I`$ and $`L`$ is maximal, we have $`L=_I`$ and $`I=J=\sqrt{I(L)}`$.
$`\mathrm{}`$
###### Corollary 2
Let $`(,[,])`$ be a Lie QD-algebroid of class $`๐`$ (i.e., a Jacobi module over $`(๐,๐(N))`$ of class $`๐`$ sections of a class $`๐`$ vector bundle) over a class $`๐`$ manifold $`N`$. Let $`SN`$ be the open support of the anchor map, i.e., $`S=\{pN:\widehat{X}(p)0\text{ for some }X\}`$. Then the map $`pp^{}=\{X:\widehat{X}(p)=0\}`$ constitutes a bijection of $`S`$ with $`\text{Sp}()`$.
Let $`\widehat{}`$ be the image of the anchor map $`\alpha :\text{Der}(๐)`$. By definition of a Jacobi module, $`\widehat{}`$ is a Lie subalgebra in $`(\text{Der}(๐),[,]_c)`$. Since $`\alpha :\widehat{}`$ is a surjective Lie algebra homomorphism, it induces a bijection of $`\text{Sp}()`$ onto $`\text{Sp}(\widehat{})`$, $`L\widehat{L}=\alpha (L)`$. Thus we get the following.
###### Corollary 3
Let $`(,[,])`$ be a Lie QD-algebroid of class $`๐`$ over a class $`๐`$ manifold $`N`$. Let $`SN`$ be the open support of the anchor map, i.e., $`S=\{pN:\widehat{X}(p)0\text{ for some }X\}`$. Then the map $`p\widehat{p}=\{\xi \widehat{}:\xi (p)=0\}`$ constitutes a bijection of $`S`$ with $`\text{Sp}(\widehat{})`$.
## 5 Isomorphisms
It is clear that any isomorphism $`\mathrm{\Psi }:_1_2`$ of the Lie algebras associated with Jacobi modules $`_i`$ over $`(R_i,๐_i)`$, $`i=1,2`$, induces a bijection $`\psi :\text{Sp}(_2)\text{Sp}(_1)`$. Since the kernels $`K_i`$ of the anchor maps $`\alpha _i:_i\widehat{_i}`$ are the intersections
$$K_i=\underset{L\text{Sp}(_i)}{}L,i=1,2,$$
$`\mathrm{\Psi }(K_1)=K_2`$, so $`\mathrm{\Psi }`$ induces a well-defined isomorphisms
$$\widehat{\mathrm{\Psi }}:\widehat{_1}\widehat{_2},\widehat{\mathrm{\Psi }}(\widehat{X})=\widehat{\mathrm{\Psi }(X)}$$
with the property
$$\widehat{L}\text{Sp}(\widehat{_1})\widehat{\mathrm{\Psi }}(\widehat{L})\text{Sp}(\widehat{_2}).$$
(8)
###### Proposition 4
If the Lie algebras $`(_i,[,]_i)`$, associated with Jacobi modules $`_i`$, $`i=1,2`$, are isomorphic, then the Lie algebras of Jacobi-hamiltonian vector fields $`\widehat{_i}`$, $`i=1,2`$, are isomorphic.
The following theorem describes isomorphisms of the Lie algebras of Jacobi-hamiltonian vector fields.
###### Theorem 7
Let $`(_i,[,]_i)`$ be a Lie QD-algebroid of class $`๐`$, over a class $`๐`$ manifold $`N_i`$, and let $`S_iN_i`$ be the (open) support of the anchor map $`\alpha _i:_i\widehat{}_i`$, $`i=1,2`$. Then every isomorphism of the Lie algebras of Jacobi-hamiltonian vector fields $`\mathrm{\Phi }:\widehat{}_1\widehat{}_2`$ is of the form $`\mathrm{\Phi }(\xi )=\phi _{}(\xi )`$ for a class $`๐`$ diffeomorphism $`\phi :S_1S_2`$.
###### Corollary 4
If the Lie algebras associated with Lie QD-algebroids $`E_i`$ of class $`๐`$, over class $`๐`$ manifolds $`N_i`$, $`i=1,2`$, are isomorphic, then the (open) supports $`S_iN_i`$ of the anchor maps $`\alpha _i:_i\widehat{}_i`$, $`i=1,2`$, are $`๐`$-diffeomorphic. In particular, $`N_1`$ and $`N_2`$ are $`๐`$-diffeomorphic provided the anchors are nowhere-vanishing.
Proof of Theorem 7.- According to Corollary 3, the isomorphism $`\mathrm{\Phi }`$ induces a bijection $`\phi :S_1S_2`$ such that, for every $`\xi \widehat{}_1`$ and every $`pS_1`$,
$$\xi (p)=0\mathrm{\Phi }(\xi )(\phi (p))=0.$$
(9)
First, we will show that $`\phi `$ is a diffeomorphism of class $`๐`$. For, let $`f๐(N_1)`$. Since the anchor map is a first-order differential operator, for every $`X_1`$ we have $`\widehat{f^2X}=2f\widehat{fX}f^2\widehat{X}`$. In particular, for any $`pN_1`$,
$$\widehat{f^2X}(p)2f(p)\widehat{fX}(p)+f^2(p)\widehat{X}(p)=0,$$
so that, due to (9),
$$\mathrm{\Phi }(\widehat{f^2X})(\phi (p))=2f(p)\mathrm{\Phi }(\widehat{fX})(\phi (p))f^2(p)\mathrm{\Phi }(\widehat{X})(\phi (p)).$$
(10)
We can rewrite (9) in the form
$$\mathrm{\Phi }(\widehat{f^2X})=2(f\psi )\mathrm{\Phi }(\widehat{fX})(f\psi )^2\mathrm{\Phi }(\widehat{X}),$$
(11)
where $`\psi =\phi ^1`$ and the both sides of (11) are viewed as vector fields on $`S_2`$. In a similar way one can get
$$\mathrm{\Phi }(\widehat{f^3X})=3(f\psi )^2\mathrm{\Phi }(\widehat{fX})2(f\psi )^3\mathrm{\Phi }(\widehat{X}).$$
(12)
To show that $`f\psi `$ is of class $`๐`$, choose $`qS_2`$ and $`X_1`$ such that $`\mathrm{\Phi }(\widehat{X})(q)0`$. Then we can choose local coordinates $`(x_1,\mathrm{},x_n)`$ around $`q`$ such that $`\mathrm{\Phi }(\widehat{X})=_1=\frac{}{x_1}`$. If $`a`$ is the first coefficient of the vector field $`\mathrm{\Phi }(\widehat{fX})`$ in these coordinates, we get out of (11) and (12) that $`(f\psi )^22a(f\psi )`$ and $`2(f\psi )^33a(f\psi )^2`$ are of class $`๐`$ in a neighbourhood of $`q`$. But
$$(f\psi )^22a(f\psi )=(f\psi a)^2a^2$$
(13)
and
$$2(f\psi )^33a(f\psi )^2=2(f\psi a)^3+3a(f\psi a)^2a^2,$$
(14)
so $`(f\psi a)^2`$ and $`(f\psi a)^3`$ are functions of class $`๐`$ in a neighbourhood of $`q`$, as the function $`a`$ is of class $`๐`$. Now we will use the following lemma which proves that $`f\psi a`$, thus $`f\psi `$, is of class $`๐`$.
###### Lemma 2
If $`g`$ is a $`๐`$-valued function in a neighbourhood of $`0๐^n`$ such that $`g^2`$ and $`g^3`$ are of class $`๐`$, then $`g`$ is of class $`๐`$.
Proof.- In the analytic cases the lemma is almost obvious, since $`g=\frac{g^3}{g^2}`$ is a meromorphic and continuous function. In the smooth case the Lemma is non-trivial and proven in \[Jo\].
$`\mathrm{}`$
To finish the proof of the theorem, we observe that $`f\psi `$ is of class $`๐`$ for all $`f๐(N_2)`$ implies that $`\psi `$, thus $`\phi =\psi ^1`$, is of class $`๐`$ and we show that $`\mathrm{\Phi }=\phi _{}`$ or, in other words, that $`\widehat{Y}(f)\psi =\mathrm{\Phi }(\widehat{Y})(f\psi )`$ for all $`f๐(N_1)`$ and all $`Y_1`$. Indeed, for arbitrary $`f๐(N_1)`$ and $`X,Y_1`$, the bracket of vector fields $`[\widehat{Y},\widehat{f^2X}]`$ reads
$`[\widehat{Y},\widehat{f^2X}]`$ $`=`$ $`[\widehat{Y},2f\widehat{fX}f^2\widehat{X}]`$
$`=`$ $`2\widehat{Y}(f)\widehat{fX}2f\widehat{Y}(f)\widehat{X}+2f[\widehat{Y},\widehat{fX}]f^2[\widehat{Y},\widehat{X}].`$
Hence, similarly as in (11),
$`\mathrm{\Phi }([\widehat{Y},\widehat{f^2X}])`$ $`=`$ $`2(\widehat{Y}(f)\psi )\mathrm{\Phi }(\widehat{fX})2(f\psi )(\widehat{Y}(f)\psi )\mathrm{\Phi }(\widehat{X})`$
$`+2(f\psi )\mathrm{\Phi }([\widehat{Y},\widehat{fX}])(f\psi )^2\mathrm{\Phi }([\widehat{Y},\widehat{X}]).`$
Comparing the above with
$$[\mathrm{\Phi }(\widehat{Y}),\mathrm{\Phi }(\widehat{f^2X})]=[\mathrm{\Phi }(\widehat{Y}),2(f\psi )\mathrm{\Phi }(\widehat{fX})(f\psi )^2\mathrm{\Phi }(\widehat{X})],$$
we get easily
$$\left(\widehat{Y}(f)\psi \mathrm{\Phi }(\widehat{Y})(f\psi )\right)\left(\mathrm{\Phi }(\widehat{fX})(f\psi )\mathrm{\Phi }(\widehat{X})\right)=0.$$
(15)
After polarizing with $`f:=f+h`$ and multiplying both sides by $`\widehat{Y}(f)\psi \mathrm{\Phi }(\widehat{Y})(f\psi )`$, we get the identity
$$\left(\widehat{Y}(f)\psi \mathrm{\Phi }(\widehat{Y})(f\psi )\right)^2\left(\mathrm{\Phi }(\widehat{hX})(h\psi )\mathrm{\Phi }(\widehat{X})\right)=0,$$
(16)
valid for all $`f,h๐(N_1)`$ and all $`X,Y_1`$. From (16) we get
$$(\widehat{Y}(f)\psi )(q)=(\mathrm{\Phi }(\widehat{Y})(f\psi ))(q)$$
for such $`q=\phi (p)S_2`$ for which in no neighbourhood of them the anchor map is a differential operator of order 0, i.e., for $`q`$ which do not belong to
$$S_2^0=\{\phi (p)S_2:\widehat{hX}(p^{})=h(p^{})\widehat{X}(p^{})\text{ for all }h๐(N_1),X_1\text{ and }p^{}\text{ close to }p\}.$$
If, on the other hand, $`qS_2^0`$, then $`\mathrm{\Phi }(\widehat{hX})(q^{})=(h\psi )(q^{})\mathrm{\Phi }(\widehat{X})(q^{})`$ for $`q^{}`$ from a neighbourhood of $`q`$, so that comparing in this neighbourhood
$$\mathrm{\Phi }([\widehat{Y},\widehat{fX}])=(\widehat{Y}(f)\psi )\mathrm{\Phi }(\widehat{X})+(f\psi )\mathrm{\Phi }([\widehat{Y},\widehat{X}])$$
with
$$[\mathrm{\Phi }(\widehat{Y}),\mathrm{\Phi }(\widehat{fX})]=\mathrm{\Phi }(\widehat{Y})(f\psi )\mathrm{\Phi }(\widehat{X})+(f\psi )[\mathrm{\Phi }(\widehat{Y}),\mathrm{\Phi }(\widehat{X})]$$
we get
$$(\widehat{Y}(f)\psi )(q)\mathrm{\Phi }(\widehat{X})(q)=\mathrm{\Phi }(\widehat{Y})(f\psi )(q)\mathrm{\Phi }(\widehat{X})(q),$$
thus
$$(\widehat{Y}(f)\psi )(q)=\mathrm{\Phi }(\widehat{Y})(f\psi )(q)$$
also for $`qS_2^0`$.
$`\mathrm{}`$
Remark.
(a) For Jacobi-Kirillov bundles with all leaves of the characteristic foliation (i.e., orbits of $`\widehat{}`$) of dimension $`>1`$ there is much simpler argument showing that $`\psi `$ is smooth than the one using Lemma 2. The difficulty in the general case comes from singularities of the โbivector fieldโ part of the anchor map and forced us to use Lemma 2.
(b) Theorem 7 has been proven for Lie algebroids in \[GG\], so the new (and difficult) here is the case of Jacobi-Kirillov bundles with non-trivial โbivector partโ of the bracket. A similar result for the Lie algebras of smooth vector fields preserving a symplectic or a contact form up to a multiplicative factor has been proven by H. Omori \[O1\]. These Lie algebras are the Lie algebras of locally hamiltonian vector fields for the Jacobi-Kirillov brackets associated with the symplectic and the contact form, respectively.
###### Corollary 5
(a) If the Lie algebras $`(๐(N_i),\{,\}_{\beta _i})`$ of the Jacobi contact brackets, associated with contact manifolds $`(N_i,\beta _i)`$, $`i=1,2`$, of class $`๐`$, are isomorphic, then the manifolds $`N_1`$ and $`N_2`$ are $`๐`$-diffeomorphic.
(b) If the Lie algebras associated with nowhere-vanishing Poisson structures of class $`๐`$ on class $`๐`$ manifolds $`N_i`$, $`i=1,2`$, are isomorphic, then the manifolds $`N_1`$ and $`N_2`$ are $`๐`$-diffeomorphic.
Janusz GRABOWSKI
Institute of Mathematics
Polish Academy of Sciences
ลniadeckich 8
P.O. Box 21
00-956 Warsaw, Poland
Email: jagrab@impan.gov.pl |
warning/0506/physics0506096.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The genesis of the Post constraint on the electromagnetic constitutive relations of linear mediums was described in detail quite recently . This structural constraint was shown to arise from the following two considerations:
* Two of the four Maxwell postulates (containing the induction fields and the sources) should be independent of the other two Maxwell postulates (containing the primitive fields) at the macroscopic level, just as the two sets of postulates are mutually independent at the microscopic level.
* The constitutive functions must be characterized as piecewise uniform, being born of the spatial homogenization of microscopic entities. Therefore, if a certain constitutive function of a homogeneous piece of a medium cannot be recognized by proper electromagnetic experimentation, the assumption of a continuously nonhomogeneous analog of that constitutive function is untenable.
Available experimental evidence against the validity of the Post constraint for linear materials was shown to be incomplete and inconclusive, in addition to being based either on the physically inadmissible premise of purely instantaneous response and/or derived from a preโmodern version of electromagnetism .
Nevertheless, solutions of very simple (frequencyโdomain) boundary value problems can be invoked very easily to claim the invalidity of the Post constraint for linear materials. Indeed, when a boundary value problem involving a homogeneous linear material is formulated to assess the validity of the Post constraint, a conflict arises between the fundamental differential equations of electromagnetism in the chosen material and a naรฏve application of the usual boundary conditions. In this paper, that conflict is easily resolved โ in favor of the Post constraint.
The organization of this paper is as follows: Section 2 contains a brief review of modern macroscopic electromagnetism, followed by a relevant presentation of linear constitutive relations in Section 3. The principal equations of a naรฏve formulation of boundary value problems are set up in Section 4, and the aforementioned conflict is presented and resolved in Section 5. The paper concludes with some remarks in Section 6.
## 2 Modern Macroscopic Electromagnetism
Let us begin with the fundamental equations of modern electromagnetism. The microscopic fields are just two: the electric field $`\underset{ยฏ}{\overset{~}{e}}(\underset{ยฏ}{x},t)`$ and the magnetic field $`\underset{ยฏ}{\overset{~}{b}}(\underset{ยฏ}{x},t)`$.<sup>2</sup><sup>2</sup>2The lowerโcase letter signifies that a field or a source density is microscopic, while the tilde $`\stackrel{~}{}`$ indicates dependence on time. Furthermore, $`ฯต_o=8.854\times 10^{12}`$ F/m and $`\mu _o=4\pi \times 10^7`$ H/m are the permittivity and the permeability of matterโfree space in the absence of an external gravitational field (which condition is assumed here). These two are accorded the status of primitive fields in modern electromagnetism, and their sources are the microscopic charge density $`\stackrel{~}{c}(\underset{ยฏ}{x},t)`$ and the microscopic current density $`\underset{ยฏ}{\overset{~}{j}}(\underset{ยฏ}{x},t)`$. Both fields and both sources appear in the microscopic Maxwell postulates
$`\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{e}}(\underset{ยฏ}{x},t)=ฯต_0^1\stackrel{~}{c}(\underset{ยฏ}{x},t),`$ (1)
$`\times \underset{ยฏ}{\overset{~}{b}}(\underset{ยฏ}{x},t)ฯต_0\mu _0{\displaystyle \frac{}{t}}\underset{ยฏ}{\overset{~}{e}}(\underset{ยฏ}{x},t)=\mu _0\underset{ยฏ}{\overset{~}{j}}(\underset{ยฏ}{x},t),`$ (2)
$`\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{b}}(\underset{ยฏ}{x},t)=0,`$ (3)
$`\times \underset{ยฏ}{\overset{~}{e}}(\underset{ยฏ}{x},t)+{\displaystyle \frac{}{t}}\underset{ยฏ}{\overset{~}{b}}(\underset{ยฏ}{x},t)=\underset{ยฏ}{0}.`$ (4)
Spatial averaging of the microscopic primitive fields and source densities yields the macroscopic Maxwell postulates
$`\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x},t)=ฯต_0^1\stackrel{~}{\rho }(\underset{ยฏ}{x},t),`$ (5)
$`\times \underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)ฯต_0\mu _0{\displaystyle \frac{}{t}}\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x},t)=\mu _0\underset{ยฏ}{\overset{~}{J}}(\underset{ยฏ}{x},t),`$ (6)
$`\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)=0,`$ (7)
$`\times \underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x},t)+{\displaystyle \frac{}{t}}\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)=\underset{ยฏ}{0},`$ (8)
which involve the macroscopic primitive fields $`\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x},t)`$ and $`\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)`$ as well as the macroscopic source densities $`\stackrel{~}{\rho }(\underset{ยฏ}{x},t)`$ and $`\underset{ยฏ}{\overset{~}{J}}(\underset{ยฏ}{x},t)`$. Equations (5)โ(8) are the fundamental (differential) equations of modern macroscopic electromagnetism. Let us note that
* all four equations contain only two fields, both primitive, and
* all four equations hold in matterโfree space as well as in matter.
Indeed, modern electromagnetism may be called EBโelectromagnetism to indicate the central role of $`\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x},t)`$ and $`\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)`$.
Equations (5)โ(8) are not, however, the textbook form of the Maxwell postulates. In order to obtain that familiar form, source densities are decomposed into free and bound components, and the bound components are then quantified through the polarization and the magnetization, both of which are in turn subsumed in the definitions of the electric induction $`\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)`$ and the magnetic induction $`\underset{ยฏ}{\overset{~}{H}}(\underset{ยฏ}{x},t)`$. Then, (5)โ(8) metamorphose into the following familiar form:
$`\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)=\stackrel{~}{\rho }_{so}(\underset{ยฏ}{x},t),`$ (9)
$`\times \underset{ยฏ}{\overset{~}{H}}(\underset{ยฏ}{x},t){\displaystyle \frac{}{t}}\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)=\underset{ยฏ}{\overset{~}{J}}_{so}(\underset{ยฏ}{x},t),`$ (12)
$`\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)=0,`$
$`\times \underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x},t)+{\displaystyle \frac{}{t}}\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)=\underset{ยฏ}{0}.`$
Here, $`\stackrel{~}{\rho }_{so}(\underset{ยฏ}{x},t)`$ and $`\underset{ยฏ}{\overset{~}{J}}_{so}(\underset{ยฏ}{x},t)`$ represent free or externally impressed source densities. Let us note that $`\underset{ยฏ}{\overset{~}{H}}(\underset{ยฏ}{x},t)`$ and $`\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)`$ do not have microscopic counterparts and therefore are not considered fundamental in modern electromagnetism.
## 3 Linear Constitutive Relations
The most general linear constitutive relations may be written as
$`\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)={\displaystyle \underset{ยฏ}{\overset{~}{\underset{ยฏ}{ฯต}}}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x}\underset{ยฏ}{x}_h,tt_h)๐\underset{ยฏ}{x}_h๐t_h}`$
$`+{\displaystyle \underset{ยฏ}{\overset{~}{\underset{ยฏ}{\alpha }}}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x}\underset{ยฏ}{x}_h,tt_h)๐\underset{ยฏ}{x}_h๐t_h}`$
$`+{\displaystyle \stackrel{~}{\mathrm{\Phi }}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x}\underset{ยฏ}{x}_h,tt_h)๐\underset{ยฏ}{x}_h๐t_h}`$ (13)
and
$`\underset{ยฏ}{\overset{~}{H}}(\underset{ยฏ}{x},t)={\displaystyle \underset{ยฏ}{\overset{~}{\underset{ยฏ}{\beta }}}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x}\underset{ยฏ}{x}_h,tt_h)๐\underset{ยฏ}{x}_h๐t_h}`$
$`+{\displaystyle \underset{ยฏ}{\overset{~}{\underset{ยฏ}{\nu }}}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\text{ }\text{}\text{ }\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x}\underset{ยฏ}{x}_h,tt_h)๐\underset{ยฏ}{x}_h๐t_h}`$
$`{\displaystyle \stackrel{~}{\mathrm{\Phi }}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\underset{ยฏ}{\overset{~}{E}}(\underset{ยฏ}{x}\underset{ยฏ}{x}_h,tt_h)๐\underset{ยฏ}{x}_h๐t_h}`$ (14)
wherein the integrals extend only over the causal values of $`(\underset{ยฏ}{x}_h,t_h)`$ in relation to $`(\underset{ยฏ}{x},t)`$. Five constitutive functions are present in the two foregoing equations: $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{ฯต}}}`$ is the permittivity tensor; $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\nu }}}`$ is the impermeability tensor; $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\alpha }}}`$ and $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\beta }}}`$ are the magnetoelectric tensors such that
$$\mathrm{Trace}\left[\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\alpha }}}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\beta }}}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)\right]0;$$
(15)
and $`\stackrel{~}{\mathrm{\Phi }}`$ may be called the Tellegen parameter.
When (13) and (14) are substituted in (9)โ(12) to retain only the primitive fields and the source densities, the resulting four equations contain $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{ฯต}}}`$, $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\alpha }}}`$, $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\beta }}}`$ and $`\underset{ยฏ}{\overset{~}{\underset{ยฏ}{\nu }}}`$ in two ways:
* by themselves, and
* through their spaceโ and timeโderivatives.
In contrast, $`\stackrel{~}{\mathrm{\Phi }}`$ does not occur by itself, but only in terms of derivatives . The elimination of this anomalous situation leads to the Post constraint
$$\stackrel{~}{\mathrm{\Phi }}(\underset{ยฏ}{x},t;\underset{ยฏ}{x}_h,t_h)0.$$
(16)
Arguments in favor of and against the Post constraint were cataloged some years ago , with the opposing arguments based on the soโcalled EH electromagnetism wherein $`\underset{ยฏ}{\overset{~}{H}}(\underset{ยฏ}{x},t)`$ is regarded as the primitive magnetic field and $`\underset{ยฏ}{\overset{~}{B}}(\underset{ยฏ}{x},t)`$ as the induction magnetic field. The EHโelectromagnetism is a preโmodern formulation that is still widely used in frequencyโdomain research. Opposing arguments of a similar nature have also been made under the rubric of the heterodox EDBHโelectromagnetism , wherein $`\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)`$ and $`\underset{ยฏ}{\overset{~}{H}}(\underset{ยฏ}{x},t)`$ are also supposed to have microscopic counterparts and are therefore also considered primitive.
## 4 Boundary Value Problems
Constitutive functions are macroscopic entities arising from the homogenization of assemblies of microscopic bound source densities, with matterโfree space serving as the reference medium . In any small enough portion of matter that is homogenizable, the constitutive functions are uniform. When such a portion will be interrogated for characterization, it will have to be embedded in matterโfree space. Typically, macroscopically homogeneous matter is characterized in the frequency domain. Hence, it is sensible to investigate if the Tellegen parameter can be determined by such a measurement.
Without loss of generality, let us consider therefore that all space is divided into two regions, $`V_+`$ and $`V_{}`$, separated by a boundary $`S`$. The region $`V_+`$ is not filled with matter, whereas the region $`V_{}`$ is filled with a spatially homogeneous, temporally invariant and spatially local matter characterized by the constitutive relations
$`\begin{array}{cc}\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{\underset{ยฏ}{ฯต}}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{\underset{ยฏ}{\alpha }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )+\mathrm{\Phi }(\omega )\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{H}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{\underset{ยฏ}{\beta }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{\underset{ยฏ}{\nu }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\mathrm{\Phi }(\omega )\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )\hfill & \end{array}\},`$ (19)
$`\underset{ยฏ}{x}V_{},`$ (20)
where $`\omega `$ is the angular frequency, and $`\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )`$ is the temporal Fourier transform of $`\underset{ยฏ}{\overset{~}{D}}(\underset{ยฏ}{x},t)`$, etc.
The frequencyโdomain differential equations
$$\begin{array}{c}\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )=0\hfill \\ \times \underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )i\omega \underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{0}\hfill \end{array}\},\underset{ยฏ}{x}V_+V_{},$$
(21)
are applicable in both $`V_+`$ and $`V_{}`$, with $`i=\sqrt{1}`$.
The remaining two Maxwell postulates in matterโfree space may be written as
$$\begin{array}{c}ฯต_0\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )=\rho _{so}(\underset{ยฏ}{x},\omega )\hfill \\ \mu _0^1\times \underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )+i\omega ฯต_0\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{J}_{so}(\underset{ยฏ}{x},\omega )\hfill \end{array}\},\underset{ยฏ}{x}V_+,$$
(22)
in terms of only the macroscopic primitive fields, with sources that are sufficiently removed from the boundary $`S`$ . The fields $`\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )`$ and $`\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )`$ in $`V_+`$ can be represented using standard techniques , and the representations of $`\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )=ฯต_0\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )`$ and $`\underset{ยฏ}{H}(\underset{ยฏ}{x},\omega )=\mu _0^1\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )`$ in $`V_+`$ then follow.
In $`V_{}`$, the remaining two Maxwell postulates are expressed as follows:
$$\begin{array}{c}\text{ }\text{}\text{ }\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )=0\hfill \\ \times \underset{ยฏ}{H}(\underset{ยฏ}{x},\omega )+i\omega \underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{0}\hfill \end{array}\},\underset{ยฏ}{x}V_{},$$
(23)
Substituting (20) therein, we obtain
$`\text{ }\text{}\text{ }\left[\underset{ยฏ}{\underset{ยฏ}{ฯต}}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{\underset{ยฏ}{\alpha }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\right]`$
$`+\mathrm{\Phi }(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )=0,\underset{ยฏ}{x}V_{},`$ (24)
and
$`\times \left[\underset{ยฏ}{\underset{ยฏ}{\beta }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )\right]+i\omega \underset{ยฏ}{\underset{ยฏ}{ฯต}}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )`$
$`+\times \left[\underset{ยฏ}{\underset{ยฏ}{\nu }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\right]+i\omega \underset{ยฏ}{\underset{ยฏ}{\alpha }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )`$
$`\mathrm{\Phi }(\omega )\left[\times \underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )i\omega \underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\right]=\underset{ยฏ}{0},\underset{ยฏ}{x}V_{}.`$ (25)
These equations simplify to
$$\text{ }\text{}\text{ }\left[\underset{ยฏ}{\underset{ยฏ}{ฯต}}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{\underset{ยฏ}{\alpha }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\right]=0,\underset{ยฏ}{x}V_{},$$
(26)
and
$`\times \left[\underset{ยฏ}{\underset{ยฏ}{\beta }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )\right]+i\omega \underset{ยฏ}{\underset{ยฏ}{ฯต}}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )`$
$`+\times \left[\underset{ยฏ}{\underset{ยฏ}{\nu }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\right]+i\omega \underset{ยฏ}{\underset{ยฏ}{\alpha }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{0},\underset{ยฏ}{x}V_{},`$ (27)
by virtue of (21). For many classes of materials and shapes of $`S`$, $`\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )`$ and $`\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )`$ in $`V_{}`$ can also be adequately represented ; and thereafter so can be $`\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )`$ and $`\underset{ยฏ}{H}(\underset{ยฏ}{x},\omega )`$ in $`V_{}`$.
In order to solve the boundary value problem, the boundary conditions
$$\begin{array}{cc}\underset{ยฏ}{B}^{norm}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{B}^{norm}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{D}^{norm}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{D}^{norm}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{E}^{tan}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{E}^{tan}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{H}^{tan}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{H}^{tan}(\underset{ยฏ}{x},\omega )\hfill & \end{array}\},\underset{ยฏ}{x}S,$$
(28)
have to be imposed on the boundary $`S`$. Here, $`\underset{ยฏ}{B}^{norm}(\underset{ยฏ}{x}\pm ,\omega )`$ indicate the normal components of $`\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )`$ on either side of $`S`$, whereas $`\underset{ยฏ}{E}^{tan}(\underset{ยฏ}{x}\pm ,\omega )`$ denote the tangential components of $`\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )`$ similarly, etc. Some resulting set of equations can then be solved to determine the scattering of an incident field by the material contained in $`V_{}`$.
Much effort is not required to solve the simplest boundary value problems. Relevant to the Post constraint, reference is made to two papers wherein the boundary $`S`$ is a specularly smooth plane of infinite extent . More complicated boundaries have also been tackled . The inescapable conclusion from examining the results of boundary value problems is that the fields scattered in $`V_+`$ by the material contained in $`V_{}`$ are affected by the Tellegen parameter (if any). Yet that conclusion is naรฏve and incorrect, as we see next.
## 5 The Conflict and Its Resolution
We have two very sharply contrasting Statements emanating from the foregoing frequencyโdomain exercise:
* The Tellegen parameter $`\mathrm{\Psi }`$ vanishes from the fundamental equations (21), (26) and (27) for the material of which the chosen scatterer is made.
* The fields scattered by the chosen scatterer contain a signature of the Tellegen parameter (if any).
In other words, the Tellegen parameter is a ghost: it does not have a direct existence in the fundamental differential equations, but its presence may be indirectly gleaned from a scattering measurement.
The ghostly nature of the Tellegen parameter is a consequence of the boundary conditions (28)<sub>2</sub> and (28)<sub>4</sub>. Even more specifically, it arises from the representations of $`\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )`$ and $`\underset{ยฏ}{H}(\underset{ยฏ}{x},\omega )`$ in $`V_{}`$. It is instructive to decompose the macroscopic induction fields as
$$\begin{array}{c}\underset{ยฏ}{D}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{D}_{actual}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{D}_{excess}(\underset{ยฏ}{x},\omega )\hfill \\ \underset{ยฏ}{H}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{H}_{actual}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{H}_{excess}(\underset{ยฏ}{x},\omega )\hfill \end{array}\},\underset{ยฏ}{x}V_{},$$
(29)
where
$$\begin{array}{c}\underset{ยฏ}{D}_{actual}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{\underset{ยฏ}{ฯต}}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{\underset{ยฏ}{\alpha }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\hfill \\ \underset{ยฏ}{H}_{actual}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{\underset{ยฏ}{\beta }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )+\underset{ยฏ}{\underset{ยฏ}{\nu }}(\omega )\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\hfill \end{array}\},\underset{ยฏ}{x}V_{},$$
(30)
are retained in (26) and (27). On the other hand,
$$\begin{array}{c}\underset{ยฏ}{D}_{excess}(\underset{ยฏ}{x},\omega )=\mathrm{\Phi }(\omega )\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )\hfill \\ \underset{ยฏ}{H}_{excess}(\underset{ยฏ}{x},\omega )=\mathrm{\Phi }(\omega )\underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )\hfill \end{array}\},\underset{ยฏ}{x}V_{},$$
(31)
are filtered out of (26) and (27) by (21) but do affect the boundary conditions (28)<sub>2</sub> and (28)<sub>4</sub>.
The fundamental differential equations in $`V_{}`$ can now be written as follows:
$$\begin{array}{c}\text{ }\text{}\text{ }\underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )=0\hfill \\ \times \underset{ยฏ}{E}(\underset{ยฏ}{x},\omega )i\omega \underset{ยฏ}{B}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{0}\hfill \\ \text{ }\text{}\text{ }\underset{ยฏ}{D}_{actual}(\underset{ยฏ}{x},\omega )=0\hfill \\ \times \underset{ยฏ}{H}_{actual}(\underset{ยฏ}{x},\omega )+i\omega \underset{ยฏ}{D}_{actual}(\underset{ยฏ}{x},\omega )=\underset{ยฏ}{0}\hfill \end{array}\},\underset{ยฏ}{x}V_{}.$$
(32)
Boundary conditions in electromagnetics emerge from the fundamental equations . Therefore, consistently with (32), the correct boundary conditions on $`S`$ are
$$\begin{array}{cc}\underset{ยฏ}{B}^{norm}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{B}^{norm}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{D}^{norm}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{D}_{actual}^{norm}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{E}^{tan}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{E}^{tan}(\underset{ยฏ}{x},\omega )\hfill & \\ \underset{ยฏ}{H}^{tan}(\underset{ยฏ}{x}+,\omega )=\underset{ยฏ}{H}_{actual}^{tan}(\underset{ยฏ}{x},\omega )\hfill & \end{array}\},\underset{ยฏ}{x}S,$$
(33)
instead of (28). Thus the correct formulation of the boundary value problem involves (33)<sub>2</sub> and (33)<sub>4</sub> instead of (28)<sub>2</sub> and (28)<sub>4</sub>.
To sum up, the conflict between Statements A and B arises from a naรฏve and incorrect formulation of the boundary value problem. The correct formulation does not contain $`\underset{ยฏ}{D}_{excess}(\underset{ยฏ}{x},\omega )`$ and $`\underset{ยฏ}{H}_{excess}(\underset{ยฏ}{x},\omega )`$ in $`V_{}`$ as well as in the boundary conditions.
## 6 Concluding Remarks
Any field that cannot survive in the fundamental differential equations is superfluous. Neither $`\underset{ยฏ}{H}_{excess}(\underset{ยฏ}{x},\omega )`$ nor $`\underset{ยฏ}{D}_{excess}(\underset{ยฏ}{x},\omega )`$ survives, and may therefore be discarded ab initio. The Post constraint thus removes the nonuniqueness inherent in (20), not to mention in (13) and (14), which can appear in two of the four Maxwell postulates in relation to the other two postulates. No wonder, de Lange and Raab could recently complete a major exercise โ whereby a multipole formulation of linear materials that was initially noncompliant with the Post constraint was made compliant.
In addition, the Post constraint also removes two anomalies: the first is that of a constitutive function not appearing by itself but only through its derivatives ; the second is that of the Tellegen โmediumโ which is isotropic (i.e., with directionโindependent properties) but wherein propagation characteristics in antiparallel directions are different.
A simple exercise shows that isolated magnetic monopoles can negate the validity of the Post constraint , but the prospects of observing such a magnetic monopole are rather remote . Furthermore, although the electromagnetic characterization of matterโfree space, even in the context of general relativity, is compliant with the Post constraint , the axion concept renders that constraint invalid . No axions have yet been detected however . Finally, available data on magnetoelectric materials seems to negate the Post constraint , but that data is faulty as it is based on the neglect of causality and a false manipulation of the Onsager principle . Needless to add, if either an isolated magnetic monopole or an axion is ever discovered, or if a magnetoelectric material is properly characterized to have the electromagnetic properties claimed for it by virtue of misapplications of various principles, the Post constraint would be invalidated and the basics of EBโelectromagnetism would have to thought anew.
Acknowledgment. Occasional discussions with Dr. E.J. Post are gratefully acknowledged. Thanks are also due to the Department of Management Communication, University of Waikato, Hamilton, New Zealand, for hospitality during a visit when this paper was written. |
warning/0506/cond-mat0506179.html | ar5iv | text | # Resistance distribution in the hopping percolation model
## I Introduction
The concepts and methods of percolation theory are widely used to explain many phenomena in physics, classical as well as quantum (for review see e.g. Refs. ShklEfr ; HavlinBook ; Ahar ). The canonical model for studying transport properties of disordered systems is percolation on a lattice. Usually it is also assumed that the conductivity between neighboring lattice sites may be defined as either finite, or zero (i.e., either conducting or insulating) without loss of generality. This model \[which we denote as the bond (or site) percolation\] has been extensively studied and is understood quite well. For the description of the nearest neighbor hopping in granular materials it is much more natural to define the conductivity between two neighboring lattice sites (labeled as โiโ and โjโ) by $`\sigma _{ij}\mathrm{exp}[r_{ij}/r_0ฯต_{ij}/k_BT]`$, where $`r_{ij}`$ is the distance between the two sites, $`r_0`$ is the scale over which the wave-function outside the grain decays, $`ฯต_{ij}`$ is the energy difference between grains, and $`T`$ is the temperature. Here we neglect the thermal hopping term (high temperature regime) and consider only nearest neighbor hopping. This behavior may be captured by a lattice model for which StrelBerk ; Halperin ; Tyc ; Le ; Sar
$$\sigma _{ij}=\sigma _0\mathrm{exp}[\kappa r(ij)],$$
(1)
where $`\kappa `$ is a measure of disorder, $`r(ij)`$ is a random number taken from uniform distribution in the range (0,1), and $`\sigma _0`$ is a dimension coefficient Halperin ; StrelBerk . We shall name this model the hopping percolation model.
One might expect that such small differences in the formulation of the disorder in the $`\sigma _{ij}`$ (namely, $`\sigma _{ij}=0,1`$ or $`\sigma _{ij}=\sigma _0\mathrm{exp}[\kappa r(ij)]`$) will lead to no important difference in the global conductance properties of these systems. Quite surprisingly, recent experiments on the conductance of granular material Cohen might suggest otherwise. Specifically, the number of red bonds (which are critical for current) expected in framework of the traditional percolation theory is proportional to $`L^{1/\nu }`$, where $`\nu `$ is the percolation correlation critical exponent and $`L`$ is the system size Con . Thus, for a typical experimental set-up of $`10^9`$ grains, one expects $`10^2`$ red bonds. On the other hand, measurements of transport through Ni granular ferromagnets, indicates a much lower number of red bonds (typically of order one) Cohen . In a recent paper StrelBerk we have attributed this difference to the fact that the estimation of the number of red bonds ($`L^{1/\nu }`$) is based on the bond percolation model, while for the hopping percolation we expect a transition to a regime of extreme strong disorder in which a single red bond governs the behavior of the system. The onset of this regime scales as to $`\kappa /L^{1/\nu }`$.
It is important to note that in contrast to the traditional bond (or site) percolation model, in which the system is either a metal or an insulator, for the hopping percolation model the system always conducts some current. Hopping conductivity (i.e., exponential local resistance) (1) is always associated with strong disorder. As was shown in Ref. StrelBerk , there are two regimes within this strong disorder: a regime which is not sensitive to the removal of a single bond, as expected from the usual percolation theory, termed the usual strong disorder regime rem . While for even stronger disorder a regime which is very sensitive to the removal of a specific single bond exists, denote as extreme disorder. In the extreme disorder regime a single bond can determine the transport properties of the entire macroscopic system StrelBerk ; Cohen .
The remainder of this paper is arranged as follows. In Section II we describe our model and the numerical approach. In Section III we present some numerical results, followed by a brief discussion in Section IV.
## II Model
We perform large-scale Monte Carlo simulations for calculating transport in these systems. We build a bond-percolating Miller-Abrahams like resistor network (see Fig. 1 and Ref. StrelBerk, ; MillerAbrahams, ; Kirkpatrick, ; Sarychev2, ), but assume the conductivity of each resistor to have the form given in Eq. (1). Then solve the corresponding set of linear Kirchhoff equations and calculate the total effective resistance $`\rho _e`$ for two dimensional (2D) and three dimensional (3D) networks StrelBerk ; MillerAbrahams ; Kirkpatrick ; Sarychev2 .
We begin by calculating the average effective conductivity $`\sigma _e`$. The approximate expression Halperin ; Ahar ; ShklEfr for the effective conductivity $`\sigma _e`$, of a random resistor network Kirkpatrick with local conductivities given by Eq. (1), in 2D is
$$\sigma _e=\sigma _0e^{p_c\kappa }.$$
(2)
In Fig. 2 we show this dependence (in terms of resistivity $`\rho _e=1/\sigma _e`$) for both site and bond percolations (see Fig. 1) for different lattice sizes. In Ref. Tyc it was shown that in the limit $`\kappa \mathrm{}`$, Eq. (2) is exact. It is easy to show that in the case of 2D random resistor bond network (for which $`p_c=0.5`$), Eq. (2) follows immediately from the Keller-Dykhne theorem and is exact for arbitrary $`\kappa `$ Keller2 . Similarly this results can be found in framework of the symmetric self-consistency effective-medium approximation (EMA) rewritten for many component composite Landauer . From Eq. (2) follows that the effective conductivity $`\sigma _e`$ depends on $`\kappa `$ and does not depend on the system size $`L`$. For finite $`L`$, Eq. (2) represents the mean conductivity over all configurations of the disordered system.
Next we study the fluctuations of the resistance, $`\rho =1/\sigma `$, from the mean value $`\rho _e=1/\sigma _e`$ for individual systems of finite size $`L`$. We perform numerical calculations of the probability distribution function $`P(\rho )`$ (i.e., the probability that the total resistance of the system is $`\rho `$) as well as the variance $`\mathrm{var}(\rho )`$ as a function of $`L`$ and $`\kappa `$. As shown in Ref. StrelBerk , the relative variance (in contrast to $`\rho _e`$), strongly depends on $`L`$ and $`\kappa `$ only through the scaled variable $`hL/\kappa ^\nu `$. Here $`\nu `$ is the critical exponent of the percolation correlation length $`\xi (pp_c)^\nu `$ (in 2D $`\nu =4/31.33`$, while in 3D $`\nu 0.88`$ ShklEfr ; HavlinBook ; Ahar ). It was also shown StrelBerk that $`h`$ describes the transition from strong disorder ($`h>1`$) to extreme disorder ($`h<1`$).
## III Results
Here we present numerical results suggesting that $`P(\rho )`$ also depends only on $`L/\kappa ^\nu `$. In order to verify and to quantify this hypothesis, we study numerically $`P(\rho )`$ for systems of different sizes $`L`$ and different disorder $`\kappa `$, but with the same value of $`h`$ (see also Ref. Wu ). In Figs. 3 and 4 we show $`P(\rho )`$ vs. $`\rho `$ for the cases of strong and extreme disorder. All data corresponding to the same parameter $`h`$ scale according to the same law \[see Figs. 3(c) and 3(g)\]. Thus, our results suggest that $`P(\rho )`$ is a function of both $`\rho /\rho _e`$ and $`h`$, i.e.,
$$P(\rho )=\frac{1}{\rho _e}f(\frac{\rho }{\rho _e};h).$$
(3)
Here $`h`$ determines the form of the function and $`P(\rho )`$ for a fixed $`h`$ depends only on $`\rho /\rho _e`$. Fig. 3(c) suggests that in the strong disorder, $`h`$ controls the width or standard deviation of the rescaled distribution. Since the standard deviation increases when $`h`$ decreases, we assume that the standard deviation $`\delta b\rho _e/h=b\rho _e\kappa ^\nu /L`$, where $`b`$ is a parameter which depends on the type of lattice. Indeed, when we plot in Fig. 3(d) $`P(\rho /\rho _e)\delta `$ vs. $`(\rho \rho _e)/\delta `$, a collapse of the two plots shown in Fig. 3(c) is obtained. The functional form obtained in Fig. 3(d) suggests that the probability distribution can be approximated by a Gaussian
$$P(\rho )(\sqrt{2\pi }\delta )^1\mathrm{exp}\left[\left(\rho \rho _e\right)^2/2\delta ^2\right].$$
(4)
Indeed, the dashed line in Fig. 3(d) represents a good fit to the Gaussian given by Eq. (4). However, Eq. (4) can not approximate the asymmetric form of $`P(\rho )`$ at extreme disorder \[see Figs. 3(g) and 3(h)\]. We suggest, as will be justified below, that $`P(\rho )`$ can be approximated (in all regimes of disorder) by the log-normal form
$$P(\rho )\frac{1}{\sqrt{2\pi }\mu \rho }\mathrm{exp}\left[\frac{\mathrm{ln}^2(\rho /\rho _e)}{2\mu ^2}\right],$$
(5)
where $`\mu =\delta /\rho _e=b\kappa ^\nu /L`$. In fact, Eq. (5) includes also the usual strong disorder case, since in the latter case $`\mathrm{ln}^2(\rho /\rho _e)(\rho /\rho _e1)^2`$ and Eq. (5) reduces to the Gaussian form (4), while at extreme disorder ($`\mu 1`$) the exponent function in Eq. (5) tends to 1, and $`P(\rho )`$ transforms to the power-like dependence $`1/\rho `$ (see Fig. 4).
From Eq. (5) it follows that at $`\mu 0`$, the distribution function $`P(\rho )`$ reduces to a delta-function: $`\underset{\mu 0}{lim}P(\rho )=\frac{1}{\rho }\underset{\mu 0}{lim}\frac{1}{\sqrt{2\pi }\mu }e^{\frac{\mathrm{ln}(\rho /\rho _e)^2}{2\mu ^2}}=\frac{1}{\rho }\delta (\mathrm{ln}\rho \mathrm{ln}\rho _e)=\delta (\rho \rho _e)`$. Therefore, at $`\mu 0`$ (i.e., $`\kappa 0`$ or $`L\mathrm{}`$) the total resistivity of the system is exactly $`\rho _0e^{p_c\kappa }`$ and has no size dependence: $`\underset{\mu 0}{lim}\rho P(\rho )๐\rho =\rho _e`$.
It should be noted that a log-normal distribution of resistances is found in quantum models of hopping conductivities (see e.g., Ref. Shapiro and references therein), while here it is demonstrated for classical exponential disorder (1). Moreover, our result (5) yield the specific analytical form of $`P(\rho )`$, which includes the dependence on $`\kappa `$ and $`L`$ for all regimes of disorder.
In Fig. 5(a) we test Eq. (5) by comparing it to simulation results. It is shown that the numerical results of the 2D resistance $`\rho P(\rho /\rho _e)/(L/\kappa ^\nu )`$ scale vs. $`(\rho /\rho _e)^{L/\kappa ^\nu }`$, as predicted by Eq. (5) for both strong and extreme disorder. A similar plot is presented in Fig. 5(b) for a 3D lattice. Although for the 3D case Eq. (2) is not exact (since Keller-Dykhne theorem exists only in 2D), nevertheless the approximated expression $`\sigma _e\sigma _0\kappa ^\nu e^{p_c\kappa }`$ is known Halperin ; Tyc , resulting in the distribution law (5). Since in Eq. (2) the parameter $`\kappa `$ appears with the prefactor $`p_c`$, we should expect that $`p_c`$ enters into the parameter $`\mu `$ of Eq. (5) as $`\mu =\alpha (p_c\kappa )^\nu /L`$. Comparing the values $`b=0.2`$ observed for the square bond percolation lattice ($`p_c=0.5`$) and $`b=0.18`$ for the cubic site percolation ($`p_c=0.3116`$), we find that $`\alpha =0.503`$. The dependence of $`\mu `$ on $`p_c`$ is in agreement with result of Ref. Kalis . These results strongly support our proposition that Eq. (5) describes well the distribution in all ranges of disorder.
The variance $`\mathrm{var}(\rho )`$ can be expressed as $`\rho ^2\rho _e^2`$, where $`\rho ^n=_{\rho _{\mathrm{min}}}^{\rho _{\mathrm{max}}}\rho ^nP(\rho )๐\rho `$, and $`\rho _e=\rho `$. For large enough $`\kappa `$, $`\rho ^n=_0^{\mathrm{}}\rho ^nP(\rho )๐\rho =\rho _0e^{n\mathrm{ln}\rho _e+\frac{1}{2}n^2\mu ^2}`$ and the relative variance takes the form
$$\left[\mathrm{var}(\rho )\right]^{1/2}/\rho _e=[e^{\mu ^2}(e^{\mu ^2}1)]^{1/2}$$
(6)
(see Ref. Aitchison ). Fig. 6 presents numerical results showing that the relative variance scales as a function of $`\mu =b\kappa ^\nu /L`$ in accordance with Eq. (6).
Next we shall present analytical arguments for the log-normal distribution (5). According to the central limit theorem Aitchison , if the values of $`\mathrm{ln}\rho `$ are normally distributed, then the values of $`\rho `$ should follow the log-normal distribution. Assuming $`\mathrm{ln}\rho =\kappa p_c`$ \[see Eq. (2)\] for all $`\rho `$, the distribution $`P(\kappa ^1\mathrm{ln}\rho )`$ is simply the distribution of the percolation threshold $`\mathrm{\Phi }(p_c)`$ which is normally distributed (e.g., Refs. Lev ). Indeed, in Fig. 5(c) we show that $`\mathrm{\Phi }[\kappa ^1\mathrm{ln}(\rho )]`$ approximately follows a normal distribution centered at $`p_c=0.5`$. Thus, the distribution of $`\rho `$ should be log-normal as in Eq. (5).
Using the above assumption $`\mathrm{ln}\rho =\mu /\kappa =\kappa p_c`$, it is possible to evaluate the distribution $`\mathrm{\Phi }(p_c)`$ and its standard deviation $`\delta _{p_c}`$. One can write a simple relation $`\mathrm{\Phi }(p_c)dp_c=\mathrm{\Phi }^{}(y)dy`$, where $`y=\kappa p_c`$, and get $`\mathrm{\Phi }(p_c)=\mathrm{\Phi }^{}(y)\frac{dy}{dp_c}=\kappa \mathrm{\Phi }(y)`$. Therefore, $`\mathrm{\Phi }(p_c)=\kappa \mathrm{\Phi }^{}(\mathrm{ln}\rho )\kappa \mu ^1\mathrm{exp}[\kappa ^2(p_c\overline{p}_c)^2/2\mu ^2]`$ $`=\delta _{p_c}^1\mathrm{exp}[(p_c\overline{p}_c)^2/2\delta _{p_c}^2]`$, where $`\overline{p}_c`$ is the mean value of the percolation threshold and $`\delta _{p_c}=\mu /\kappa `$ is the standard deviation of $`\mathrm{\Phi }(p_c)`$. This form of $`\mathrm{\Phi }(p_c)`$ is supported by our numerical simulations shown in Figs. 5(c) and 5(d).
This specific form for $`\delta _{p_c}=b\kappa ^{\nu 1}/L`$ in the hopping percolation model should be compared to $`\delta _{p_c}=L^{1/\nu }`$ known for the bond percolation model Lev ; Con . This further emphasizes the differences between the bond percolation model considered in Ref. Lev and the hopping percolation model considered here.
## IV Summary
In summary, we find the specific form of the resistance distribution in the hopping percolation model. For all ranges of strong disorder $`\kappa `$ and lattice sizes $`L`$, the distribution is log-normal and depends only on the ratio $`\kappa ^\nu /L`$, where $`\nu `$ is the correlation exponent for the bond percolation case. Assuming the relation $`\rho =\mathrm{exp}(\kappa p_c)`$ for finite systems leads to a variance of $`p_c`$, $`\delta _{p_c}=\kappa ^{\nu 1}/L`$, which is different from $`\delta _{p_c}=L^{1/\nu }`$ known for the bond percolation model Lev . Our results may be relevant to ac conductivity measurements in such systems. By appropriate choice of frequency one can detect regions of size smaller than $`\kappa ^\nu `$, where a crossover in behavior from extreme to usual strong disorder behavior is expected.
###### Acknowledgements.
This research was supported in part by grants from the US-Israel Binational Science Foundation, the Israel Science Foundation, and the KAMEA Fellowship program of the Ministry of Absorption of the State of Israel. |
warning/0506/astro-ph0506094.html | ar5iv | text | # GRAPE-SPH Chemodynamical Simulation of Elliptical Galaxies II: Scaling Relations and the Fundamental Plane
## 1 INTRODUCTION
The internal structure of galaxies, the spectrophotometric, chemical, and dynamical properties at various locations within a galaxy, is determined by the processes of galaxy formation and evolution. Stars in a galaxy are fossils; the star formation and chemical enrichment history of the galaxy are imprinted on their kinematics and chemical abundances. The SAURON project with the William Herschel Telescope (Bacon et al. (2001); Emsellem et al. (2004)) is providing wide-field mapping of the kinematics and stellar populations of nearby galaxies, which will certainly give stringent constraints on galaxy formation and evolution. Multiobject and integral field spectrographs are being developed also on 8-10m ground-based telescopes, which will provide the time evolution of such internal structure. To infer the physical evolution processes from the observational data, it is necessary to construct a realistic model, i.e., a three-dimensional chemodynamical model, and to compare the theoretical predictions with such observational data.
How elliptical galaxies form is a long-standing matter of debate. Two competing scenarios for the formation of elliptical galaxies have been proposed; the monolithic collapse (e.g., Larson (1974); Arimoto & Yoshii (1987)), and the major merger (e.g., Toomre (1977); Kauffmann, White & Guiderdoni (1993); Baugh, Cole & Frenk (1996); Steinmetz & Navarro (2002)). In Kobayashi (2004, hereafter K04), we constructed a self-consistent three-dimensional chemodynamical model of ellipticals, introducing various physical processes associated with the formation of stellar systems; radiative cooling, star formation, feedback of Type II and Ia supernovae (SNe II and SNe Ia), and stellar winds (SWs), and chemical enrichment. We then argued that both formation processes should arise to explain the observed variation in radial metallicity gradients. The metallicity is enhanced in the central dense region, and the metallicity gradients are generated. However, because merging events weaken the metallicity gradients, and because the secondary star burst induced by the mergers is not enough to regenerate them, galaxies that form monolithically have steeper gradients, while galaxies that undergo major mergers have shallower gradients. Therefore no correlation is found between the mass and gradients, as in the observation (Kobayashi & Arimoto (1999)).
While the internal structure of elliptical galaxies is greatly affected by their merging histories, their global properties should be determined from their masses according to the scaling relations. Differences in merging history may provide the scatter in these relations. In this paper, we investigate whether our simulated galaxies follow the observed correlations; the Faber-Jackson relation, the Kormendy relation, the colour-magnitude relation, the mass-metallicity relation and the fundamental plane (FP).
The FP is a correlation of early-type galaxies with $`2+n`$ parameters (e.g., Djorgovski & Davis (1987); Dressler et al. (1987)) that reflect the internal structures, and is a clue to understand formation and evolution of early-type galaxies. One possible interpretation of the FP, defined by central velocity dispersions $`\sigma _0`$, absolute effective radii r<sub>e</sub>, and surface brightnesses within an effective radius SB<sub>e</sub> attributes it to a correlation of the mass-to-light ratio $`M/L`$ to the total luminosity, or equivalently, to the total galaxy mass (e.g., Faber et al. (1987)). The dependence of $`M/L`$ on the mass and luminosity stems from the stellar metallicity and/or age (Pahre, Djorgovski & de Carvalho (1998)). However, elliptical galaxies may not be homologous along the FP, and there remains a dispersion that is not due to observational error. This may be caused by metallicity, age, and/or dynamical disturbance. The FP is observed up to $`z0.5`$ in clusters (Kelson et al. (1997), 2000), which is understood as the evidence for passive evolution since $`z<1`$. For field galaxies, a comparable correlation is observed up to $`z1`$ (Treu et al. (2001); Gebhardt et al. (2003); van de Ven, van Dokkum & Franx (2003); van der Wel et al. (2004)). The larger scatter than in clusters and the zero-point offset at higher redshifts can be interpreted as an age difference, although other possibilities such as dynamical disturbance have not yet been discussed.
The details of our GRAPE-SPH chemodynamical code were described in K04, and we briefly summarize them in ยง2. Following the discussion in K04, we vary the parameters controlling star formation and the initial mass function, and we show how these affect global properties and their correlations (ยง3.2). With the best parameter set, we show the scaling relations of simulated galaxies comparing with the observations. We focus on the fundamental plane, and discuss the origin of the scatter in ยง3.3. ยง4 and ยง5 respectively contain the discussion and our conclusions.
## 2 CHEMODYNAMICAL MODEL
The characteristics of our GRAPE-SPH code may be summarized as follows (see K04 for the detail).
i) The SPH method (Monaghan (1992) for a review) is adopted, and the gravity is calculated in direct summation using the special purpose computer GRAPE (Sugimoto et al. (1990)). The SPH formulation used in the code is almost the same as in Navarro & White (1993). The GRAPE-SPH code was originally written by Nakasato & Nomoto (2003), and is highly adaptive in space and time through individual smoothing lengths and individual timesteps. The calculations were done with the GRAPE5 system in the National Astronomical Observatory of Japan and the GRAPE6 of the University of Tokyo.
ii) Radiative cooling is computed using a metallicity-dependent cooling function. For primordial gas (\[Fe/H\] $`<5`$), we compute the cooling rates using the two-body processes of H and He, and free-free emission, as in Katz, Weinberg & Hernquist (1996). For metal enriched gas (\[Fe/H\] $`5`$), we use a metallicity-dependent cooling function computed with the MAPPINGS III software (Sutherland & Dopita (1993)). In this cooling function, the elemental abundance ratios are set to be constant for given \[Fe/H\] according to the relations found in the solar neighborhood. \[O/Fe\]$`=0.5`$ for Galactic halo stars for \[Fe/H\] $`1`$, and solar values for \[Fe/H\] $`0`$. We interpolate between these values for $`1<`$ \[Fe/H\] $`<0`$.
iii) Our star formation criteria are the same as in Katz (1992); (1) converging flow; $`(๐)_i<0`$, (2) rapid cooling; $`t_{\mathrm{cool}}<t_{\mathrm{dyn}}`$, and (3) Jeans unstable gas; $`t_{\mathrm{dyn}}<t_{\mathrm{sound}}`$. The star formation timescale is proportional to the dynamical timescale ($`t_{\mathrm{sf}}\frac{1}{c}t_{\mathrm{dyn}}`$), where $`c`$ is our star formation timescale parameter. We also adopt the probability criterion (Katz (1992)); A random number between 0 and 1 is compared with the probability $`P1\mathrm{exp}\left[\frac{\mathrm{\Delta }t_{\mathrm{sf}}}{t_{\mathrm{sf}}}\right]`$ in a time interval $`\mathrm{\Delta }t_{\mathrm{sf}}=2`$ Myr.
If a gas particle satisfies the above star formation criteria, a fractional part of the mass of the gas particle turns into a star particle. Since an individual star particle has a mass of $`10^{57}M_{}`$, it dose not represent a single star, but an association of many stars. The mass of the stars associated with each star particle is distributed according to an initial mass function (IMF). We adopt a power-law IMF, $`\varphi (m)m^x`$ (the slope $`x=1.35`$ gives the Salpeter IMF), which is invariant to time and metallicity.
In ยง3.2, we discuss the dependence of our results on our free parameters; the star formation timescale $`c`$ and the slope of the IMF $`x`$. In order to reproduce the observed radius-magnitude relation and mass-metallicity relation, we will choose $`c=0.1`$ and $`x=1.35`$ as a standard model.
iv) For the feedback of energy and heavy elements, we do not adopt the instantaneous recycling approximation. Via SWs, SNe II, and SNe Ia, thermal energy and heavy elements are ejected from an evolved star particle as functions of time, and are distributed to all surrounding gas particles out to a constant radius of 1 kpc. The ejected energy of each SW, SN II, and SN Ia are $`0.2\times 10^{51}`$ erg depending on metallicity, $`1.4\times 10^{51}`$ erg, and $`1.3\times 10^{51}`$ erg. We distribute this feedback energy in purely thermal form, although a fraction of it (given by a free parameter $`f_{\mathrm{kin}}`$) can be distributed in kinetic form as a velocity perturbation to the gas particles (see Navarro & White (1993)). As shown in Fig.14 of K04, if we adopt $`f_{\mathrm{kin}}=0.1`$, the star formation efficiency is lower, the surface brightness of the final galaxy decreases at the centre, and metal-rich gas blows out, resulting in effective radii which are too large and metallicity gradients which are too shallow.
For the metals, the mass-dependent nucleosynthesis yields of SNe II and SNe Ia are taken from Nomoto et al. 1997ab. The upper and lower mass limits of the IMF are $`0.05`$ and $`120M_{}`$, respectively. The progenitor mass ranges of SWs and SNe II are $`8120M_{}`$ and $`850M_{}`$, respectively. For SNe Ia, we adopt the single degenerate scenario with the metallicity effect (Kobayashi et al. (1998), 2000), where the progenitors are the Chandrasekhar WDs with an initial mass of $`38M_{}`$, and the lifetimes are determined from the lifetimes of the secondary stars with $`0.91.5M_{}`$ and $`1.82.6M_{}`$ for the red-giants and main-sequence systems, respectively.
v) The photometric evolution of a star particle is identical to the evolution of a simple stellar population. Spectra $`f_\lambda `$ are taken from Kodama & Arimoto (1997) as a function of age $`t`$ and metallicity $`Z`$.
vi) The initial condition is a slowly rotating sphere with a CDM initial fluctuation generated by the COSMICS package (Bertschinger (1995)). The cosmological parameters are set to be $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_m=1.0`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, and $`\sigma _8=1.0`$. The initial angular momentum is added as rigid rotation with the constant spin parameter $`\lambda 0.02`$. We set different resolutions and total mass; the total mass of $`10^{12}M_{}`$ (baryon fraction of $`0.1`$) with comoving radius of $`1.5`$ Mpc, and with $`N10000`$ and $`60000`$ (the half for gas and the rest for dark matter), which are the same as in K04. The mass of a dark matter particle is $`1.8\times 10^8M_{}`$ and $`3.0\times 10^7M_{}`$, and the mass of a gas particle is $`2.0\times 10^7M_{}`$ and $`3.3\times 10^6M_{}`$, respectively. In this paper, we add a new sample for cD galaxies with the total mass of $`10^{13}M_{}`$ from wider initial conditions with radius $`3`$ Mpc and $`N60000`$. This mass resolution is similar to the lower resolution of K04 sample.
## 3 RESULTS
### 3.1 EVOLUTION HISTORIES
Simulating the chemodynamical evolution of 74 fields with different cosmological initial conditions, we obtain 128 galaxies at the present time (i.e., $`t=13`$ Gyr). Depending on the CDM initial fluctuations, in some cases one galaxy forms in the centre of the field, in the others one galaxy and several subgalaxies. We select galaxies having stellar masses in a 20 kpc sphere larger than $`4.5\times 10^7M_{}`$. Although many less-massive subgalaxies form, we discard them because our resolution is not enough to study them in detail. We summarize the number of runs and the resulting galaxies in Table 1.
Different galaxies undergo different evolution histories. The difference is seeded in the initial conditions. Galaxies form through the successive merging of subgalaxies with various masses. These vary between a major merger at one extreme and a monolithic collapse of a slowly rotating gas cloud at the other. We classify galaxies into the 5 classes according to their merging histories, defining a major merger with $`f>0.2`$ at $`z<3`$; \[E1\] monolithic, \[E2\] assembly, \[E3\] minor merger, \[E4\] major merger, and \[E5\] multiple major mergers (see K04 for the detail). The numbers of galaxies of each class is slightly different from the adopted parameter set as summarized in Table 2, although the galaxy identification method is the same in K04. The fractions of major merger galaxies among elliptical galaxies are 46%, 54%, and 59% respectively for the model A, B, and C. This is because the number of merging events increase at $`1<z<3`$ with slower star formation (smaller $`c`$) and weaker feedback (smaller $`x`$).
As well as the observed dwarfs, the variation of the star formation histories is seen for our simulated dwarfs with $`M_{\mathrm{V},\mathrm{tot}}>19`$ mag. We classify these into the 4 classes; \[D1\] initial starburst, \[D2\] continuous star formation, \[D3\] continuous star formation with recent starburst, and \[D4\] recent starburst. Observationally, the first class of galaxies would be dwarf ellipticals, the others are dwarf irregulars. The star formation rate (SFR) is truncated in \[D1\] because of the supernova feedback. The intermittent SFR is induced by the gas accretion and/or interaction with other galaxies. In the following sections, we show the results both for giant and dwarf galaxies, but we should note that our numerical resolution in the dwarf galaxies is not enough; the particle number in a galaxy is small ($`200`$ in the worst case) and the gravitational softening ($`1.0`$ kpc for the low resolution) is comparable to the size of galaxies.
All simulated ellipticals form with an initial starburst at $`z>2`$ with the typical timescale of $`12`$ Gyr. The SFR decreases because the gas is exhausted in the galaxy. The secondary starburst is induced by the accretion of gas clumps and/or the merging of gas-rich galaxies. Not all merging events induce a secondary starburst; the fraction of merging events induce such a starburst is about $`10\%`$, depending on the gas mass of the secondary galaxy. The initial starburst is always larger than the secondary one.
As discussed in K04, such a truncating SFR is due to the artificial cut-off of mass accretion caused by the vacuum boundary of our initial conditions. However, if we simulate wider region, the star formation continues longer, and thus colours tend to be too blue. The mass accretion can be continued but the star formation should have stopped by some process by $`z2`$ in cluster ellipticals (e.g., Kodama & Arimoto (1997); Ellis et al. (1997); Kodama et al. (1998); Stanford, Eisenhardt & Dickinson (1998); Brown et al. (2000)) and $`z1`$ in field ellipticals (e.g., Silva & Bothun (1998); Schade et al. (1999); Brinchmann & Ellis (2000); Daddi, Cimatti & Renzini (2000); Im et al. (2002)). An analogous truncation is required such as tidal stripping and effects of active galactic nuclei (AGN) (see ยง4). Eventually, although the dark matter halo may be affected, the stellar population does not change so much.
Contrary to the observational constraints, galactic winds are hard to generate in our simulations, and star formation does not completely stop by $`z=0`$ in most simulated elliptical galaxies. At the centre of the present-day galaxy, the dynamical potential is so deep that the gas density is high. In such regions, super metal-rich stars ($`Z10Z_{}`$) keep on forming in the simulation. In several dwarf galaxies with $`M<`$$`10^9M_{}`$, weak galactic winds can be seen even if $`f_{\mathrm{kin}}=0`$. By the present epoch, roughly $`25\%`$ of the baryons are blown away by the input of thermal energy from supernovae, but most of the heavy elements remain locked into stars in the galaxy (see ยง3.2).
### 3.2 PARAMETER DEPENDENCES
We show the parameter dependence of the global properties of a galaxy for varying star formation parameter $`c`$ and slope of the IMF $`x`$ in Figures 1 and 2. The present values are summarized in Table 3. We adopt identical initial condition for a non-major merger galaxy (ID 782389913, \[E1\] monolithic-like) with high resolution. The model A (solid lines, $`c=1.0,x=1.10`$) is the same as in K04, where the feedback is stronger than the others because the star formation takes place at earlier epoch before the gas cloud collapses enough. With the model B (dashed lines, $`c=0.1,x=1.10`$), the star formation takes place more slowly, and thus the stellar metallicity increases a lot. With the model C (dotted lines, $`c=0.1,x=1.35`$), which is adopted in the following sections as a standard model, the stellar yield is reduced with smaller $`x`$, and thus the feedback is weaker than in the other models. Here we define the galaxy as $`r<r_{200}`$, i.e., the region with the mean density higher than $`200\rho _\mathrm{c}`$ at each time.
Although the time evolution of the size and total mass (panels a and b) do not change at all, the stellar mass (panel c) varies with $`c`$ and $`x`$ by a factor of $`20\%`$ and $`5\%`$, respectively, in the sense that stronger feedback decreases the stellar mass. The baryon fraction (panel d) is controlled by $`c`$. Although the dark matter contraction is independent of these parameters, $`c`$ changes the peak epoch of the initial starburst ($`0.3`$ and $`0.6`$ Gyr respectively for $`c=1.0`$ and $`0.1`$), which determines when and where stars form in the contracting gas cloud. On the other hand, the present gas fraction (panel e) does not depend on $`c`$, but on $`x`$. This is because the return fraction is smaller with $`x=1.35`$ by a factor of $`2040\%`$ than for $`x=1.10`$ depending on the lower mass limit $`M_{\mathrm{}}`$ of the IMF; 0.33 ($`x=1.35`$, $`M_{\mathrm{}}=0.05M_{}`$, model C), 0.42 ($`1.35`$, $`0.1M_{}`$), 0.45 ($`1.10`$, $`0.05M_{}`$, model A), and 0.52 ($`1.10`$, $`0.1M_{}`$, model B). How much gas turns to stars finally depends both on $`c`$ and $`x`$ because some fraction of the gas is ejected from the galaxy.
Although the residual gas mass in the galaxy mainly depends on $`x`$, its temperature is also affected by $`c`$. Figure 2 shows the time evolution of the gas mass in the galaxy and in the wind for all gas (panels a, d), hot gas with $`T>10^6`$ K (panels b, e), and cold gas with $`T<10^4`$ K (panels c, f). With the strongest feedback (model A, solid lines), the half of the residual gas is hot, and no cold gas exists. With a longer timescale of star formation (model B, dashed lines), the residual gas can cool and the temperature decreases. With model C (dotted lines), the gas cooling is not efficient because of lower metallicity.
In this simulation, although ten times larger amount of gas exists outside the galaxy ($`r>2r_{200}`$), some of them never accrete onto the galaxy. We define the wind gas as the gas particles that have been inside $`r<r_{200}`$ and are outside $`r>2r_{200}`$ at present. The wind gas mass is almost the half of the outside gas mass. Although the temperature of the wind gas is different (panels e, f), the ejected wind mass by the present is independent of these parameters (panel d) with an ejected wind fraction, the ratio of the wind mass to the accreted baryon mass $`M_{\mathrm{wind}}/M_{\mathrm{acc}}`$, of $`25\%`$. The wind efficiency, the ratio of the wind mass to the total stellar mass $`M_{\mathrm{wind}}/M_{}`$, depends on $`c`$, and is $`0.4`$ and $`0.35`$ for $`c=1.0`$ and $`0.1`$. This means that it is easier for gas particles to be blown away when the gas collapse has not yet finished.
The biggest effect of $`x`$ is on the metallicity, and the mean stellar metallicity of the whole of the galaxy (panel f) is $`2.2`$, $`2.7`$, and $`0.7Z_{}`$ respectively for model A, B, and C. The observational estimate using the metallicity gradients (Kobayashi & Arimoto (1999)) is $`0.51`$ solar. Although $`x=1.35`$ is favored to meet this observational constraint under our star formation and feedback schemes, we cannot rule out other IMF. In principle, $`x=1.10`$ gives stronger feedback, which can reduce the stellar metallicity. However, this feedback mechanism do not work well in our simulations and almost all the metals produced by supernovae are locked into stars. Because the metallicity-dependent cooling function is included in our simulation, the metal enriched gas easily cools and loses the energy needed to escape. In spite of the strong dependence of metallicity on $`x`$, the ejected metal fraction (the ratio of metal mass in the wind to the total produced metals $`M_{\mathrm{Z},\mathrm{wind}}/M_{\mathrm{Z},\mathrm{acc}}`$) depends only on $`c`$, and is $`2\%`$ and $`1\%`$ for $`c=1.0`$ and $`0.1`$, respectively. With larger $`c`$, metals start to be lost before the galaxy collapses enough, and thus more metals can be ejected. We should note, however, that these numbers are much smaller than expected from observations of X-ray gas in clusters. These suggest that the ejected metal fraction should be larger than two-thirds (Renzini (2002)).
As discussed in K04, the galaxies of the K04 simulation are too extended, which causes an offset in the radius-magnitude relation. This is because star formation takes place too early before the gas can accrete towards the centre. This problem can be solved by changing $`c`$. Figure 3 shows the star formation rates derived from the ages of stars in the present-day galaxy (top panel), the surface brightness profile (middle panel), and the oxygen abundance gradients (bottom panel). The peak epoch of the initial star burst is delayed from $`0.3`$ Gyr with $`c=1.0`$ (solid lines) to $`0.6`$ Gyr with $`c=0.1`$ (dashed and dotted lines). With $`c=0.1`$, the surface brightness increases at the centre, which results in smaller $`r_\mathrm{e}`$ (4.3, 2.7, and 2.2 kpc respectively for models A, B, and C). The metallicity gradient does not change.
We show scaling relations for our simulated galaxies in Figure 4. The effective radius $`r_\mathrm{e}`$ is derived by fitting the Vaucouleursโ law to the projection of $`|Z|100`$ kpc on the $`X`$-$`Y`$ plane (see K04 for the detail). Figure 4 shows the relations of effective radius $`r_\mathrm{e}`$ (upper panels) and luminosity weighted metallicity (lower panels) to total V-magnitude. The solid and dashed lines show the observed relations for giant ellipticals (Pahre (1999)) and dwarf ellipticals, respectively (Binggeli, Sandage & Tarenghi 1984). With model A (left panels), there is an off-set in the radius-magnitude relation because galaxies are too extended. The slope of the mass-metallicity relation is the same as observed, although the metallicity is too high. With model B (middle panels), the radius becomes smaller by a factor of three, which agrees well with the observed radius-magnitude relation. However, the metallicity is still too high and the slope of the mass-metallicity relation is too shallow. With model C (right panels), the radius-magnitude relation remains in good agreement with observations, and the metallicity decreases to meet the observations. Although the dispersion is larger, the simulated galaxies also follow the observed mass-metallicity relation.
### 3.3 CORRELATIONS
#### 3.3.1 Scaling Relations
For elliptical galaxies, it is well known that there are various correlations among physical scale parameters. The best known of these is the Faber-Jackson (1976) relation $`L\sigma ^n`$, where the slope is $`n4`$, but with a variation, depending on the sample definition (Kormendy & Djorgovski (1989)). A correlation between the de Vaucouleursโ effective radius $`r_\mathrm{e}`$ and surface brightness SB<sub>e</sub> was found by Kormendy (1977), where more luminous galaxies have larger $`r_\mathrm{e}`$ and fainter SB<sub>e</sub>. However, the scatter of SB<sub>e</sub> for a given $`\sigma `$ is quite large, and the SB$`{}_{\mathrm{e}}{}^{}\sigma `$ diagram corresponds to a face-on view of the fundamental plane (Kormendy & Djorgovski (1989)). For dwarf galaxies, these relations are different; the $`M_Br_\mathrm{e}`$ relation has a shallower slope, and the $`M_B\mathrm{SB}_\mathrm{e}`$ relation reverses with luminous dwarfs having brighter SB<sub>e</sub> (Binggeli et al. (1984)).
Figure 5 shows the scaling relations between total stellar mass $`M_{}`$ in $`r_{200}`$, total V-band luminosity $`M_{\mathrm{V},\mathrm{tot}}`$ derived from the de Vaucouleursโ law, central velocity dispersion $`\sigma _0`$ in 2 kpc, effective radius $`r_\mathrm{e}`$, and mean surface brightness $`\mathrm{SB}_\mathrm{e}`$ within $`r_\mathrm{e}`$ in the V-band. All points are for the simulated galaxies, and massive and luminous galaxies are on the right side in all diagrams. The symbols show the merging histories for elliptical galaxies and star formation histories for dwarf galaxies as listed in the figure caption.
(panels a, b) The total stellar mass ($`r<r_{200}`$) correlates well with the total luminosity derived from the de Vaucouleursโ fit (panel a) and with the central velocity dispersion (panel b). The correlations for the stellar mass measured in $`r_\mathrm{e}`$ have a scatter which is twice as large.
(panel c) The Faber-Jackson relation $`L\sigma ^n`$ shows a smaller dispersion than observed (gray points). The slope is $`n=3.8`$ for this observation, but is steeper with $`n=2.6`$ for the simulation. This is because the simulated dwarf galaxies tend to have smaller $`\sigma _0`$ because of lack of resolution. The slope for $`21<M_V<19.5`$ mag is consistent with $`2.7`$ both for the simulation and the observation.
(panels d, e, f) The mass-effective radius relation: Massive/luminous ellipticals have larger effective radii. The simulated galaxies follow the observed relation in the Coma cluster (Pahre (1999), solid line). The dispersion is not so small as $`0.7`$ dex, but is comparable to the observation with $`0.5`$ dex (gray points). In the simulation, the surface brightness of the central part is smeared by the gravitational softening, which causes an uncertainty in fitting the de Vaucouleursโ law. For dwarf ellipticals, the observed relation has a shallower slope than that of giant ellipticals (Binggeli et al. (1984), dashed line). This tendency can be seen in the simulated dwarfs although the scatter is larger.
(panels g, h, i) The mass-surface brightness relation: For giant ellipticals, massive/luminous galaxies tend to have smaller surface brightness, but this relation is reversed for dwarf ellipticals. These tendency can be reproduced in the simulation, but the dispersion is very large.
(panel j) The surface brightness-effective radius relation: The simulated giant ellipticals follow the observed relation where larger galaxies have lower surface brightnesses in $`r_\mathrm{e}`$. The scatter is almost the same as observed. For dwarf galaxies, the observed relation is rectangular, and larger dwarfs have higher surface brightnesses. The simulated dwarf galaxies populate the same side of the observed relation, but the direction of the relation is different from observation. The effective radii of the simulated dwarfs tend to be too large, which is due to the lack of resolution.
The observed scaling relations are reproduced in the simulations. The scatter exists even if the uncertainties of the simulations are taken account. The origin of the scatter is clearly demonstrated by the symbols. The galaxies that form monolithically (filled circles and squares) have smaller effective radii $`r_\mathrm{e}`$ and thus brighter surface brightness SB<sub>e</sub>, while the galaxies that undergo major mergers (open circles and squares) have larger $`r_\mathrm{e}`$ and thus fainter SB<sub>e</sub>. This is because the merging events destroy the galaxy structures and make the radius larger and larger. The dynamical information is not fully wiped out but is blurred by a merger, as is shown using the difference of the energies of particles before and after merging events in Figs.12 and 13 of K04. Therefore, we conclude that the scatter of the scaling relations stems from differences in merging histories.
#### 3.3.2 Mass-Metallicity Relation
The colour-magnitude relation of elliptical galaxies is usually interpreted as luminous galaxies having higher stellar metallicities. This is supported by the observation that a colour-magnitude relation with the same slope is found for high-redshift cluster ellipticals (e.g., Kodama & Arimoto (1997)). Previous studies showed that the line index (Mg<sub>2</sub>)<sub>0</sub> correlate with the velocity dispersion $`\sigma _0`$ at the galaxy centre (e.g., Davies et al. (1987); Bender, Burstein & Faber 1993) and also with total absolute magnitude (e.g., Burstein et al. (1988)). The relation with $`\sigma _0`$ is tighter, but still shows some intrinsic scatter. The same (Mg<sub>2</sub>)<sub>0</sub>-$`\sigma _0`$ relation is found both for cluster and field ellipticals (Bernardi et al. (1998)). Therefore, the mass-metallicity relation of ellipticals is a common relation independent of environment and time, which contains important information on the star formation and feedback processes during the early stages of galaxy formation.
Figure 6 shows the sequences of stellar populations against the galaxy mass. The mass tracers are the total absolute V-magnitude and the central velocity dispersion. The characteristics of stellar populations are expressed by B-V and V-K colours, stellar abundances of oxygen \[O/H\] and iron \[Fe/H\], and stellar age, all of which are measured within $`r_\mathrm{e}`$ and weighted by V-luminosities. Compared with observation, the mass-metallicity relations (panels e-h) are weak, with shallower slope and larger scatter in the simulations, although the average is consistent. These are because the thermal feedback of supernovae is not enough to terminate star formation in the SPH simulations. Especially, ejection of the metal enriched gas do not work well. Although the steep slope of the observational relation requires not only more metal ejection in less-massive galaxies, but also more metal production in massive galaxies, such process cannot occur even if we change the IMF (ยง3.2). For the scatter, we find no significant dependence on merging history; it is mainly caused by the age differences because the luminosity weighted metallicity is affected by bright young populations formed in the secondary starbursts. The relations for iron show larger scatter because iron is produced mainly by SNe Ia with longer lifetimes and affected by late star formation (see ยง4 for the abundance ratios).
Nonetheless, the majority of stars in the simulated giant galaxies are formed in an initial starburst, and the luminosity-weighted ages are as old as $`710`$ Gyr. No relation is found between age and galaxy mass (panels i-j). For dwarf galaxies, the ages decrease to $`38`$ Gyr with a large scatter and a trend with mass can be seen, which is consistent with the observation using $`H\gamma _\sigma `$ index (Yamada et al (2004)). The scatter stems from the differing star formation histories of dwarfs. Among dwarf galaxies, dwarf ellipticals that formed by an initial star burst (\[D1\], asterisks) have larger ages up to $`58`$ Gyr.
Therefore, the simulated colour-magnitude relations (panels a-d) have shallower slope and a larger scatter than the observations. At the reddest edge, we find the relation where massive ellipticals have red colours as observed. Redder colours like V-K show smaller scatter because the origin of the scatter is the younger stars in the simulated galaxies.
#### 3.3.3 Mass-to-Light Ratios
Figure 7 shows (a) baryon fractions, (b) gas fractions, (c) stellar mass-to-light ratios, and (d) total mass-to-light ratios against stellar masses. All of them are measured in spherical regions with the radius of $`2r_\mathrm{e}`$, which includes $`60\%`$ of the total stellar mass. The baryon fractions (panel a) are $`0.35`$ in massive galaxies, and decrease toward $`0.1`$ in dwarf galaxies. In the galaxy, the baryon fraction increases toward the galaxy centre, but is $`0.5`$ at the most in the central $`2`$ kpc. Even at the galaxy centre, equal amounts of dark mass exist with the baryons in the simulated galaxies. The gas fractions (panel b) are less than $`5\%`$ for the galaxies with $`M_{}>10^{10}M_{}`$, and increases towards $`50\%`$ in dwarf galaxies. Therefore, the stellar mass to the total mass ratio well correlate with the galaxy mass. More stars form in more massive systems.
The stellar mass-to-light ratio (panel c) is derived in terms of the SSP model, which should be consistent with the observational estimates; for ellipticals, $`M_{}/L`$ \[$`M_{}/L_{}`$\] $`58`$ and $`69`$ in the V-band and the B-band, respectively. A trend can be found that massive ellipticals have larger $`M_{}/L`$, although the scatter is large. For dwarfs, $`M_{}/L`$ decreases to $`26`$ and $`25`$ in the V-band and the B-band, respectively, which is due to younger ages and lower metallicities than ellipticals. However, in our simulation, the contribution of dark matter is quite large even at the galaxy centre. The total mass-to-light ratios (panel d) are as large as $`2040`$ and $`2070`$ respectively for the simulated ellipticals and dwarfs. Observationally, the mass of dark matter can be estimated with the X-ray hot gas; the dark matter mass can be several times larger than the stellar mass for several galaxies even at the galaxy centre, which results in the total mass-to-light ratios of $`20`$ (Matsushita et al. (1998)).
#### 3.3.4 Fundamental Plane
Figure 8 shows the V-band fundamental plane shown in the $`\kappa `$-space (Bender, Burstein & Faber (1992)). The parameters $`\kappa _1`$, $`\kappa _2`$, and $`\kappa _3`$ express masses, surface brightnesses, and mass-to-light ratios, respectively, and are defined as $`\kappa _1(2\mathrm{log}\sigma _0+\mathrm{log}r_\mathrm{e})/\sqrt{2}`$, $`\kappa _2(2\mathrm{log}\sigma _0+2\mathrm{log}I_\mathrm{e}\mathrm{log}r_\mathrm{e})/\sqrt{6}`$, and $`\kappa _3(2\mathrm{log}\sigma _0\mathrm{log}I_\mathrm{e}\mathrm{log}r_\mathrm{e})/\sqrt{3}`$, where $`\sigma _0`$ is the central velocity dispersion and $`I_\mathrm{e}10^{0.4(\mathrm{SB}_\mathrm{e}27)}`$. The solid line shows the observed relation for the V-band (Pahre (1999)), and we reproduce the observed relations from the B-band to the near infrared.
The $`\kappa _1`$-$`\kappa _2`$ diagram (lower panel) is the face-on view of the fundamental plane. There is no correlation between masses and surface brightnesses, and the simulated galaxies cover the similar region to the observed giant and dwarf ellipticals (gray points, Pahre (1999)). Dwarf galaxies populate the region with small masses and faint surface brightnesses.
The $`\kappa _1`$-$`\kappa _3`$ diagram (upper panel) is the edge-on view. The simulated galaxies follow the observed relation with a shallow slope (solid line, $`\kappa _3=0.171\kappa _1+0.143`$, Pahre (1999)), where more massive ellipticals have large โmass-to-light ratiosโ. Since the baryon fraction is as small as $`0.5`$ even in the galaxy centre, the โmass-to-light ratiosโ expressed by $`\kappa _3`$ is affacted by the dark matter content. The rms fitting of the simulated ellipticals (dotted line), $`\kappa _3=0.247\kappa _10.035`$, is in agreement with the observation, although the slope is slightly steeper and the zero-point is larger by 0.1 dex than observed. Dwarf galaxies lie above the relation, where the total mass-to-light ratios are larger than for giant ellipticals (Fig.7d).
An intrinsic scatter exists along the fundamental plane. The origin of the scatter is clearly shown by the symbols; merger galaxies (open circles and squares) have smaller $`\kappa _2`$ and larger $`\kappa _3`$ than non-merger ellipticals (filled circles and squares). Figure 9 shows the histograms of the deviations from the fundamental plane for the non-major merger (\[E1\]-\[E3\], gray area) and major merger galaxies (\[E4\]-\[E5\], hatched area). These distributions are different, and major merger galaxies have larger $`\kappa _3`$ than non-major merger galaxies at given $`\kappa _1`$. The thick dashed line is for the simulated dwarf galaxies, which have much larger $`\kappa _3`$. Therefore, the origin of the scatter along the fundamental plane is found to lie in differences in merging history. As discussed for Fig. 5, the galaxies that undergo major mergers tend to have larger $`r_\mathrm{e}`$ and fainter $`I_\mathrm{e}`$. There is no significant change in $`\sigma _0`$ and total luminosity $`L`$. From the definitions, these result in smaller $`\kappa _2`$ ($`L^2\sigma _0^2r_\mathrm{e}^5`$) and larger $`\kappa _3`$ ($`L^1\sigma _0^2r_\mathrm{e}`$).
What is the origin of the fundamental plane? In our simulation, the slope of the $`\kappa _1`$-$`\kappa _3`$ relation is originated from the combination of the metallicity, age, and the dark matter content. i) Metallicities: The slope mainly stems from the metallicity effect; metallicities are higher for massive galaxies (Fig.6ef), which results in smaller total luminosities and larger stellar mass-to-light ratios (Fig.7c). The scatter of the mass-metallicity relation is not small, and the scatter of the FP can be reduced by the other effects. ii) Ages: No relation is found between age and mass (Fig.6ij), and galaxies are as old as $`10`$ Gyr coevally at $`3.5<\kappa _1<4`$. At $`\kappa _1<3.5`$, however, some low-mass metal-rich galaxies are younger and have $`\kappa _3`$ as small as the other metal-poor galaxies (e.g., the galaxies at $`\kappa _1=3.25`$ and $`3.13`$ on the dotted line). iii) Dark matter content: Since the baryon fraction is larger for massive galaxies (Fig.7a), the total mass-to-light ratio is smaller for massive galaxies, which is the opposite of the FP. However, the old and metal-rich galaxies at $`\kappa _13.43.5`$ have larger baryon fraction and less dark matter content, which results in smaller $`\kappa _3`$ than the massive galaxies at $`\kappa _14`$.
## 4 DISCUSSION
The hydrodynamical simulation including star formation involves an uncertainty that is how to determine the star formation parameter $`c`$. We show here that it can be constrained from the scaling relation of galaxies. Because the star formation timescale controls when and where stars form in a contracting gas cloud in a dark matter halo, the star formation timescale determines not only the ages of stars, but also the size of the galaxy. To reproduce the observed radius-magnitude relation of elliptical galaxies, $`c`$ is constrained to be $`0.1`$ in our model, i.e., the local star formation timescale is ten times longer than the dynamical timescale. The global star formation timescale is found to be $`12`$ Gyr, which is longer than the $`0.1`$ Gyr that is commonly adopted in one-zone models (e.g., Kodama & Arimoto (1997)). For spiral galaxies, the global star formation timescale should be several Gyr to meet various observations such as the metallicity distribution function of the Milky Way Galaxy. Physical processes that change the star formation timescale might be the existence of rotation, the suppression of star formation due to the UV background radiation, and environmental effects in clusters.
We should note that mass accretion and star formation are truncated artificially by the initial conditions in our simulations. However, in observed ellipticals, star formation should be truncated at $`z2`$ by some process, and an analogous truncation is required, as discussed in ยง3.1. In our model, supernova feedback is not enough to stop the star formation, even if we increase the feedback energy with different IMF (ยง3.2). Hypernovae, which eject ten times larger energy than normal supernovae, may increase the supernova feedback. Tidal stripping and ram-pressure stripping should affect, but it may be hard to explain the uniformity of cluster and field ellipticals. The AGN feedback can suppress the star formation effectively. If the relation between the black hole mass and the bulge luminosity (Magorrian et al. (1998)) suggests that the AGN activity increases following the merging of galaxies, the AGN feedback may become effective around this redshift.
To explain the lack of gas in present-day ellipticals, and to explain the heavy elements in the intracluster medium (Ciotti et al. (1991)), a galactic wind seems indispensable. In our simulation, however, galactic winds do not occur in large galaxies, and thus star formation never terminates completely. This causes the large scatter of 0.3 dex in the B-V colour-magnitude relation (Fig. 6). This problem arises from the SPH method and the feedback scheme. Including kinetic feedback $`f_{\mathrm{kin}}>0`$ as in other simulations (Kawata & Gibson (2003)) does not seem to be a good solution. With $`f_{\mathrm{kin}}=0.1`$, the surface brightness decreases at the centre of our ellipticals, and metal-rich gases blow out forming new stars at large radii. This results in effective radii which are too large and metallicity gradients which are too shallow to reproduce the observations. Changing the IMF slope is not good solution either in our model. $`x=1.10`$ gives stronger feedback, and the ejected wind mass does increase. Unfortunately, the ejected metal fraction does not increase, and the overall stellar metallicity becomes too high. The AGN feedback can not help solving problem because metals should be ejected during the star formation at very early epoch. It takes time to generate super massive black hole by the merging of galaxies. Because of the metallicity dependence of gas cooling, enriched gas is easily turn to stars.
Because of the existence of two distinct types of supernova explosion that produce different elements on different timescales, the abundance ratios of the stellar population can be used to put constraints on star formation histories. SNe II, which are the core collapse-induced explosions of short-lived massive stars ($`>\mathrm{\hspace{0.17em}8}M_{}`$), produce more O and Mg relative to Fe (i.e., \[O/Fe\] $`>0`$) with a timescale of $`10^{68}`$ yr, while SNe Ia, which are the thermonuclear explosions of accreting white dwarfs in close binaries, produce mostly Fe and little O with a timescale of $`0.520`$ Gyr. (The yields relative to solar value are almost the same between O and Mg both for SNe II and Ia.) Observationally, it has often been claimed that Mg is overabundant in elliptical galaxies (Worthey, Faber & Gonzalez (1992); Thomas & Maraston (2003)). Moreover, Fisher, Franx & Illingworth (1995) showed that ellipticals with larger $`\sigma _0`$ tend to have larger \[Mg/Fe\]<sub>0</sub>. However, in our simulation, luminosity weighted \[O/Fe\] spans $`0.3`$ to $`0.1`$ without any correlation with mass, although dwarf galaxies have larger \[O/Fe\] by 0.1 dex. (In model A with stronger feedback and the flatter IMF with $`x=1.10`$, \[O/Fe\] spans $`0.1`$ to $`0.1`$.) Star formation needs to be terminated by efficient feedback in the simulations. We should note, however, that it may still be difficult to explain why \[O/Fe\] is larger in more massive galaxies. The timescale and duration of star formation should be longer for massive galaxies because the deep dynamical potential keeps gas cooling. Some possibilities have been suggested (e.g., Matteucci (1994)); i) the slope of the initial mass function may be different in massive ellipticals, ii) the nucleosynthesis yields of SNe II may be different, namely, iron production may be different because of lower energies and/or larger fall-back, iii) the binary frequency may be smaller and less SNe Ia may occur in massive ellipticals, and iv) the metal enriched wind may cause selective mass loss, so that iron enriched gas can be ejected efficiently before it is consumed in forming the next generation of stars.
## 5 CONCLUSIONS
We study the formation and evolution of galaxies with a GRAPE-SPH chemodynamical model that includes various physical processes associated with the formation of stellar systems; radiative cooling, star formation, feedback from SNe II, SNe Ia, and SWs, and chemical enrichment. We simulate 74 slowly-rotating spherical fields with CDM initial fluctuations (spin parameter $`\lambda 0.02`$), and obtain 128 galaxies with stellar masses in the range $`10^{912}M_{}`$ (74 ellipticals, 45 dwarfs, and 9 cD galaxies). In our scenario, galaxies form through the successive merging of subgalaxies. The merging histories are various with differences seeded in the initial conditions. In some cases, galaxies form through the assembly of gas rich small galaxies, and the process looks like a monolithic collapse. In other cases, the final galaxies form through a major merger of preexisting galaxies. Major mergers are defined as those with mass ratio $`f>0.2`$ occurring at $`z<3`$.
Internal structure such as metallicity gradients is greatly affected by merging histories, while the global properties are determined from overall masses according to the scaling relations. Assuming that the star formation timescale is ten times longer than the local dynamical timescale (i.e., $`c=0.1`$), we succeed in reproducing the observed global scaling relations, e.g., the Faber-Jackson relation, the Kormendy relation, the colour-magnitude relation, the mass-metallicity relation and the fundamental plane. The different relations for ellipticals and dwarfs could be reproduced, although simulated dwarfs have larger effective radii than observed because of the lack of resolution. The luminosity-weighted ages of dwarfs span in wide range, $`38`$ Gyr, depending on their star formation histories, while ellipticals are as old as $`710`$ Gyr independent of their mass.
Adopting the Salpeter IMF ($`x=1.35`$), we could reproduce the mass-metallicity relations both for the central stellar metallicity and for the mean stellar metallicity of the whole of the galaxy. However, the slope is shallower and the scatter is larger than observed, which are because the feedback is not so effective that most metals are locked into stars in the simulation. The colour-magnitude relation also shows a larger scatter because the star formation does not terminated completely in the simulations.
An intrinsic scatter exists along the fundamental plane, and the origin of the scatter in the simulation lies in differences in merging history. Galaxies that undergo major mergers tend to have larger effective radii and fainter surface brightnesses, which result in larger $`\kappa _1`$ (expressing masses), smaller $`\kappa _2`$ (surface brightnesses), and larger $`\kappa _3`$ (mass-to-light ratios).
We examine the dependence of our results on the star formation parameter $`c`$ and the slope $`x`$ of the initial mass functions. Although the time evolution of the size and total mass do not change at all, the stellar mass depends on these parameters in the sense that stronger feedback (i.e., larger $`c`$ and smaller $`x`$) decreases the stellar mass. We found that $`c`$ controls when and where stars form in the contracting gas cloud, thus changing the baryon fraction and determining the effective radius at given mass. With the model to reproduce the observed mass-radius relation, the baryon fractions are $`0.3`$ and $`0.1`$, and the total mass-to-light ratios are $`2040`$ and $`2070M_{}/L_{}`$ respectively for the simulated ellipticals and dwarfs. The biggest effect of $`x`$ is on the metallicity, but it also changes the gas fraction and the residual gas mass in the galaxy.
On the other hand, the wind gas mass, which is defined as the gas particles that have been inside $`r<r_{200}`$ and are outside $`r>2r_{200}`$ at present, does not depend on these parameters, and $`25\%`$ of the accreted baryons can be blown away. The wind efficiency, the ratio of the wind gas mass to the total stellar mass, depends on $`c`$, and is $`0.4`$ and $`0.35`$ for $`c=1.0`$ and $`c=0.1`$, respectively. However, most heavy elements end up locked into stars in the galaxy. The ejected metal fraction depends only on the star formation timescale, and is $`2\%`$ even if we take the quickest star formation rate (i.e., $`c=1.0`$) under our scheme. To explain the metals detected in the intracluster medium in galaxy clusters, the feedback scheme should be improved so as that the enriched gas can blow away efficiently. Changing the IMF do not help in solving this problem under our star formation and feedback schemes.
## Acknowledgments
This paper is a part of the Ph.D. thesis of C. Kobayashi in the Astronomy Department of the University of Tokyo. I would like to thank the supervisor, K. Nomoto, N. Arimoto, and S.D.M. White for detailed suggestions. I am grateful to N. Nakasato, J. Makino, T. Kodama, V. Springel, F. van den Bosch, and A. Renzini for fruitful discussions. I also thank to the Japan Society for Promotion of Science for a financial support, and to the National Observatory of Japan for the GRAPE system. |
warning/0506/cond-mat0506604.html | ar5iv | text | # Theory of Anisotropic Hopping Transport due to Spiral Correlations in the Spin-Glass Phase of Underdoped Cuprates
## I Introduction
One of the central issues in the physics of the high-temperature superconductors is the nature of the ground state at low doping. In particular, the possible co-existence of ordering tendencies in the spin and charge sectors at low temperature is currently being actively investigated. Orenstein ; Sachdev ; Kivelson1 Experimentally at low temperature La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> (LSCO) co-doped with Nd (LNSCO) exhibits static lattice deformation Tranquada as well as incommensurate (IC) magnetic order at doping $`x1/8`$. The lattice deformation indicates the presence of static charge order (stripes). Recently similar behavior has also been found in La<sub>2-x</sub>Ba<sub>x</sub>CuO<sub>4</sub> (LBCO), LBCO $`x1/8`$. However for $`x<1/8`$ and in Nd-free LSCO, charge order has not been detected, while IC magnetism persists down to $`x0.02`$. Wakimoto ; Fujita Thus while IC magnetic order is generically observed in the underdoped regime, static charge order seems to be confined to the neighborhood of $`x=1/8`$.
Theoretically the IC magnetism in the cuprates is often modeled as originating from static charge stripes, and since these are not universally present, one is forced to introduce additional concepts, such as โfluctuatingโ stripes, electronic liquid crystals, etc. Kivelson1 ; Kivelson2 The microscopic origin of such states is still not clear and is currently being debated. An alternative explanation for the IC magnetism which follows naturally from the $`tJ`$ model, is the formation of a non-collinear spiral state, partially relieving the frustration due to the hole motion. SS Such a state would not co-exist with charge order, as can be shown in the context of effective Landau theory. Zachar However it is still possible that in the presence of anisotropic Dzyaloshinski-Moriya interactions the spin spiral could cause weak lattice modulation with period generically half that of the spiral. While the spiral ground state has a tendency to be unstable toward phase separation, we have shown recently that in the extended $`tt^{}t^{\prime \prime }J`$ model the spiral can be stabilized by the presence of the additional hoppings. SK The spiral description was successfully applied to explain magnetic properties of LSCO, such as the location of the elastic neutron scattering peaks and the change of the incommensurability direction by 45 across the superconductor-insulator boundary ($`x=0.055`$). SK1 The spiral was also proposed as a candidate for the ground state in the spin-glass (insulating) phase for $`0.02<x<0.055`$. HCS ; Juricic ; SK1
In the present work we address transport properties within the spiral framework. Our main motivation comes from the recent experimental data in the spin-glass phase of LSCO showing transport anisotropies as large as 50% to 70%, both in DC and AC measurements. Ando ; Dumm It has been argued that these data provide indirect support to the notion of fluctuating stripes or electronic liquid crystals, but no quantitative theory exists that takes these concepts into account. Kivelson1 Recent infrared experiments that carefully identify phonon modes in LSCO once again find no charge ordering tendencies in the spin-glass phase of this material. Padilla Therefore we take the point of view that the ground state has a spiral spin structure, and calculate the transport anisotropy in the variable-range hopping (VRH) regime for low temperature and frequency, where the anisotropy has the largest value. We show that the spatial anisotropy of the hole wave-function in a spiral translates into anisotropy of the hopping transport. Our main result is that for microscopic parameters appropriate for LSCO (within the $`tt^{}t^{\prime \prime }J`$ model), the magnitude of the anisotropy is large (40%-90%, depending on temperature), and close to the one found experimentally. Thus we demonstrate that the transport anisotropy data can be explained quantitatively within the spiral theory which does not involve any tendency of the holes to self-organize into charge stripes. The anisotropy was also analyzed in Ref. Juricic, on the basis of topological defect scattering in a spiral; however the results are applicable to the quasi-metallic (higher temperature) regime where the anisotropy is very small, of the order of several percent. We emphasize that in the present work we consider the low-temperature, strongly-localized regime where it is clear from experiment that the anisotropy is the largest (10-20 times larger than the value in the metallic regime). Finally we mention that the problem was also addressed within the spin-charge separation scenario which provides a description of the pseudogap phase Marchetti and could lead to effective โinsulatingโ behavior; however within this approach the transport anisotropy is linked to the magnetic correlation length anisotropy which can be taken from experiment but cannot be explicitly calculated.
### I.1 Summary of Previous Results and Notation
Our starting point is the description of the spin-glass phase ($`0.02<x<0.055`$) developed in Ref. SK1, which is briefly summarized below. First, a single hole resides near the points $`๐ค_0=(\pm \pi /2,\pm \pi /2)`$, with a quadratic dispersion around them $`ฯต_๐ค\frac{\beta _1}{2}k_1^2+\frac{\beta _2}{2}k_2^2`$, where $`k_1`$ is perpendicular to the face of the magnetic Brillouin zone and $`k_2`$ is parallel to the face. Within the effective $`tt^{}t^{\prime \prime }J`$ model for LSCO, the parameters are taken to be: $`t/J=3.1,t^{}/J0.5,t^{\prime \prime }/J0.3`$, where $`J125\text{meV}`$. From now on we measure energies in units of $`J`$ (i.e. set $`J=1`$) and lengths in units of the lattice spacing $`a`$ (we set $`a=1`$). By using the self-consistent Born approximation one finds that these parameters lead to an almost isotropic dispersion $`\beta _1\beta _2=\beta 2.2`$, and a quasiparticle residue around the nodal points $`Z0.34`$. A hole trapped by the Coulomb field of the Sr ion generates a spiral distortion of the Nรฉel background, which can be parametrized as: $`|\text{i}=e^{i\theta (๐ซ_๐ข)๐ฆ\sigma /2}|,|\text{j}=e^{i\theta (๐ซ_๐ฃ)๐ฆ\sigma /2}|`$, $`\text{ i }\text{ โupโ sublattice},\text{ j }\text{ โdownโ sublattice}`$, where $`๐ฆ`$ is an arbitrary unit vector perpendicular to the quantization axis of the states $`|`$, $`|`$. The angle $`\theta (๐ซ)`$, measuring deviations from collinearity is given by: SK1
$$\theta (๐ซ)=\frac{Zt}{\sqrt{2}\pi \rho _s}\frac{(๐_\pm ๐ซ)}{r^2}\left[1e^{2\kappa r}(1+2\kappa r)\right],$$
(1)
where $`\rho _s0.18`$ is the spin stiffness, and $`1/\kappa 35`$ is the localization length of the orbital part of the wave function. The value of $`1/\kappa `$ is extracted from experimental data, SK1 and is generally expected to increase with doping. The unit vector $`๐_\pm =\frac{1}{\sqrt{2}}(1,\pm 1)`$, in the usual square-lattice coordinate system. At finite doping ($`0.02<x<0.055`$) the interaction of the long-range dipole distortions from holes trapped at different Sr ions leads to spiral magnetic order, SK1 characterized by average $`\overline{\theta }๐_\pm ๐ซ`$, as shown in Fig. 1(a). In a perfect square lattice the spiral can be directed along any diagonal $`(1,\pm 1)`$. However from the location of the elastic neutron scattering peaks Wakimoto the incommensurability is determined to be along the orthorhombic $`\widehat{b}`$ direction. In our picture this means that the orthorhombic deformation pins the direction of the spiral, as shown in Fig. 1, and thus we set $`๐=๐_+=\frac{1}{\sqrt{2}}(1,1)`$. In what follows the exact nature of the pinning mechanism, to be discussed elsewhere, will not be important.
Experimentally the magnetic correlation length is finite $`\xi 620`$ (decreases with increasing doping). Wakimoto ; Fujita On the theoretical side it was argued that $`\xi `$ is finite due to topological defects that lead to frustration of the long-range spiral order. HCS ; SK1 Since the localization length is less than the magnetic scale, $`1/\kappa <\xi `$, we expect that the effective description of transport properties, developed below in terms of the one-hole wave function, is quantitatively valid also at finite (small) doping (i.e. the topological frustration mechanism does not affect our considerations). It should be noted that the above inequality while being explicit is also the strongest (most restrictive) condition we could give. Its refinement would have to come out of a detailed theory of IC magnetism in the spin-glass phase (i.e. a self-consistent theory that takes into account both the effective disorder and the spiral formation, generated by the doped holes). Purely theoretical arguments aside, transport measurements in the doping range 0%-4% suggest that the system is in a strongly-localized regime at low temperatures, and we thus expect our calculations to be valid as long as doping is not too close to the insulator-metal boundary (at 5.5%).
The rest of the paper is organized as follows. In Section II we analyze the properties of a localized hole in a spiral background, and in Section III we use the holeโs wave function to calculate the in-plane transport anisotropies in the variable-range hopping regime. Section IV contains our conclusions.
## II Localized hole in the presence of magnetic spiral correlations
The coupling $`H_{SP}`$ between the spin of the magnetic background (angle $`\theta (๐ซ)`$) and the orbital wave-function of the hole $`\chi (๐ซ)`$ generates the spiral and has the form: SK1
$$H_{SP}=\sqrt{2}Zt(๐\theta )\chi ^2(๐ซ)d^2r.$$
(2)
The effective Schrรถdinger equation is then:
$$\left(\frac{\beta }{2}^2\frac{q^2}{r}\sqrt{2}Zt(๐\theta )\right)\chi (๐ซ)=ฯต\chi (๐ซ).$$
(3)
Here the Coulomb potential of the $`Sr`$ ion which keeps the hole localized is $`\frac{q^2}{r}=\frac{q_0^2}{_er}`$, where $`q_0`$ is the unit charge and $`_e`$ is the effective dielectric constant known to be quite large in the copper oxides, $`_e30100`$ (increases with doping). Kastner
The last term in (3) determines the anisotropy of $`\chi (๐ซ)`$ and is given by
$`\sqrt{2}Zt(๐\theta )={\displaystyle \frac{Z^2t^2}{\pi \rho _s}}{\displaystyle \frac{1}{r^2}}[f(\kappa r)g(\kappa r)\mathrm{cos}(2\phi )],`$
$`g(z)=1(1+2z+2z^2)e^{2z},f(z)=2z^2e^{2z}.`$ (4)
The coordinate system is chosen so that $`๐`$ is parallel to the $`\widehat{b}`$ axis and $`\phi `$ is the polar angle (Fig. 1). The value of $`\kappa `$ (and consequently $`ฯต`$) can be found by a variational minimization of the total energy: SK1
$$\kappa =\frac{2q^2/\beta }{1\frac{\mathrm{\Lambda }}{2}},ฯต=\beta \kappa ^2/2,\mathrm{\Lambda }\frac{Z^2t^2}{\pi \beta \rho _s}.$$
(5)
We will use $`\kappa 0.30.4`$, consistent with experiment, SK1 rather than rely on (5) which contains the uncertainty related to the exact value of the dielectric constant. The dimensionless parameter $`\mathrm{\Lambda }`$, as defined in (5), characterizes the coupling between the orbital motion of the hole and the deformation of the spin background. Substituting the values of $`Z,\beta `$ appropriate for LSCO (discussed in Section I.A), we obtain $`\mathrm{\Lambda }1`$, which will be used from now on.
Let us look for solution of (3) in terms of the series in angular harmonics
$$\chi (\kappa r,\phi )=\chi _0(\kappa r)+\chi _2(\kappa r)\mathrm{cos}(2\phi )+\chi _4(\kappa r)\mathrm{cos}(4\phi )+\mathrm{}$$
(6)
Then Eq. (3) leads to a set of coupled equations for $`\chi _i,i=0,2,4,..`$. We have solved these equations numerically and the results for the first two harmonics are presented in Fig. 2. The next harmonics are found to be small, for example $`\chi _4/\chi _0<2\times 10^2`$, and can be in fact safely neglected. At large distances all the wave-functions behave as $`\chi _ie^{\kappa r}`$, while at finite distances this behavior is substantially modified, as shown in Fig. 2.
## III Transport anisotropies in the variable-range hopping (VRH) regime
### III.1 In-plane DC Resistivity Anisotropy
Below characteristic temperature $`T_{VRH}2030\text{K}`$, the low-temperature resistivity of LSCO in the spin-glass phase can be described by the 2D version of the Mott VRH formula $`\rho \mathrm{exp}(T_0/T)^{1/3}`$. Keimer ; Kastner ; Lai ; Suzuki Here $`T_0`$ depends somewhat on doping and sample quality and generally decreases when doping increases (and thus conduction becomes easier). As an estimate, for example at 4% doping the data of Ref. Keimer, are well fit with $`T_0500\text{K}`$; Lai analyzing the curves of Ref. Suzuki, we obtain $`T_0200300\text{K}`$, and from Ref. Ando, we have extracted $`T_0200\text{K}`$.
Within the spiral framework the physical idea behind the resistivity anisotropy is quite simple. According to (6) the hole wave function acquires an elliptic deformation induced by the spin spiral, as shown schematically in Fig. 1(b). Hence hopping in the $`\widehat{b}`$-direction is less probable than hopping in the $`\widehat{a}`$-direction. Using (6), and essentially following the derivation of Mottโs result, Mott we obtain for the resistivity anisotropy:
$`{\displaystyle \frac{\rho _b(T)}{\rho _a(T)}}={\displaystyle \frac{_0^{2\pi }\mathrm{sin}^2(\phi )\mathrm{\Phi }^2(\kappa r_T,\phi )๐\phi }{_0^{2\pi }\mathrm{cos}^2(\phi )\mathrm{\Phi }^2(\kappa r_T,\phi )๐\phi }},`$
$`\mathrm{\Phi }(\kappa r,\phi ){\displaystyle \frac{\chi (\kappa r,\phi )}{\chi _0(\kappa r)}}.`$
Here $`r_T=\frac{1}{2\kappa }(T_0/T)^{1/3}`$ is the VRH length, Mott and in order for the approach to be justified we must certainly have $`\kappa r_T>1`$. Eq. (III.1) reflects the difference in the wave-function overlap for the two electric field directions. The factors $`\mathrm{sin}^2(\phi ),\mathrm{cos}^2(\phi )`$ arise from the fact that the conductivity varies with the square of the carrier jump distance projection in the electric field direction. Mott Since the holeโs dispersion is isotropic ($`\beta _1=\beta _2=\beta `$), the exponential part of the wave function is isotropic and consequently the resistivity anisotropy is expected to arise from the anisotropy of the exponential prefactors, whose overlap leads to (III.1). Both the VRH length and $`T_0`$ thus remain isotropic within this formulation.
It is clear from (III.1) that as $`T`$ decreases (i.e. $`\kappa r_T`$ increases) the anisotropy grows, due to the increase of $`|\chi _2/\chi _0|`$ (Fig. 2(Inset)) and consequently the more pronounced angular dependence of $`\mathrm{\Phi }(\kappa r_T,\phi )`$. The results are summarized in Fig. 3, where also the evolution of the zero temperature anisotropy $`\rho _b(T=0)/\rho _a(T=0)`$ as a function of $`\mathrm{\Lambda }`$ is shown for completeness in the inset. At very low temperature ($`T1\text{K}`$, Lai i.e. $`T/T_010^3`$) we would have to take into account a crossover to the Coulomb gap regime, which however would practically not influence the curve in Fig. 3.
The data of Ref. Ando, were taken at temperatures $`T>10\text{K}`$, meaning that the lowest ratio $`T/T_00.05`$ (we take $`T_0200\text{K}`$). As the temperature increases beyond $`T_{VRH}2030\text{K}`$ when $`\kappa r_T1`$, the approach based on VRH conduction ceases to be valid, as the conduction mechanism changes to impurity band conduction and eventually quasi-metallic behavior. In the most relevant low-temperature range below $`T_{VRH}`$ (corresponding to largest anisotropy), both the calculated magnitude of $`\rho _b/\rho _a`$ and its temperature dependence are very close to the experimental results, Ando although in these experiments the temperature is not low enough to penetrate the โdeepโ VRH regime $`T/T_00.05`$ where the anisotropy should increase even further. It should be also noted that there exists quite a bit of uncertainty in the determination of $`T_0`$ and hence in the determination of the exact value of $`\rho _b/\rho _a`$ from Fig. 3.
### III.2 AC Resistivity Anisotropy
At finite frequency and temperature the calculation of the AC conductivity is a very complicated problem. However in the โquantumโ limit when the frequency $`\omega T`$, the VRH AC conduction is expected to be dominated by resonant absorption by singly occupied pairs (without involvement of phonons), and is usually relevant in doped semiconductors. Efros The VRH length in the AC regime (neglecting for a moment the angular dependence of the states) is logarithmic: Efros $`r_\omega =(1/\kappa )\mathrm{ln}(2|ฯต|/\omega )`$, which is the main difference from the DC case. From now on we denote $`\mathrm{\Omega }\omega /(2|ฯต|)`$. The formula for $`r_\omega `$ follows from the fact that upon evaluation of the conductivity, the most effective pairs are the ones that satisfy: Efros $`\omega =2I(r)`$, where $`I(r)=I_0e^{\kappa r}`$ is the overlap integral, and a typical estimate of the prefactor is $`I_0|ฯต|`$. The functional dependence of the conductivity in 2D at low frequency $`\mathrm{ln}(1/\mathrm{\Omega })1`$, in the leading logarithmic approximation, is given by the Mott-Shklovskii-Efros expression: $`\sigma (\omega )\omega r_\omega ^3[\omega +q^2/r_\omega ]`$, where the second term takes into account the Coulomb interaction in the resonant pair. Efros
The generalization of the Mott-Shklovskii-Efros formula to the anisotropic case is straightforward, as it amounts to taking into account the non-exponential (angle-dependent) part of the wave-function, leading to the replacement (with logarithmic accuracy): $`\mathrm{ln}(1/\mathrm{\Omega })\mathrm{ln}\left(\mathrm{\Phi }(\kappa r_\omega ,\phi )/\mathrm{\Omega }\right)L_{\omega ,\phi }`$, where $`\mathrm{\Phi }`$ is defined in (III.1). Taking also into account the expressions for $`ฯต,\kappa `$ (5), we obtain for the resistivity anisotropy:
$`{\displaystyle \frac{\rho _b(\omega )}{\rho _a(\omega )}}={\displaystyle \frac{_0^{2\pi }\mathrm{sin}^2(\phi )(\omega ,\phi )๐\phi }{_0^{2\pi }\mathrm{cos}^2(\phi )(\omega ,\phi )๐\phi }},`$ (8)
$`(\omega ,\phi )=L_{\omega ,\phi }^2\left(1+{\displaystyle \frac{2\mathrm{\Omega }}{1\mathrm{\Lambda }/2}}L_{\omega ,\phi }\right).`$
The expression (8) is plotted in Fig. 4. Due to the logarithmic dependence, in the theoretical limit of zero frequency the anisotropy would vanish (albeit very slowly), i.e. $`\rho _b(\omega =0)/\rho _a(\omega =0)=1`$. However one must keep in mind that the above expressions are not valid at arbitrary low frequencies, as we must have $`\omega >T`$. Generally we find that the (maximum) AC anisotropy of 30-40% is somewhat smaller than the DC anisotropy.
When attempting to compare our results to the experimental data of Ref. Dumm, we realize that $`\omega `$ and $`T`$ are not sufficiently low, nor is the difference between them sufficiently high to justify the separation into โquantumโ and โthermalโ VRH conduction and consequently the Mott-Shklovskii-Efros approach. The temperature used is $`T=13\text{K}`$ and at the lowest frequency of $`\omega =20\text{cm}^12.5\text{meV}`$, estimating for the hole energy $`|ฯต|=\beta \kappa ^2/2510\text{meV}`$, we have $`\mathrm{\Omega }0.130.25`$. Moreover, upon increasing the frequency to $`80100\text{cm}^1`$, a change of behavior from insulating to conducting takes place, a broad peak develops, Dumm and consequently it is not clear to us that the AC data are ever in a clean VRH regime. In spite of all this the magnitude of the calculated anisotropy agrees reasonably well with experiment. It should be experimentally possible to lower the temperature (towards the mK range) as well as $`\omega `$ in order to probe the VRH AC anisotropy more reliably and compare with our theory.
## IV Discussion and Conclusions
We would like to reiterate that in our theory there is no tendency of the holes to form charge stripe-like structures. Nevertheless the resistivity shows anisotropic behavior (and so does the uniform magnetic susceptibility, to be discussed separately). It is particularly important, in our view, that the calculations presented in this work produce quantitative results, and we would thus hope that the successes of the present formulation would stimulate further exploration of the spiral and similar scenarios, not involving any charge-ordering tendencies. The present work is relevant to the spin-glass regime of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, $`0.02<x<0.055`$. However we have also argued SK ; SK1 that the spin spiral structure could exist in the superconducting state $`x>0.055`$. The theory naturally explains why the incommensurate direction, determined from elastic neutron scattering, rotates by 45 exactly at the insulator-superconductor transition point. On the other hand it is usually argued that at $`x=1/8`$ both in charge-ordered LNSCO Tranquada and LBCO, Hucker non-collinear spiral order is not consistent with experiment. We would indeed not claim that the spiral ground state is stable at that particular doping since in fact the spiral becomes commensurate with the lattice and therefore significant changes in the ground state could occur due to the spin-lattice coupling.
In conclusion, starting from a spiral ground state which unambiguously follows from the extended $`tJ`$ model, we have calculated the in-plane anisotropy of the low-temperature DC and low-frequency AC resistivity of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> in the spin-glass phase. The theory has no fitting parameters and the calculations are performed for the variable-range hopping regime. Within the spiral description the transport anisotropy is due to the anisotropy of the hole wave-function since the hole hops in a non-collinear (spiral) spin background. The AC anisotropy reasonably agrees with experiment in spite of the fact that the data are on the border of the VRH regime. The experimental data for the DC resistivity are well within the VRH regime and here both the calculated magnitude and temperature dependence of the anisotropy agree with experiment very well.
###### Acknowledgements.
We are grateful to L. Benfatto, J. Haase, and D. Poilblanc for valuable discussions and comments. V.N.K. acknowledges the support of the Swiss National Fund. |
warning/0506/astro-ph0506478.html | ar5iv | text | # Evidence for Evolution or Bias in Host Extinctions of Type 1a Supernovae at High Redshift
## 1 Introduction
Type 1a supernovas are candidates for standard astrophysical candles, from which the relation of redshift $`z`$ and distance can be estimated. In a universe of constant expansion the โHubble plotโ made from magnitudes and redshifts should be a straight line. Data is now available for a wide range of redshifts up to 1.755 (Schmidt et al. 1998; Garnavich et al. 1998; Perlmutter et al. 1998; Riess et al. 1998; Perlmutter et al. 1999; Knop et al. 2003; Tonry et al. 2003; Barris et al 2004). The Hubble diagrams derived from supernovae have indicated an upward bending curve, interpreted as acceleration of the expansion rate, along with even more complicated features of โjerkโ. It is important to explore other interpretations, including possible evolution of supernova or host galaxy characteristics with redshift. Many papers have explored non-cosmological explanations (Coil et al. 2000; Leibundgut 2001; Sullivan et al. 2003; Riess 2004). Meanwhile, the high redshift host galaxies have significantly different morphologies compared to those at low redshifts (Abraham & van den Bergh 2001; Brinchmann et al. 1998; van den Bergh 2001). Dust and related extinction characteristics may certainly depend on redshift (Totani & Kobayashi 1999). Furthermore the abundance ratios of the progenitor stars may be different at different redshifts (Hรถflich et al. 2000). Several studies emphasize that evolution effects cannot be ruled out (Falco et al. 1999; Aguirre 1999; Farrah et al. 2004; Clements et al. 2004).
In this paper we find evidence for evolution or bias in the extinction parameters used to pre-process the data. If the effect is due to bias, extinctions have been overestimated, which makes supernovas appear more dim. Yet just the same phenomenon could occur from a real physical effect in which the actual host extinctions are correlated with the deviation of magnitudes from model fits.
### 1.1 Background
Traditional Hubble diagrams represent the relation of observed flux $``$ to the luminosity of the source $``$,
$`={\displaystyle \frac{}{4\pi d_L^2}},`$ (1)
where $`d_L`$ is the so-called luminosity distance. The distance modulus $`\mu _p=mM`$, where $`m`$ and $`M`$ are the apparent and absolute magnitudes respectively, is
$$\mu _p=5\mathrm{log}d_L+25,$$
(2)
where the luminosity distance $`d_L`$ is in megaparsecs.
The process of converting observed data into the supernova magnitudes reported actually contains an additive parameter, called the extinction coefficient $`A`$. Extinction may depend on frequency, designated by $`A_B`$, $`A_R`$, etc. The units of $`A`$ are magnitude. In practice $`A`$ shifts the supernova magnitude $`m_0`$ deduced from light-curves to a reported magnitude (โextinction corrected magnitudeโ) $`m=m_0A`$. Our galaxy contributes extinction, as do the additional extinction effects associated with supernova host galaxies, which are more model dependent.
Riess et al (2004) discovered 16 Type Ia supernovas at high redshifts and compiled a 157 source โgoldโ data set held to be of the highest reliability. Extinctions are listed in Riess et al (2004) for all except 24 sources among this โgoldโ set.
## 2 Analysis
Riess et al focus on the differences of magnitudes $`\mathrm{\Delta }\mu `$ relative to the traditional Hubble plot. In Fig. 1 we show the residuals $`\mathrm{\Delta }\mu `$ versus the extinction coefficients $`A_V`$, for all the sources for which extinctions are known. There is a clear correlation. The sense of correlation is that points with $`\mathrm{\Delta }\mu >0`$, lying above the straight line Hubble plot, tend to have small or even negative extinction, and points lying below the straight line tend to have large extinction. A precedent for examining correlations of residuals is given in Williams et al., (2003).
Residuals depend on the baseline model from which they are measured. Fig.1 uses the FRW model and โconcordanceโ parameters $`\mathrm{\Omega }_M=0.27`$, with $`\mathrm{\Omega }_\mathrm{\Lambda }=0.73`$ under the constraint $`\mathrm{\Omega }_k=0`$. This is one of the baselines cited by Riess et al (2004). Here $`\mathrm{\Omega }_M`$ is the matter density, $`\mathrm{\Omega }_\mathrm{\Lambda }`$ the vacuum energy density and $`\mathrm{\Omega }_k=1\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$. The class of FRW models predicts the luminosity distance as
$$d_L=\frac{c(1+z)}{H_0|\mathrm{\Omega }_k|^{1/2}}\mathrm{sinn}\left\{|\mathrm{\Omega }_k|^{1/2}_0^z๐z\left[(1+z)^2(1+\mathrm{\Omega }_Mz)z(2+z)\mathrm{\Omega }_\mathrm{\Lambda }\right]^{1/2}\right\}$$
(3)
Here $`\mathrm{sinn}`$ denotes $`\mathrm{sinh}`$ for $`\mathrm{\Omega }_k>0`$, $`\mathrm{sin}`$ for $`\mathrm{\Omega }_k<0`$ and is equal to unity for $`\mathrm{\Omega }_k=0`$. Parameters are fit by minimizing $`\chi ^2`$, defined by
$$\chi ^2=\underset{i}{}\frac{(\mu _p^i\mu _0^i\mu _{p0})^2}{(\delta \mu _0^i)^2},$$
(4)
where $`\mu _p^i`$ and $`\mu _0^i`$ are the theoretical and observed distance moduli respectively and $`\delta \mu _0^i`$ are the reported errors. Our notation includes the intercept parameter $`\mu _{p0}`$ (not always explicit in the literature). The Hubble constant $`H_0`$ and fit parameters such as the zero point are not reported in Riess et al (2004), which states that they are irrelevant and arbitrarily set for the sample presented here. We verify (Riess et al, 2004) $`\chi ^2=178`$ for the concordance parameters cited above, along with the other $`\chi ^2`$ values for several other studies, presented below.
### 2.1 Quantification
We quantify the correlation of extinctions with residuals with the correlation coefficient $`R(\mathrm{\Delta }\mu ,A_V)`$, also simply $`R`$, defined by
$`R(\mathrm{\Delta }\mu ,A_V)={\displaystyle \frac{\underset{i}{}(\mathrm{\Delta }\mu _i\overline{\mathrm{\Delta }}\mu )(A_{V,i}\overline{A}_V)}{\sigma _{\mathrm{\Delta }\mu }\sigma _{A_V}}},`$ (5)
where $`\overline{\mathrm{\Delta }}\mu ,\sigma _{\mathrm{\Delta }\mu }`$ are the means and standard deviation of the $`\mathrm{\Delta }\mu `$ set, with corresponding meaning for $`\sigma _{A_V},\overline{A}_V`$. The correlation $`R(\mathrm{\Delta }\mu ,A_V)=0.439`$ for the concordance parameters cited above, excluding the 24 sources for which extinctions are not known. The integrated probability (confidence level, $`P`$-value) to find correlations equal or larger in a random sample is $`4.2\times 10^7`$.
To investigate whether the correlation of extinctions with residuals might be a model artifact, we decided to fit several other models cited by Riess et al (2004). The results of these fits are shown in Table 1. For example, under the best fit model with $`\mathrm{\Omega }_M=0.31,\mathrm{\Omega }_\mathrm{\Lambda }=0.69=1\mathrm{\Omega }_M`$ then $`R(\mathrm{\Delta }\mu ,A_V)=0.434`$ with probability $`P=5.6\times 10^7`$.
From Fig. 1 we see that the correlation is strongest for large values of $`A_V`$. For example, for the best fit parameters $`(\mathrm{\Omega }_M=0.31,\mathrm{\Omega }_\mathrm{\Lambda }=0.69)`$ we find that excluding the four sources with $`A_V>0.8`$ the correlation coefficient goes down to $`R(\mathrm{\Delta }\mu ,A_V)=0.28`$ with $`P=1.5\times 10^3`$. Retaining the 139 points with $`A_V0.5`$ yields $`R(\mathrm{\Delta }\mu ,A_V)=0.18`$. We do not have a particular reason to entertain these cuts except to make the correlation go away. At the risk of complicating interpretation, one can try dividing the residuals by the data pointโs uncertainty. This is an uncertain trial because a fundamental issue is the uncertainty in the extinction coefficients, which is unavailable from the literature. Fig. 2 shows the correlation with error bars assigned to the residuals.<sup>1</sup><sup>1</sup>1We thank an anonymous referee for this suggestion. The figure shows that most of the data with $`A_V>0.3`$ lies below 0, indicating bias. Division by uncertainty only reduces $`R(\mathrm{\Delta }\mu /\sigma ,A_V)0.37`$ for the gold set, an effect of having introduced noise.
We next examine whether the correlation seen in the residuals depends on redshift. We divide the data as equally as possible in a large redshift sample ($`z0.41`$, 78 sources) and a low redshift sample ($`z<0.41`$, 79 sources). (The cut $`z0.46`$ was identified by the Hubble team as a transition region.) For the low redshift sample we find $`R(\mathrm{\Delta }\mu ,A_V)=0.509`$, $`P=1.2\times 10^5`$, compared to the high redshift sample yielding $`R(\mathrm{\Delta }\mu ,A_V)=0.378`$, $`P=3.7\times 10^3`$. Although statistics have been diluted, it is clear that the two samples show different behavior, with the correlation being much more significant in the low redshift sample.
Questions then branch along three lines: (1) The assignment of extinctions by present schemes may contain hidden bias. (2) There may be a real physical effect at work, and (3) Systematic errors might be re-evaluated in order to ameliorate the significance of the correlation.
$``$ 1: A seldom discussed but established bias exists in the assignment of $`A_V`$ from the fits to light curves. We find it highlighted by the Berkeley group (Perlmutter et al 1999, especially the Appendix). The scheme used starts with a conditional probability $`P(A|A_{dat})`$, where $`A_{dat}`$ is the extinction from the best fit to the light curve data. A prior probability $`P_0(A_{dat})`$ is assumed, and from Bayesโ Theorem the probability of $`A`$ after seeing the data is estimated. The value of $`A`$ is chosen to โmaximize the probability of $`A`$โ given the combined information from the prior and the data.
The method introduces an extra dependence on the choice of priors. For prior distributions centered at small host extinction, the work of Hatano (1998) is cited, based on Monte Carlo estimates from host galaxies of random orientation. Freedom is used to formulate a one-sided prior distribution with support limited to $`A>0`$. This make a bias in the combination of assuming $`A>0`$ for the priors (fluctuations could do otherwise) and the detailed way in which $`A_{dat}`$ is assigned. This bias tends to cause the same signal as dimming or acceleration (Perlmutter et al (1999)). As of 1999 the outcomes of this bias were stated to be less than 0.13 magnitude.
Yet one would need an absolute standard to evaluate any bias reliably. Subsequently the method itself has evolved (Riess et al. 2004), citing an iterative โtraining procedureโ we have not found described in detail. A few points now have $`A_V<0`$.
There is evidently a further bias in taking data from the peak of the proposed distribution. It is not the same thing as sampling the proposed distribution randomly. Iteration of a procedure taking from the peak tends to drive a Bayesian update procedure towards a narrow distribution centered at the peak. In some renditions this may cause systematic errors of fluctuations to evolve towards becoming underestimated.
$``$ 2: It is possible that the extinction correlation is a signal of physical processes of evolution with redshift. It is impossible to adequately summarize the literature discussing this possibility. Aguirre (1999) made a comparatively early study with a balanced conclusion that extinction models might cause some of the effects interpreted as acceleration. Drell, Loredo & Wasserman (2000) concentrate on this question, concluding that the methodology of using type 1a supernovas as standard candles cannot discriminate between evolution and acceleration. Farrah et al. (2004) (see also Clements et al. 2004) cite a history of work scaling optical frequency extinction with the sub-millimeter wavelength observations (Hildebrand 1983; Casey 1991; Bianchi 1999). They report extinction for 17 galaxies with $`z=0.5`$ with sub-millimeter wavelengths. While stating consistency with local extinctions at the $`1.3\sigma `$ level, they add โIt does however highlight the need for caution in general in using supernovae as probes of the expanding Universe, as our derived mean extinction, $`A_V=0.5\pm 0.17`$, implies a rise that is at face value comparable to the dimming ascribed to dark energy. Therefore, our result emphasizes the need to accurately monitor the extinction towards distant supernovae if they are to be used in measuring the cosmological parameters.โ The trend of Farrahโs observation is same as the correlation seen in the supernova data, and remarkably, the corrections we obtain empirically in various fits (below) almost all amount to 0.5 magnitude or less. The fact that low redshift objects show higher correlation implies that there is a higher tendency to overestimate extinctions of these sources in comparison to the sources at higher redshifts. Since the estimated extinctions show no correlation with redshift, this suggests that the true low redshift extinctions, on the average, may be smaller in comparison to the extinctions of high redshift sources. Nevertheless the question of evolution of the sources remains open and will not be resolved here.
$``$ 3: Perhaps the means of assigning extinction coefficients are reasonable on average, but statistical fluctuations have given a false signal. Then the error bars on the extinction coefficients come to be re-examined. Inasmuch as this is coupled to the entire chain of data reduction, it is beyond the scope of this paper.
### 2.2 Empirically Corrected Extinctions
Without engaging in physical hypotheses of extinction, it is reasonable to test whether a different extinction model can give a satisfactory fit to the data. We studied a corrected value $`A_V(\delta )`$ depending on the parameter $`\delta `$ by the simple rule
$$A_V(\delta )=(1+\delta )A_V.$$
(6)
We then determine $`\delta `$ by the best fit to the cosmological model. The best fit $`\delta `$-values and the corresponding $`\chi ^2`$ values for different models are given in Table 2. Parameter $`\delta `$ produces a huge effect of more than 23 units of $`\chi ^2`$.
There are many ways to compare the new and old fits. As a rule, the model with $`\chi ^2`$ per degree of freedom ($`dof`$, the number of data points minus the number of parameters) closest to unity is favored. Since the new fits decrease $`\chi ^2`$ by 20-some units with one additional parameter, the significance of revising the extinction values is unlikely to be fortuitous. For example the model with $`\mathrm{\Omega }_M=0.27`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.73`$ gives $`\chi ^2/dof=1.14`$ and $`1.01`$ without correction ($`\delta =0`$) and with correction ($`\delta =0.42`$). As a broad rule in comparing data sets, the difference $`\mathrm{\Delta }\chi ^2`$ should be distributed by $`\chi _\nu ^2`$, where $`\nu =1`$ is the number of parameters added. The naive $`p`$-value or confidence level to find $`\mathrm{\Delta }\chi ^2=23`$ in $`\chi _1^2`$ is 1.6 $`\times 10^6`$. Thus introducing $`\delta `$ would be well-justified simply to improve the poor fit of $`\chi ^2178`$ without ever seeing the extinction correlation with residuals. Values of $`\delta `$ for all models are found to be negative, suggesting that the host extinction values given in Riess et al (2004) are overestimates.
It is interesting and significant that the new residuals, computed relative to the revised fits, show negligible correlation with host extinction. This is seen in Fig. 3, which shows the $`R`$ values on the same plot as $`\chi ^2/dof`$. The fact that $`R`$ vanishes when $`\delta `$ meets the best-fit value is significant. It is far from trivial, as $`R`$ concerns an independent set of numbers, the $`A_V`$ values, not directly used in calculating $`\chi ^2`$.
Figure 4 shows the residuals versus corrected host extinction after including the correction term. The reduction in correlation $`R`$ comes with an increased scatter in $`A_V(\delta )`$ at large $`A_V(\delta )`$, which is not unexpected.
It is also interesting to ask whether host extinction might have some dependence on the luminosity distance $`d_L`$. It is hard to imagine no evolution at all, and we explored a linear ansatz. The linear model is
$$A_V(\delta ,d_L)=(1+\delta )A_V+\delta _1d_L.$$
(7)
We add that when a model of evolution is introduced, the cosmological interpretation might be disturbed, so that the outcomes must be taken in context. More cannot be anticipated because the fits themselves will choose $`\delta _1`$. Fit parameters and $`\chi ^2`$ values are given in Table 3.
#### 2.2.1 Is Acceleration Supported?
Accelerating models show no need for the $`\delta _1`$ term. Assuming acceleration, the fits (Table 3) show that reducing extinction values by about 40% explains the data better, and removes an alarming correlation. On the other hand the matter-dominated model $`(\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0)`$ shows interesting sensitivity to $`\delta _1`$. In Fig.5 we compare the sensitivity of different fits to parameter $`\mathrm{\Omega }_M`$. With $`\delta _1=0`$ constrained, the effects of $`\delta `$ are rather orthogonal to those of $`\mathrm{\Omega }_M`$, so that the region $`\mathrm{\Omega }_M0.3`$ is favored whether or not there is a significant correlation $`R`$. Yet varying $`\delta _1`$ greatly broadens acceptable values of $`\mathrm{\Omega }_M`$, while maintaining the $`R0`$ effect of $`\delta `$. The significance depends on oneโs hypothesis: if one chooses $`\mathrm{\Omega }_M=1`$ a-priori, parameter $`\delta _1`$ is traded for parameter $`\mathrm{\Omega }_M`$. The overall probability of either hypothesis is only in part determined by the $`p`$-value of the data given the distribution: the rest depends on oneโs prior beliefs in evolution, which we will not pursue. It is fair to say that the revised fits give more leeway to matter-dominated models on statistical grounds.
In all cases fits are driven to $`A_VA_V(\delta 0.4)0.6A_V`$, either simply to improve $`\chi ^2/dof`$, or to remove the correlation with residuals.
To conclude, analysis using reported extinction coefficients is well known to produce good fits to acceleration of the expansion rate. However the extinctions show correlation with residuals with random chance probability using two independent tests, the extinction correlation and $`\chi ^2`$ values, both below the level of $`10^6`$. The hypothesis that extinction coefficients should be corrected empirically provides substantially improved fits to the data, while also eliminating significant correlation of residuals. A model of linear evolution yields interesting effects of high statistical significance correlated with redshift. The studies indicate either bias in host extinction assignments or evolution of the source galaxies. The significance of acceleration itself cannot be resolved on the basis of these studies, but might be revised, depending on oneโs priors. We suggest that observers report uncertainties in their assignment of extinction parameters, both in the future and for the existing data sets.
Acknowledgments: Research supported in part under DOE Grant Number DE-FG02-04ER14308. This work was completed when PJ was visiting the National Center for Radio Astrophysics, Pune. He thanks Prof. V. Kulkarni for kind hospitality. JP thanks Hume Feldman and Ruth Daly for discussions.
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warning/0506/cond-mat0506529.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Gauge theories and their associated concepts have played important roles not only in elementary particle physics but also in condensed matter physics. For example, the conventional superconducting phase transition is characterized as a change of the gauge dynamics of the electromagnetic U(1) gauge symmetry. Also, for a variety of strongly-correlated electron systems, it has been recognized that their phase structures and the properties of low-energy excitations are naturally described by using the terminology of gauge theories. As such low-energy โquasi-particlesโ, the composite fermions/bosons in the fractional quantum Hall states and the holons and spinons in the t-J model of high-$`T_c`$ cuprates have been proposed. It was argued that their unconventional properties like fractionality and their existence itself may be explained by a confinement-deconfinement phenomenon of the gauge dynamics of the effective gauge theories derived from the original models. This interesting idea is still controversial, but certainly warrants further investigation.
Most of the problems in gauge-theoretical studies on the strongly-correlated electron systems reduces to studying the phase structures of gauge theories coupled to gapless relativistic/nonrelativistic matter fields. In the elementary particle physics, it is generally believed that the phase structures of gauge systems coupled to matter fields are not easy to study analytically because introduction of matter fields results in the lack of simple order parameters such as Wilson loops for pure gauge systems. Furthermore, inclusion of gapless fermions make numerical simulations a difficult task since one must face quite nonlocal interactions generated by integrating over fermion variables. Sometimes couplings to gapless matter fields change the universality class of the gauge system under consideration from that of the pure gauge model. A good example is the four-dimensional (4D) QCD coupled with light quarks in which the number of light quarks strongly influences its phase structure.
In this paper, we address this problem of gauge dynamics of coupled systems. Specifically, we are interested in the U(1) lattice gauge theory (LGT) with gapless/gapful and relativistic/nonrelativistic matter fields in three dimensions (two spatial dimensions at zero temperature). This theory of course covers the important model QED<sub>3</sub>. It appears also as a main part of the effective gauge theory of the strongly-correlated electron systems mentioned above, and so plays an important role in studying these systems. In the ordinary 3D compact U(1) pure gauge system (i.e., without matter fields) with local interactions, it is established that only the confinement phase is realized because of instanton condensation. For the case with additional massless (gapless) matter fields coupled to the U(1) gauge field, recent studies give controversial results on the possibility of a deconfinement phase; it is supported in Ref. whereas it is denied in Ref..
Our approach to this problem is by (i) introducing an effective theory of the original theory and (ii) studying its phase structure numerically. As the effective theory we use a 3D U(1) pure LGT with nonlocal interactions among gauge variables. These nonlocal interactions are along the temporal direction and mimic the effect of matter fields. We consider exponentially-decaying interactions for massive matter fields (i.e., fields with gaps) and power-decaying interactions for massless (gapless) fields. We shall see that certain cases of power-decaying interactions exhibit second-order phase transitions which separate the confinement phase and the deconfinement phase. The existence of the deconfinement phase in the effective theory indicates that it is realized in the original model if the number of gapless matter fields is sufficiently large. This result is in agreement with the results of Ref..
The rest of the paper is organized as follows. In section 2, we briefly survey the effect of matter fields in several aspects. In section 3, the original model and its effective nonlocal gauge model are explained. Section 4 is devoted for numerical calculations. We calculate the internal energy, specific heat, expectation values of Polyakov loops, Wilson loops, and density of instantons. All these quantities indicate a second-order confinement-deconfinement phase transition (CDPT) for the gauge models with sufficiently long-range correlations among the gauge variables. In section 5, we study tractable low-dimensional spin models, which are obtained by simple reduction of the gauge degrees of freedom in the nonlocal effective model. Then we obtain an intuitive picture of the CDPT of the present long-range U(1) gauge theories. Section 6 is devoted for conclusion. In Appendix, the effective nonlocal model is studied by the low- and high-temperature expansions. The analytic expressions of the internal energy and the specific heat are in good agreement with the numerical calculations in Sect.4.
## 2 Effects of Matter Fields on Gauge Dynamics
In this section, we briefly review the effects of matter fields upon the U(1) gauge dynamics in several aspects.
### 2.1 Weak-coupling regime in noncompact/compact U(1) gauge theory
One may generally expect that inclusion of massless matter fields to a system of gauge field may drastically change the gauge dynamics at long wavelengths. For the case that relativistic and massless matter fields are coupled to a U(1) gauge field $`A_\mu (x)(\mu =0,1,2)`$ in 3D, the one-loop radiative correction of matter fields at weak gauge couplings generates the following nonlocal term in the effective action of $`A_\mu (x)`$;
$$\mathrm{\Delta }Ae^2d^3xd^3y\underset{\mu ,\nu }{}F_{\mu \nu }(x)\frac{1}{|xy|^2}F_{\mu \nu }(y),F_{\mu \nu }=_\mu A_\nu _\nu A_\mu .$$
(2.1)
It is obvious that the above term strongly suppresses fluctuations of $`A_\mu (x)`$ at long distances due to the factor $`|xy|^2`$. The above nonlocal terms are leading at low energies and momenta if its effective coupling constant does not vanish at the infrared limit. Then the potential energy $`V(r)`$ between two charges separated by distance $`r`$ changes drastically,
$`V(r)\mathrm{log}rV(r){\displaystyle \frac{1}{r}}.`$ (2.2)
In the compact U(1) gauge theory, topologically nontrivial excitations (e.g. instantons) of $`A_\mu (x)`$ appear whose effect at small gauge coupling $`e`$ can be estimated by replacing the field strength $`F_{\mu \nu }(x)`$ in Eq.(2.1) by $`F_{\mu \nu }(x)2\pi n_{\mu \nu }(x)`$, where $`n_{\mu \nu }`$ is an integer field whose rotation measures the instanton number(density) $`\rho (x)=ฯต_{\mu \nu \lambda }_\mu n_{\nu \lambda }(x)`$. In the gauge model with the usual Maxwell term, the potential energy between a pair of instantons at distance $`r`$ is Coulombic, $`V_{\mathrm{ins}}(r)1/r`$. However, the long-range action (2.1) modifies $`V_{\mathrm{ins}}(r)`$ to a long-range one,
$`V_{\mathrm{ins}}(r){\displaystyle \frac{1}{r}}V_{\mathrm{ins}}(r)\mathrm{log}r.`$ (2.3)
Recently, it was argued that a gas of charged particles with the long-range interaction like $`\mathrm{log}r`$ exhibits a phase transition between a dilute gas of dipoles and a plasma. This result implies that instantons in the long-range gauge theories form dipole pairs and do not condense in the weak-coupling region. Since the confinement phase of gauge dynamics requires a condensation of isolated instantons (a plasma phase), this result leads to a deconfinement phase of gauge dynamics there.
### 2.2 The 3D CP<sup>N-1</sup> model on the criticality
As an example of the nonperturbative effect of matter fields, let us consider the CP<sup>N-1</sup> model in a 3D continuum. The action of the model is given as
$$A_{\mathrm{CP}}=d^3x\left[\frac{1}{f}|(_\mu +iA_\mu )z|^2+\sigma \left(|z|^21\right)\right],$$
(2.4)
where $`z_a(x)(a=1,2,\mathrm{},N)`$ is the CP<sup>N-1</sup> field satisfying the local constraint, $`_{a=1}^N|z_a(x)|^2=1`$ for each $`x`$ via the Lagrange multiplier field $`\sigma (x)`$, and $`A_\mu (x)`$ is the auxiliary U(1) gauge field.
At large $`N`$, the $`1/N`$ expansion is reliable and predicts a second-order phase transition at the critical coupling $`f=f_c`$ . For $`f>f_c`$, the model is in the disordered-confinement phase in which $`\sigma (x)0`$, whereas for $`f<f_c`$, the model is in the ordered-Higgs phase in which $`\sigma (x)=0`$ and the CP<sup>N-1</sup> field has a nonvanishing expectation value like $`z_N(x)=v_00`$. As a result, the gauge field $`A_\mu `$ acquires a finite mass ($`v_0`$) by the Anderson-Higgs mechanism, and the low-energy excitations are gapless $`z_a(x)(aN)`$ fields.
Recently, considerable interests have been paid on the question how the gauge field behaves just at the critical point $`f=f_c`$. In the leading order of the $`1/N`$ expansion, it is shown that all the components $`z_a(x)(a=1,2,\mathrm{},N)`$ are massless and the gauge field $`A_\mu `$ acquires the nonlocal โkinetic termโ of Eq.(2.1). This implies that, at the critical point, the nonperturbative fluctuations of $`A_\mu `$ like instantons are suppressed, so the gauge dynamics at $`f=f_c`$ is in the Coulomb phase with the potential $`V(r)1/r`$ as Eq.(2.2).
Then we have numerically studied the 3D CP$`{}_{}{}^{1}+`$ U(1) LGT from the above point of view. The numerical results of the CP<sup>1</sup> model show a second-order transition and the confinement phase is realized at the critical point. However, calculations of CP<sup>N-1</sup> model with $`N=3,4,5`$ show that the topologically nontrivial configurations are suppressed more as $`N`$ increases. We expect that there exists a critical value of $`N=N_c`$, and the deconfinement phase is realized at the critical point for $`N_c<N`$ in compatible with the large-$`N`$ analysis. We note that a similar phenomenon of inducing a deconfinement phase by a plenty of massless matter fields has been established in lattice QCD in 4D. For a sufficiently large number of light flavors $`N_f>7`$, the model stays in the deconfinement phase even the pure gauge term is missing.
### 2.3 Nonrelativistic fermions in strongly-correlated electron systems
Nonrelativistic fermions are distinguished from relativistic fermions by the properties; (i) they propagate only in the positive direction of the imaginary time in path-integral formulation, and (ii) they form a Fermi surface (line). In strongly-correlated electron systems, one faces nonrelativistic fermions not only in the original models but also in their effective gauge models. In studying the fractionalization phenomena of electrons like the charge-spin separation (CSS) in high-temperature superconductivity and the particle-flux separation (PFS) in fractional quantum Hall systems , we regard an electron $`C_x`$ (we suppress the spin index for simplicity) at the site $`x`$ of a 2D spatial lattice as a composite of a fermion $`A_x`$ and a boson $`B_x`$ as
$`C_x`$ $`=`$ $`A_x^{}B_x.`$ (2.5)
To assure the correct physical space composed of $`C_x`$, one imposes the local constraint for the physical states $`|\mathrm{phys}`$ as
$`\left(A_x^{}A_x+B_x^{}B_x1\right)|\mathrm{phys}=0\mathrm{for}\mathrm{each}x.`$ (2.6)
The hopping term of electrons (e.g., the $`t`$-term of the t-J model) may be rewritten via a โdecouplingโ as
$`C_{x+i}^{}C_x+\text{H.c}`$ $`=`$ $`B_{x+i}^{}A_{x+i}A_x^{}B_x+\text{H.c.}`$ (2.7)
$``$ $`(B_{x+i}^{}W_{xi}B_x+A_{x+i}^{}W_{xi}^{}A_x+\text{H.c.})|W_{xi}|^2+\mathrm{}.`$
$`W_{xi}`$ is an auxiliary complex field defined on the link $`(x,x+\widehat{i})`$ where $`i=1,2`$ is the direction index of the lattice. Its phase degree of freedom $`U_{xi}W_{xi}/|W_{xi}|`$ behaves as a spatial component of a compact U(1) gauge field $`U_{xi}[=\mathrm{exp}(i\theta _{xi})]`$. $`U_{xi}`$ represents the binding force of the constituents $`A_x`$ and $`B_x`$. In the path-integral formalism, the partition function $`\mathrm{Tr}_C\mathrm{exp}[\beta H(C)](\beta T^1)`$ has the following representation:
$`Z`$ $`=`$ $`{\displaystyle \underset{x,\tau }{}d\overline{A}_x(\tau )dA_x(\tau )dB_x(\tau )\underset{\mu =0,1,2}{}dU_{x\mu }(\tau )\mathrm{exp}(A)},`$
$`A`$ $`=`$ $`{\displaystyle _0^\beta }d\tau {\displaystyle \underset{x}{}}[i\theta _{x0}\overline{A}_x(_0i\theta _{x0}+\mu _A)A_x\overline{B}_x(_0i\theta _{x0}+\mu _B)B_x`$ (2.8)
$`+t{\displaystyle \underset{i}{}}(\overline{A}_{x+\widehat{i}}U_{xi}(\tau )B_x+\mathrm{H}.\mathrm{c}.)+A_{\mathrm{int}}],`$
where $`\tau ([0,\beta ])`$ is the continuum imaginary time, $`A_x(\tau )`$ is a Grassmann number expressing fermions, $`\mu _{A(B)}`$ is the chemical potential, and $`t`$ is the hopping amplitude. The field $`\theta _{x0}`$ is a Lagrange multiplier field to enforce the constraint (2.6). It may be viewed as the time-component of gauge field. In fact, in the discrete-time formulation on the 3D lattice it appears as the exponent of gauge variable in the $`\tau `$-direction, $`U_{x0}=\mathrm{exp}(i\theta _{x0})`$.
When the gauge dynamics of $`U_{xi}`$ is realized in the confinement phase, the constituents $`A_x`$ and $`B_x`$ are bound within electrons and the relevant low-energy quasi-particles are described by the electron operators $`C_x`$ themselves. On the other hand, if the gauge dynamics is realized in the deconfinement phase, the binding force is weak and these constituents are no more bound in each electron but dissociate each other. This is just the the fractionalization(separation) phenomena of electrons.
In Ref. we have studied the system of Eq.(2.8) at finite temperatures ($`T`$) by the hopping expansion in the temporal gauge $`\theta _{x0}=0`$. After integrating over $`A_x`$ and $`B_x`$ in powers of $`t`$, one obtains the effective interactions. Up to $`O(t^2)`$, by restoring the temporal components $`U_{x0}`$, one gets the following effective interaction among $`U_{x\mu }`$,
$`\mathrm{\Delta }A`$ $``$ $`\delta (1\delta )t^2{\displaystyle \underset{x,i}{}}{\displaystyle _0^\beta }๐\tau {\displaystyle _0^\beta }๐\tau ^{}V_{x_{},i}(\tau ,\tau ^{})+\mathrm{c}.\mathrm{c}.,`$
$`V_{x_{},i}(\tau ,\tau ^{})`$ $``$ $`\overline{U}_{x_{},\tau ^{},i}\mathrm{exp}\left(i{\displaystyle _\tau ^\tau ^{}}๐\tau ^{^{\prime \prime }}[\theta _{x_{}+\widehat{i},0}(\tau ^{^{\prime \prime }})\theta _{x_{},0}(\tau ^{^{\prime \prime }})]\right)U_{x_{},\tau ,i},`$ (2.9)
where $`\delta =A_x^{}A_x`$ is the concentration of fermions. The corresponding terms are illustrated in Fig.1.
The interactions among $`U_{x\mu }(\tau )`$ in $`\mathrm{\Delta }A`$ are quite nonlocal in the $`\tau `$-direction, and $`\mathrm{\Delta }A`$ favors the ordered configurations of $`U_{x\mu }`$, so the deconfinement phase. In the previous papers, we argued that this is the essence of the mechanism of fractionalization phenomena like CSS and PFS. By mapping these gauge models approximately to a spin model, we concluded that the U(1) gauge dynamics is realized in the deconfinement phase at the low-$`T`$ region below certain critical line $`T_c(\delta )`$. Thus the possible deconfienment phase in the model similar to Eq.(2.9) is to support the electron fractionalization phenomena.
### 2.4 Chern-Simons term by Dirac fermions
On considering the effects of matter fields upon 3D gauge dynamics, there is an important difference between fermionic and bosonic matte fields. That is, relativistic fermions have a possibility to generate the Chern-Simons (CS) term $`\mathrm{\Delta }A_{\mathrm{CS}}`$ via radiative corrections:
$`\mathrm{\Delta }A_{\mathrm{CS}}=c{\displaystyle d^3xฯต_{\mu \nu \rho }A_\mu _\nu A_\rho },`$ (2.10)
which violates the parity symmetry. In particular, as the mass term $`m_D\overline{\psi }\psi `$ of 3D two-component spinor Dirac fermion field $`\psi `$ violates the parity invariance, it generates the CS term with the coefficient $`c\mathrm{sgn}(m_D)`$ . On the other hand, scalar fields do not renormalize the coefficient of the Chern-Simons term. In the perturbation theory, the CS term is the leading term at long distances and low energies. So it may change the phase structure of the gauge system, in particular that of the compact gauge systems. One can intuitively expect appearance of a deconfinement phase since the CS term suppreses fluctuations of the gauge field.
However in most of the effective gauge theories of strongly-correlated electron systems, the CS term is not spontaneously generated, because the most of these systems including the t-J model preserves the parity invariance. For example, in the flux state of the Heisenberg antiferromagnetic spin model and the t-J model in the slave fermion representation, relativistic Dirac fermions appear as low-energy excitations, but the parity invariance is preserved in the effective theory because the lattice fermions appear in doublets with opposite signatures of masses, which correspond to four-component spinor Dirac field in the continuum. Then the CS coefficient cancels with each other in the radiative correction as $`c_{m_D=\pm m}\mathrm{sgn}(m_D)=0`$.
Though nonperturbative investigation of the Chern-Simons gauge theory is very important by itself<sup>1</sup><sup>1</sup>1For example, phase structure of SU(N) Maxwell-CS gauge theory has been studied via frustrated Heisenberg spin model without parity invariance, which is a low-energy effective model of strongly-coupled SU(N) gauge theory of fermions . Existence of a deconfinement phase transition is suggested there., we shall focus on parity-invariant lattice gauge theory in the rest of discussions in this paper.
## 3 Nonlocal U(1) Lattice Gauge Theory
In the previous section, we have seen various approaches to study the effect of matter fields upon U(1) gauge dynamics. To confirm these results, one must examine the validity of the approximations employed there. In particular, to check the validity of the hopping expansion at $`T=0`$ is quite important for studies on strongly-correlated electron systems.
Keeping this problem in mind, we start this section with a U(1) LGT coupled with matter fields and introduce its effective nonlocal LGT. Let us consider a 3D cubic lattice (i.e., a 2D lattice with a discrete imaginary time). The gauge field $`U_{x\mu }(\mu =0,1,2)`$ is defined on the link $`(x,x+\widehat{\mu })`$ between the pair of nearest-neighbor sites $`x`$ and $`x+\widehat{\mu }`$. The partition function $`Z`$ is given by the following functional integral,
$`Z`$ $`=`$ $`{\displaystyle \underset{x}{}d\overline{\varphi }_xd\varphi _x\underset{x\mu }{}dU_{x\mu }\mathrm{exp}(A)},`$
$`A`$ $`=`$ $`{\displaystyle \underset{x,y}{}}\overline{\varphi }_x\mathrm{\Gamma }_{xy}(U)\varphi _y+A_U,`$
$`A_U`$ $`=`$ $`q{\displaystyle \underset{x,\mu <\nu }{}}(\overline{U}_{x\nu }\overline{U}_{x+\widehat{\nu },\mu }U_{x+\widehat{\mu },\nu }U_{x\mu }+c.c.),`$ (3.1)
where $`x=(x_0,x_1,x_2)`$ is the site-index of the 3D lattice of the size $`V=N_0N_1N_2`$ with the periodic boundary condition, $`\mu `$ is the imaginary-time index ($`\mu =0`$) and spatial direction indices ($`\mu =1,2`$), $`\varphi _x`$ is the matter field on $`x`$, $`U_{x\mu }=\mathrm{exp}(i\theta _{x\mu })(\pi <\theta _{x\mu }\pi )`$ is the $`U(1)`$ gauge variable on the link $`(x,x+\widehat{\mu })`$, and $`q`$ is inverse gauge coupling constant. $`\mathrm{\Gamma }_{xy}(U)`$ represents the local minimal couplings of $`\varphi _x`$ to $`U_{x\mu }`$. For example, for a bosonic matter field,
$$\underset{x,y}{}\overline{\varphi }_x\mathrm{\Gamma }_{xy}(U)\varphi _y=t\underset{x,\mu }{}[\overline{\varphi }_{x+\widehat{\mu }}U_{x\mu }\varphi _x+\mathrm{H}.\mathrm{c}.]+M^2\underset{x}{}\overline{\varphi }_x\varphi _x,M^2=6+m^2,$$
(3.2)
where $`m^2`$ is the mass in unit of the lattice spacing and $`6`$ in $`M^2`$ is the number of links emanating from each site.
After integrating over the matter field $`\varphi _x`$, effective gauge model is obtained, which includes all contributions from $`\varphi _x`$ to the gauge dynamics,
$`Z`$ $`=`$ $`{\displaystyle \underset{x\mu }{}dU_{x\mu }\mathrm{exp}\left[f\mathrm{Tr}\mathrm{log}\mathrm{\Gamma }_{xy}(U)+A_U\right]},`$ (3.3)
where $`f`$ is a parameter counting the statistics and internal degrees of freedom of $`\varphi _x`$. Due to the $`(\mathrm{Tr}\mathrm{log}\mathrm{\Gamma }_{xy}(U))`$ term, the effective gauge theory becomes nolocal. For relativistic matter fields, a formal expression of the effective gauge theory action is obtained by the hopping expansion and it is expanded as a sum over all the closed random walks $``$ (loops including backtrackings) on the 3D lattice, which represent world lines of particles and antiparticles as
$$\mathrm{Tr}\mathrm{log}\mathrm{\Gamma }_{xy}(U)=\underset{}{}\frac{\gamma ^{L[]}}{L[]}\underset{(x\mu )}{}U_{x\mu }.$$
(3.4)
$`L[]`$ is the length of $``$, and $`\gamma =(6+m^2)^1`$ is the hopping parameter. There are many different random walks that have the same shape of a closed loop on the lattice. Each random walk in such a family may have a different starting point and/or backtrackings. This degeneracy cancels out the denominator $`L[]`$ in Eq.(3.4). For the constant gauge-field configuration $`U_{x\mu }=1`$, the expansion in (3.4) is logarithmically divergent $`\mathrm{log}m`$ as $`m0`$ due to the lowest-energy zero-momentum mode.
Below we shall study a slightly more tractable model than that given by Eq.(3.3). It is suggested by the hopping expansion (3.4), and obtained by retaining only the rectangular loops extending in the $`\tau `$-direction in the loop sum and choosing their expansion coefficients as follows;
$`Z_๐ฏ`$ $`=`$ $`{\displaystyle \underset{x\mu }{}dU_{x\mu }\mathrm{exp}(A_๐ฏ)},`$
$`A_๐ฏ`$ $`=`$ $`g{\displaystyle \underset{x}{}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{\tau =1}{\overset{N_0}{}}}c_\tau (V_{x,i,\tau }+\overline{V}_{x,i,\tau })+A_S,`$
$`V_{x,i,\tau }`$ $`=`$ $`\overline{U}_{x+\tau \widehat{0},i}{\displaystyle \underset{k=0}{\overset{\tau 1}{}}}\left[\overline{U}_{x+k\widehat{0},0}U_{x+\widehat{i}+k\widehat{0},0}\right]U_{xi},`$
$`A_S`$ $`=`$ $`g\lambda {\displaystyle \underset{x}{}}(\overline{U}_{x2}\overline{U}_{x+\widehat{2},1}U_{x+\widehat{1},2}U_{x1}+\mathrm{c}.\mathrm{c}.),`$ (3.5)
where $`g`$ is the (inverse) gauge coupling constant, and $`V_{x,i,\tau }`$ is the product of $`U_{x\mu }`$ along the rectangular $`(x,x+\widehat{i},x+\widehat{i}+\tau \widehat{0},x+\tau \widehat{0})`$ of size $`(1\times \tau )`$ in the $`(i0)`$ plane. See Fig.2.
In $`A_S`$, we have retained only the single-plaquette coupling with the coefficient $`(g\lambda )`$. For the nonlocal coupling constant $`c_\tau `$, we consider the following three cases;
$`c_\tau `$ $`=`$ $`\{\begin{array}{cc}\tau ^\alpha ,\hfill & \mathrm{power}\mathrm{law}\mathrm{decay}(\mathrm{PD}\alpha )(\alpha =1,2,3),\hfill \\ e^{m\tau },\hfill & \mathrm{exponential}\mathrm{decay}(\mathrm{ED}),\hfill \\ 1,\hfill & \mathrm{no}\mathrm{decay}(\mathrm{ND}).\hfill \end{array}`$ (3.9)
The power $`\alpha =1`$ in the PD model in (3.9) reflects the effect of the relativistic massless excitations without dimensional parameters. In fact, this $`c_\tau `$ generates a logarithmically divergent action for $`U_{x\mu }=1`$ explained below Eq.(3.4) as one can see from the relation, $`_\tau \mathrm{exp}(m\tau )\tau ^1\mathrm{log}(1/m)`$. The action for $`m=0`$ is then proportional to $`_\tau \tau ^1\mathrm{log}N_0`$ for finite $`N_0`$. On the other hand, the ED model contains the parameter $`m`$ with mass dimension and simulates the case of massive matter fields. (We used Eq.(3.9) instead of $`\mathrm{exp}(m\tau )/\tau `$ to make the comparison with the PD case more definitive.) The ND model corresponds to gauge model coupled to nonrelativistic fermions with a Fermi surface (or Fermi line) as it is seen from Eq.(2.9).
## 4 Numerical Results
In this section, we report our numerical calculations of the internal energy, the specific heat, the Polyakov lines, etc., and determine the phase structure of the models. Most of our interest concerns a possible CDPT in the 3D compact U(1) gauge models with the nonlocal interactions.
In the MC simulations, we consider the isotropic lattice, $`N_\mu =N(\mu =0,1,2)`$, with the periodic boundary condition up to $`N=32`$, where the limit $`N\mathrm{}`$ corresponds to the system on a 2D spatial lattice at $`T=0`$. For the mass of the ED model, we set $`m=1`$. For the spatial coupling $`\lambda `$ scaled by $`g`$, we consider the two typical cases $`\lambda =0`$ (i.e., no spatial coupling) and $`\lambda =1`$.
### 4.1 Internal energy and specific heat
First we calculate the following โinternal energyโ $`E`$ and the โspecific heatโ $`C`$ per site;
$`Z_๐ฏ`$ $`=`$ $`{\displaystyle [dU]\mathrm{exp}(A_๐ฏ)}\mathrm{exp}(FV),`$
$`E`$ $``$ $`{\displaystyle \frac{1}{V}}A_๐ฏ={\displaystyle \frac{1}{V}}{\displaystyle \frac{g}{Z_๐ฏ}}{\displaystyle \frac{dZ_๐ฏ}{dg}}=g{\displaystyle \frac{dF}{dg}},`$
$`C`$ $``$ $`{\displaystyle \frac{1}{V}}(A_๐ฏA_๐ฏ)^2=g^2{\displaystyle \frac{d^2F}{dg^2}}.`$ (4.1)
Here we note that the definition of $`C`$ (4.1) is the response of $`E`$ under the variation of $`g`$ and is different from the conventional specific heat that measures the response under the variation of temperature itself. The latter contains extra terms associated with the change of $`N_0`$. Because these terms behave less singular than the variance of $`E`$, they are irrelevant in searching for phase transitions.
In Fig.3, we present $`E`$ and $`C`$ for $`\lambda =1`$ vs the nonlocal coupling $`g`$. In Appendix we calculate $`E`$ and $`C`$ by the high-temperature expansion (HTE) for small $`g`$ and by the low-temperature expansion (LTE) for large $`g`$. The MC results of $`E`$ and $`C`$ are consistent with these expansions. For comparison with HTE (e.g., for $`0<g<0.1`$ in the PD-1 model ($`\alpha =1`$)), see Fig.15 in Appendix. For comparison with LTE, the leading result of LTE is shown by the straight lines in Fig.3(a,b) and $`C=1`$ in Fig.3(c,d).
First, let us see the PD-1 model ($`\alpha =1`$) in detail. $`E`$ of Fig.3(a) connects the results of HTE and LTE. Since the LTE is an expansion around $`V_{x,i,\tau }=1`$, this behavior implies $`V_{x,i,\tau }1`$ for large $`g`$. $`C`$ of Fig.3(c) shows that its peak develops as the system size $`N`$ increases. These two points indicate that the PD-1 model exhibits a second-order phase transition separating the disordered (confinement) phase and the ordered (deconfinement) phase at $`g=g_c0.17`$ which is determined from the data of $`N=24`$. The existence and the nature of the phase transition will be confirmed later by the measurement of the Polyakov lines and the instanton density as we show in the following subsections.
We notice that the location of the peak in $`C`$ of Fig.3(c) shifts to smaller $`g`$ as the system size $`N`$ increases in the direction opposite to the usual second-order phase transitions. This behavior reflects the fact that the couplings among $`U_{x\mu }`$ increases effectively as $`N`$ increases due to the additional terms in the summation over $`N_0`$ in the action even if one fixes the overall constant $`g`$. This is the characteristic nature of nonlocal interactions in strong contrast to local interactions.
On the contrary, in the ED model of Fig.3(d), the peak of $`C`$ does not develop as $`N`$ increases, showing no signals of a second-order transition. It may have a higher-order transition or just a crossover. Similar behavior of $`C`$ is observed in the ordinary U(1) gauge systems with local actions which have only the confinement phase. The physical meaning of the above โcrossoverโ in the ED model shall become clear by studying instantons in Sect.4.3.
It is quite interesting to see whether the data of $`C`$ for $`N=8,16`$ and $`24`$ in Fig.3(c) exhibit the finite-size scaling behavior. To this end, let us introduce a parameter
$`ฯต(gg_{\mathrm{}})/g_{\mathrm{}}`$ (4.2)
where $`g_{\mathrm{}}`$ is the critical gauge coupling of the infinite system at $`N\mathrm{}`$. Then let us assume that the correlation length $`\xi `$ scales as $`\xi ฯต^\nu `$ with a critical exponent $`\nu `$. We also expect that the peak of $`C`$ diverges as $`C_{\mathrm{peak}}ฯต^\sigma `$ as $`N\mathrm{}`$ with another critical exponent $`\sigma `$. The finite-size scaling hypothesis predicts that the specific heat $`C(ฯต,N)`$ for sufficiently large $`N`$ scales as
$$C(ฯต,N)=N^{\sigma /\nu }\varphi (N^{1/\nu }ฯต),$$
(4.3)
where $`\varphi (x)`$ is a certain scaling function. In Fig.4, we present $`\varphi (x)`$ determined by using the data in Fig.3(c) with $`\nu =1.2,\sigma /\nu =0.25`$ and $`g_{\mathrm{}}=0.11.`$ This result indicates that the finite-size scaling law holds quite well. Cbonsidering the errors in the data we estimate the values of scaling parameters and $`g_{\mathrm{}}`$ as
$`\nu =1.21.3,\sigma /\nu =0.250.26,g_{\mathrm{}}=0.100.12.`$ (4.4)
The simulations of the PD-1 and ED models with $`\lambda =0`$ give similar behaviors of $`E`$ and $`C`$ as the $`\lambda =1`$ case, preserving the above phase structure for $`\lambda =1`$. That is, the PD-1 model ($`\lambda =0`$) has a CDPT, while the ED model ($`\lambda =0`$) has no transition.
Let us next consider the ND model. In Fig.5, we present $`C`$ for $`\lambda =0`$ and 1, which show strong signals of a second-order phase transition. However, the value of the critical coupling for $`\lambda =1`$, $`g_c0.1(N=8),0.045(N=16),0.03(N=24)`$, decreases very rapidly as the system size increases. This behavior is explained by the increase of effective coupling explained above. Then one may think that $`g_c0`$ as $`N\mathrm{}`$, i.e., only the deconfinement phase survives in the ND model. This expectation is consistent with the fact that the coefficient $`Q_2_\tau c_\tau ^2`$ of HTE, $`C=2(2Q_2+\lambda ^2)g^2+O(g^3)`$ in Eq.(A.13) diverges as $`Q_2N\mathrm{}`$ for the ND model, i.e., the radius of convergence in the HTE is zero. On the contrary, $`Q_2`$ is finite for the PD-1 model, assuring us $`g_c0`$. This is supported also by the scaling analysis given above.
From these results for the various cases of long-range interaction, it seems that there exists a critical power $`\alpha =\alpha _c`$ below which the CDPT takes place. In Fig.6 we show $`C`$ of the PD model with $`\alpha =2`$ and $`\alpha =3`$.
We obtain a very interesting result, i.e., for the PD-2 system ($`\alpha =2`$) with the nonvanishing spatial coupling $`\lambda =1`$, there exists a second-order CDPT as in the PD-1 case, whereas in the PD-3 case ($`\alpha =3`$) the peak in $`C`$ does not develop as $`N`$ increases, hence no signals of CDPT. Furthermore, careful study of the PD-2 case shows that the second-order CDPT disappears for vanishing spatial coupling $`\lambda =0`$. Thus we conclude that the critical power is $`\alpha _c=2`$ and the spatial coupling $`\lambda `$ controls the existence of the CDPT. The above results will be confirmed by the study of the Polyakov line in the following subsection.
### 4.2 Polyakov lines
We have argued the possible CDPT by measuring the thermodynamic quantities like $`E`$ and $`C`$. In order to study the CDPT in more details and further the nature of gauge dynamics in each phase, it is useful to calculate order parameters in gauge theory like Polyakov lines and Wilson loops.
First, let us introduce the Polyakov lines $`P_x_{}`$ for each spatial site $`x_{}(x_1,x_2)`$ and study their spatial correlations $`f_P(x_{})`$,
$`P_x_{}={\displaystyle \underset{x_0=1}{\overset{N_0}{}}}U_{x_{},x_0,0},x_{}=(x_1,x_2),`$
$`f_P(x_{})=\overline{P}_x_{}P_0.`$ (4.5)
Since the present model (3.5) contains no long-range interactions in the spatial directions, $`f_P(x_{})`$ is expected to supply a good order parameter to detect a possible CDPT. In the deconfinement phase, the fluctuations of $`U_{x0}`$ are small, which implies an order in $`f_P(x_{})`$, i.e., we expect $`f_P(x_{})0`$ as $`x_{}`$ large in the deconfinement phase.
In Fig.7, we present $`f_P(x_{})`$ for the PD-1 and ED models. The PD-1 model of Fig.7(a) clearly exhibits an off-diagonal long-range order, i.e., $`lim_x_{}\mathrm{}f_P(x_{})0`$ for $`g0.20`$, whereas the ED model of Fig.7(b) does not for all $`g`$โs. To see this explicitly, we plot in Figs.7(c) and (d) the order parameter $`p(f_P(x_{}^{\mathrm{MAX}}))^{1/2}`$ for the PD-1 model, where $`x_{}^{\mathrm{MAX}}N/\sqrt{2}`$ is the distance at which $`f_P`$ becomes minimum due to the periodic boundary condition. $`p`$ of the PD-1 model($`\lambda =0`$) starts to develop continuously from zero at $`g=g_c0.15`$. The size dependence of $`p`$ shows a typical behavior of a second-order phase transition. Thus the gauge dynamics of the PD-1 model is realized in the deconfinement phase for $`g>g_c`$, whereas it is in the confinement phase for $`g<g_c`$. In contrast, the ED model stays always in the confinement phase. These results including the value of $`g_c`$ are in good agreement with those derived from the data of $`E`$ and $`C`$ given in Fig.3.
Let us turn to the PD-2 model with and without the spatial coupling. In Fig.8, we show the result of $`p`$ for $`\lambda =1`$ and $`\lambda =0`$. We observe that the model with $`\lambda =1`$ shows a typical behavior of the second-order phase transition as the system size is increased, whereas the case of $`\lambda =0`$ does not. From this result and the observation of $`C`$ in the previous subsection, we conclude that the CDPT exists in the PD-2 model $`(\lambda =1)`$ whereas it disappears in the PD-2 model $`(\lambda =0)`$.
In Fig.9 we present $`f_P(x_{})`$ of the PD-3 model. It is obvious that there is no long-range order in the PD-3 model both for $`\lambda =0,1`$ and only the confinement phase is realized for all $`g`$.
### 4.3 Wilson loops
Let us turn to study of the Wilson loops. For ordinary pure and local gauge systems, the Wilson loop $`W[๐]`$ along a closed loop $`๐`$ on the lattice is a good order parameter to study the gauge dynamics; $`W[C]`$ obeys the area law in the confinement phase and the perimeter law in the deconfinement phase;
$`W[๐]{\displaystyle \underset{๐}{}}U_{x\mu }\{\begin{array}{cc}\mathrm{exp}(aS[๐]),\hfill & \mathrm{area}\mathrm{law},\hfill \\ \mathrm{exp}(a^{}L[๐]),\hfill & \mathrm{perimeter}\mathrm{law},\hfill \end{array}`$ (4.8)
where $`S[๐]`$ is the minimum area of a surface, the boundary of which is $`๐`$, and $`a`$ and $`a^{}`$ are constants. For a (local) gauge theory containing matter fields of the fundamental charge, $`W[๐]`$ cannot be an order parameter because the matter fields generate the terms $`_๐U_{x\mu }`$ with coefficients $`\mathrm{exp}(bL[๐])`$ in the effective action. However, in the present model (3.5), the nonlocal terms are restricted only along the temporal direction, so it is interesting to measure $`W[๐]`$ for the loops lying in the spatial (1-2) plane. If $`W[๐]`$ obeys the perimeter law, fluctuations of the spatial component of gauge field is small.
In Fig.10, we plot $`W[๐]`$. For the PD-1 model in Fig.10(a), the data at $`g=0.25`$ seem to prefer the perimeter law. For the ED model in Fig.10(b), the area law fits $`W[๐]`$ better than the perimeter law at $`g=1.5`$; a considerably larger value than $`g1.0`$ at the peak of $`C`$. This suggests that the area law holds in the ED model at all $`g`$. These observations are consistent with the previous results on the (non)existence of the CDPT. We conclude that the Wilson loops in the spatial plane can be used as an order parameter of gauge dynamics in the present model.
The result that the PD-1 case of the simplified model (3.5) exhibits the CDPT strongly suggests that the original model (3.3) of massless matter fields also has the deconfinement phase, because the isotropic distribution of the nonlocal gauge couplings in the original model should give the similar effect of suppression of fluctuations of $`U_{x\mu }`$ as those in the temporal direction in Eq.(3.5). This expectation is supported by the measurement of the spatial Wilson loop given above.
### 4.4 Instantons
It is well known in a 3D continuum space-time that instanton (monopole) configurations of U(1) gauge field carry nontrivial topological numbers, and the instanton density serves as an index to express the disorderness of gauge field. To see the details of th gauge dynamics of the present model and to support the conclusions obtained in the previous subsections, let us study instantons on the lattice. We employ the definition of the instanton density $`\rho _x`$ at the site $`x`$ in U(1) lattice gauge theories by DeGrand and Toussaint . We introduce the โvector potentialโ $`\theta _{x\mu }`$ as the exponent of $`U_{x\mu }=\mathrm{exp}(i\theta _{x\mu })`$ $`[\theta _{x\mu }(\pi ,\pi )]`$. Then, the magnetic flux $`\mathrm{\Theta }_{x,\mu \nu }`$ penetrating plaquette $`(x,x+\mu ,x+\mu +\nu ,x+\nu `$) is expressed as
$`\mathrm{\Theta }_{x,\mu \nu }\theta _{x\mu }+\theta _{x+\mu ,\nu }\theta _{x+\nu ,\mu }\theta _{x\nu },(4\pi <\mathrm{\Theta }_{x,\mu \nu }<4\pi ).`$ (4.9)
We decompose $`\mathrm{\Theta }_{x,\mu \nu }`$ into its integer part $`2\pi n_{x,\mu \nu }`$ ($`n_{x,\mu \nu }`$ is an integer) and the remaining part $`\stackrel{~}{\mathrm{\Theta }}_{x,\mu \nu }`$ $`\mathrm{\Theta }_{x,\mu \nu }(\text{mod}\mathrm{\hspace{0.33em}2}\pi `$) uniquely,
$$\mathrm{\Theta }_{x,\mu \nu }=2\pi n_{x,\mu \nu }+\stackrel{~}{\mathrm{\Theta }}_{x,\mu \nu },(\pi <\stackrel{~}{\mathrm{\Theta }}_{x,\mu \nu }<\pi ).$$
(4.10)
Physically speaking, $`n_{x,\mu \nu }`$ describes the Dirac string whereas $`\stackrel{~}{\mathrm{\Theta }}_{x,\mu \nu }`$ describes the fluctuations around it. The quantized instanton charge $`\rho _x`$ at the cube around the site $`\stackrel{~}{x}=x+(\widehat{1}+\widehat{2}+\widehat{3})/2`$ of the dual lattice is defined as
$`\rho _x`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu ,\nu ,\rho }{}}ฯต_{\mu \nu \rho }(n_{x+\mu ,\nu \rho }n_{x,\nu \rho })`$ (4.11)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{\mu ,\nu ,\rho }{}}ฯต_{\mu \nu \rho }(\stackrel{~}{\mathrm{\Theta }}_{x+\mu ,\nu \rho }\stackrel{~}{\mathrm{\Theta }}_{x,\nu \rho }),`$
where $`ฯต_{\mu \nu \rho }`$ is the complete antisymmetric tensor. $`\rho _x`$ measures the total flux emanating from the monopole(instanton) sitting at $`\stackrel{~}{x}`$. Roughly speaking, $`\rho _x`$ measures the strength of nonperturbative gauge fluctuations around $`\stackrel{~}{x}`$. For the local 3D U(1) compact lattice gauge theory without matter fields, the average density
$`\rho {\displaystyle \frac{1}{V}}{\displaystyle \underset{x}{}}|\rho _x|,`$ (4.12)
is known to behave as $`\rho \mathrm{exp}(cg)`$ ($`c`$ is a constant) if the instanton action $`cg`$ is large; The instanton gas stabilizes the confinement phase for all the gauge coupling.
In Fig.11 we present $`\rho `$ as a function of $`g`$ in the PD-1 and ED models. It decreases as $`g`$ increases more rapidly in the PD-1 model than in the ED model. This difference in behavior is consistent with the result that the PD-1 model exhibits a second-order transition, while the ED model does not. The $`\lambda `$ coupling enhances the rate of decrease in $`\rho `$ as one expects since the spatial coupling enhances the ordered deconfinement phase. In the ED model with $`\lambda =1`$, $`\rho `$ is fitted by $`\mathrm{exp}(cg)`$ in the dilute (large $`g`$) region, and the smooth increase for smaller $`g`$ indicates a crossover from the dilute gas of instantons to the dense gas, just the behavior similar to the case of pure and local lattice gauge theory.
In Ref. we presented the figure (Fig.3) which is very similar to Fig.11. However, the former was plotted by using the definition $`\rho V^1_x(1\delta _{0,\rho _x})`$, the average of occupation number of instantons (1 for $`\rho _x0`$ and 0 for $`\rho _x=0`$) per site instead of the instanton density of Eq.(4.12) itself. The curves in two figures are almost indistinguishable, because the configirations with $`|\rho _x|2`$ are rather rare.
In Fig.12 we present snapshots of $`\rho _x`$ for the PD-1 model with $`\lambda =1`$. Fig.12(a) is a dense gas whereas Fig.12(b) is a dilute gas. They are separated at $`g_c0.20`$, the location of the peak of $`C`$ for $`N=16`$. In Fig.12(b), instantons mostly appear in dipole pairs at nearest-neighbor sites, $`\rho _x=\pm 1,\rho _{x\pm \mu }=1`$, while in Fig.12(a), they appear densely and it is hard to determine their partners. In both cases, the distributions $`\rho _x`$ have no apparent anisotropies like column structures. However, the orientations of dipoles in Fig.12(b) are mostly ($`92\%`$) in the temporal direction as expected from the nonlocal interactions of Eq.(3.5).
Next, let us examine the difference of gauge dynamics in various cases of nonlocal interactions in detail. To this end, it is useful to measure new observables made out of gauge-field configurations; nonlocal instantons elongated in the temporal direction with the length $`t=2,3,\mathrm{},N1`$. The density of these instantons is defined by
$`\rho _{x,t}`$ $``$ $`{\displaystyle \underset{\tau =0}{\overset{t1}{}}}\rho _{x+\tau \widehat{0}},`$
$`\rho _t`$ $``$ $`\rho _{x,t}.`$ (4.13)
In Fig.13 we plot $`\rho _t`$ vs $`g`$ for various cases of the exponent $`\alpha `$ and $`\lambda `$.
We notice that $`\rho _t`$โs in Fig.13(a,b,c) behave quite differently from $`\rho _t`$โs in Fig.13(d,e). For the former cases, all the curves of $`\rho _t`$ with different $`t`$โs decrease similarly as $`g`$ increases \[except for $`t=1`$ in (c)\], whereas for the latter cases, each $`\rho _t`$ decreases in a different manner. The present nonlocal instanton density $`\rho _t`$ is capable to distinguish the cases (a,b,c) exhibiting a CDPT and the cases (d,e) without transition (cross over) through its $`t`$ dependence. This is consistent with our common understanding that the second-order phase transition is a collective phenomenon at which wild fluctuations of all the variables in the system in the disordered phase start to be reduced coherently and also scale-invariantly near the transition point.
## 5 Effective Models at Large $`g`$ and $`\lambda =0`$
In this section, we study a 1D XY spin model for the spatial gauge variables $`U_{xi}`$ and also a 2D XY spin model for the temporal ones $`U_{x0}`$, both of which are regarded as effective models for the present nonlocal model at large $`g`$ and $`\lambda =0`$ obtained by a simple reduction of the dynamical degrees of freedom. Two models are complementary each other and study of them helps us to understand the properties of the present nonlocal gauge model studied by the numerical simulations in section 4. In particular, these models are capable to explain the existence of the deconfinement (ordered) phase in the PD-1 model at large $`g`$ and the nonexistence of it in the other PD and ED models with $`\lambda =0`$.
### 5.1 The 1D XY model with spatial variables $`U_{xi}`$
Let us imagine the situation in which fluctuations of the temporal gauge fields $`U_{x0}`$ are small so that one may replace $`U_{x0}=\mathrm{exp}(i\theta _{x0})`$ by its average as follows;
$`U_{x0}u(0<u<1).`$ (5.1)
This is expected for sufficiently large $`g`$. Furthermore we put $`\lambda =0`$, i.e., there are no direct interactions among spatial gauge variables $`U_{xi}=\mathrm{exp}(i\theta _{xi})`$. Then the system is decoupled to 1D subsystems defined at each spatial link $`(x,x+i)`$. The subsystem at $`(x,x+i)`$ is described by the U(1) angles $`\theta _j(\theta _{xi})`$ (we write $`x_0j=1,\mathrm{},N_0`$ and suppress the suffix $`x_1,x_2,i`$). Its energy takes the form of a 1D XY spin model with nonlocal interactions with couplings $`c_\tau `$. The partition functions of the subsystem and the total system are given as follows;
$`Z_{1\mathrm{D}\mathrm{X}\mathrm{Y}}={\displaystyle \underset{j=1}{\overset{N_0}{}}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d\theta _j}{2\pi }}\mathrm{exp}\left[g{\displaystyle \underset{j=1}{\overset{N_0}{}}}{\displaystyle \underset{\tau =1}{\overset{N_0}{}}}u^{2\tau }c_\tau \mathrm{cos}(\theta _{j+\tau }\theta _j)\right],`$
$`Z_๐ฏ(g:\mathrm{large},\lambda =0)(Z_{1\mathrm{D}\mathrm{X}\mathrm{Y}})^{2N_1N_2}.`$ (5.2)
To study the correlation function $`\mathrm{exp}[i(\theta _r\theta _0)]`$, we make the harmonic approximation for large $`g`$; (i) expand the cosine term up to the quadratic term, (ii) extend the range of $`\theta _j`$ to $`(\mathrm{},\mathrm{})`$ and (iii) neglect topologically nontrivial configurations. We expect that the third approximation in the above is justified by the long-range ferromagnetic interactions in Eq.(5.2). Then we have
$`Z_{1\mathrm{D}\mathrm{X}\mathrm{Y}}`$ $``$ $`{\displaystyle \underset{j=1}{\overset{N_0}{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\theta _j}{2\pi }}\mathrm{exp}\left[{\displaystyle \frac{g}{2}}{\displaystyle \underset{j=1}{\overset{N_0}{}}}{\displaystyle \underset{\tau =1}{\overset{N_0}{}}}u^{2\tau }c_\tau \left(\theta _{j+\tau }\theta _j\right)^2\right]`$
$``$ $`{\displaystyle \underset{k=1}{\overset{N_0}{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\stackrel{~}{\theta }_k\mathrm{exp}\left[{\displaystyle \underset{k=1}{\overset{N_0}{}}}G_k\stackrel{~}{\theta }_k^2\right],`$
$`G_k`$ $`=`$ $`g{\displaystyle \underset{\tau }{}}u^{2\tau }c_\tau [1\mathrm{cos}(k\tau )],`$
$`f(r)`$ $``$ $`\mathrm{exp}[i(\theta _r\theta _0)]=\mathrm{exp}[{\displaystyle \frac{1}{2}}(\theta _r\theta _0)^2]`$ (5.3)
$`=`$ $`\mathrm{exp}[{\displaystyle \underset{k}{}}G^1(k)(1\mathrm{cos}(kr))],`$
where $`\stackrel{~}{\theta }_k`$ is the Fourier-transformed variable of $`\theta _j`$. For the standard local coupling $`c_\tau =\delta _{\tau 1}`$, the $`k`$-sum in the exponent of $`f(r)`$ gives $`๐k\mathrm{cos}(kr)/k^2r`$, which implies $`f(r)\mathrm{exp}[(gu^2)^1r]0`$ as $`r\mathrm{}`$ due to the severe infrared fluctuations. For the nonlocal cases we have
$`f(r)`$ $``$ $`\mathrm{exp}\left[{\displaystyle \underset{k}{}}{\displaystyle \frac{1\mathrm{cos}(kr)}{g_\tau u^{2\tau }c_\tau \tau ^2k^2}}\right]\mathrm{exp}(Mr),`$
$`M`$ $`=`$ $`{\displaystyle \frac{1}{g_\tau u^{2\tau }c_\tau \tau ^2}}\{\begin{array}{cc}0,\hfill & \mathrm{PD}1\hfill \\ \frac{1u^2}{gu^2},\hfill & \mathrm{PD}2\hfill \\ \frac{1}{g\mathrm{ln}(1u^2)},\hfill & \mathrm{PD}3\hfill \\ \frac{(1h)^3}{gh(1h+3h^2h^3)},(h=u^2e^1)\hfill & \mathrm{ED}\hfill \end{array}.`$ (5.8)
Namely, the order in the correlation function survives as $`N_0\mathrm{}`$ for the PD-1 case if $`u0`$, whereas the order is destroyed in the other cases even for $`u0`$. This result is just consistent with the numerical results for specific heat and the Polyakov lines studied in section 4. For $`\lambda =0`$ case, the ordered-deconfinement phase is realized only in the PD-1 model at large $`g`$.
### 5.2 The 2D XY model with temporal variables $`U_{x0}`$
The discussion in the previous subsection is supported by considering another effective model that focuses on the dynamics of the temporal gauge field $`U_{x0}`$. To obtain it, we first replace the spatial gauge field $`U_{xi}`$ by its average,
$`U_{xi}v(0<v<1).`$ (5.9)
This replacement is supported for the PD-1 model by the result of the 1D XY model (5.8) in the previous subsection. Then we fix the gauge to the temporal gauge by setting $`U_{x0}=1`$ except for $`x_0=N_0`$ because of the periodic boundary condition. By keeping only the terms in the energy involving $`U_{x_{},N_0,0}`$, we have
$`U_{x_{},N_0,0}`$ $``$ $`\mathrm{exp}(i\phi _x_{}),`$
$`Z_{2\mathrm{D}\mathrm{X}\mathrm{Y}}`$ $`=`$ $`{\displaystyle \underset{x_{}}{}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d\phi _x_{}}{2\pi }}\mathrm{exp}(A_{2\mathrm{D}\mathrm{X}\mathrm{Y}}),`$
$`A_{2\mathrm{D}\mathrm{X}\mathrm{Y}}`$ $`=`$ $`gv^2{\displaystyle \underset{x_{}}{}}{\displaystyle \underset{i=1,2}{}}{\displaystyle \underset{\tau =1}{\overset{N_0}{}}}c_\tau \mathrm{cos}(\phi _{x_{}+i}\phi _x_{}).`$ (5.10)
This is just the standard 2D XY spin model having the nearest-neighbor couplings among the XY spins $`\mathrm{exp}(i\phi _x_{})`$. Here the XY spin correlation function corresponds to the correlation function of the Polyakov lines studied in section 4,
$`\mathrm{cos}(\phi _x_{}\phi _0_{})=f_P(x_{}).`$ (5.11)
The spin stiffness in the effective 2D XY model is given as
$`gv^2Q_1\{\begin{array}{cc}\mathrm{},\hfill & \mathrm{PD}1\hfill \\ \frac{gv^2\pi ^2}{6},\hfill & \mathrm{PD}2\hfill \\ \frac{gv^2\pi ^4}{90},\hfill & \mathrm{PD}3\hfill \\ \frac{gv^2}{e1},\hfill & \mathrm{ED}\hfill \end{array}.`$ (5.16)
It is well known that the 2D XY model with finite spin stiffness has no long-range order, although the Kosterlitz-Thouless transition is possible. Only in the PD-1 case, the spin stiffness diverges for finite $`v`$ and the model may have the order $`U_{x_{},N_0,0}0`$. This result (5.16) supports the procedure (5.1) for the 1D XY model in the previous subsection for the PD-1 case. On the other hand, Eq.(5.8) supports the replacement (5.9) as we mentioned before. Thus studies of these two XY models give us the conclusion that only the PD-1 model at large $`g`$ has the ordered (Coulomb) phase for the $`\lambda =0`$ case.
Studies of the two effective XY models in this section predict that the cases of PD-2, PD-3 and ED with $`\lambda =0`$ have no Coulomb phase, as it was verified by the previous numerical calculations. Inclusion of $`\lambda `$ term, however, generates direct interactions between the spatial variables $`U_{xi}`$. For small $`\lambda `$, $`U_{xi}`$โs ($`i=1,2`$) can be integrated out perturbatively in powers of $`\lambda `$. The obtained effective model contains nonlocal and multi-body interactions of $`\phi _x_{}`$โs which prefer the ferromagnetic order of $`\phi _x_{}`$. Therefore the effective model for $`\lambda 1`$ may not belong to the same universality class of the 2DXY model with local interactions. Then the above result does not contradict the numerical result in the previous section, which shows the deconfinement phase exists in the PD-2 model with $`\lambda =1`$.
## 6 Conclusion
In this paper, we studied the nonlocal compact U(1) gauge theory on the 3D lattice, which โsimulatesโ gauge models coupled with massless/massive matter fields. The main contributions of the present paper may be the following two points: (i) MC simulations are feasible within reasonable computer time even for 3D lattice gauge theories with nonlocal interactions along one direction, and (ii) the measurements of $`E`$,$`C`$, Polyakov lines, Wilson loops, and local and nonlocal instantons give rise to clear and consistent results on the phase structure of the model. In particular, they distinguish the cases of ND($`\lambda =1,0`$), PD-1($`\lambda =1,0`$), PD-2($`\lambda =1`$) with a CDPT from the other cases of PD-2($`\lambda =0`$), PD-3,ED without CDPT in a definitive manner.
As explained in Sect.2.3, the results obtained in this paper are quite important for studies of the strongly-correlated electron systems like the high-$`T_c`$ cuprates, the fractional quantum Hall effect, quantum spin models, etc. For example, in the t-J model of high-$`T_c`$ superconductivity, by using the hopping expansion of holons and spinons at finite $`T`$ with the continuous imaginary time, we derived an effective gauge theory, which is highly nonlocal in the temporal direction. The obtained effective theory has a similar action as Eq.(3.5) with $`c_\tau =`$constant and $`gn`$ where $`n`$ is the density of matter fields(holons and spinons). This corresponds to the ND model. Although the above effective gauge model is obtained for the system at finite $`T`$, we expect that a similar gauge model appears as an effective model at $`T=0`$. The result that the ND model has a CDPT strongly suggests that the t-J model has the corresponding phase transition into the deconfinement phase, which is nothing but the charge-spin separated phase.
In the deconfinement phase, all the three variables $`U_{x\mu }(\mu =0,1,2)`$ are stable, having small fluctuations. The stability of the Lagrange multiplier $`U_{x0}`$ means that the constraint (2.6) is not respected by the holons and spinons, the low-energy excitations of the system, whereas the stability of $`U_{xi}`$ indicates that the gauge interaction between the holons and spinons can be treated perturbatively. Therefore, quasi-particles in the CSS state are the holons, spinons and weakly interacting gauge bosons. Of course, in order to present a definite โproofโ of the CSS, it is necessary to investigate a gauge system with full isotropic nonlocal interaction, because the integration over holons and spinons generates nonlocal interactions not only in the temporal direction but also in the spatial directions and their combinations.
Another interesting model related with the present one is the U(1) Higgs model coupled with the nonlocal gauge field. At present, it is believed that there is no phase transition in the 3D U(1) gauge-Higgs model with the ordinary local action if the Higgs field has the fundamental charge and its radial fluctuations are suppressed. However, the situation may be changed by nonlocal gauge interactions. The existence of the deconfinement phase in the present nonlocal gauge system without Higgs fields suggests that all the three phases, i.e., the confinement, Coulomb and Higgs phases, may be realized in the 3D nonlocal gauge model with a local coupling to a Higgs field. The deconfinement phase of the present model corresponds to the Coulomb phase. This problem is closely related with โdoped holesโ in the algebraic spin liquid which may be realized in certain antiferromagnetic spin models and materials in the spatial 2D lattice. We shall report on these problems in a separate publication.
Acknowledgement
One of the authors (K.S.) thanks the members of Department of Physics, Kanazawa University for their hospitality delivered to him during his stay.
## Appendix A High- and Low-Temperature Expansions
In this appendix we study the behavior of the nonlocal gauge theory (3.5) in two regions of $`g`$ by analytic methods; the region of small $`g(<<1)`$ by high-temperature expansion (HTE) in Sect.A.1 and the region of large $`g(>>1)`$ by low-temperature expansion (LTE) in Sect.A.2. Once one obtains an approximate expression for the partition function $`Z_๐ฏ`$, one can calculate various thermodynamic quantities like $`E`$ and $`C`$ in Eq.(4.1).
### A.1 High-temperature expansion (HTE) for small $`g`$
Let us consider the case of small $`g`$. Since the action $`A_๐ฏ`$ is proportional to $`g`$, one may expand the partition function $`Z_๐ฏ`$ of (3.5) in powers of $`g`$ as
$$Z_๐ฏ=[dU]\mathrm{exp}(A_๐ฏ)=[dU]\underset{n=0}{\overset{\mathrm{}}{}}\frac{(A_๐ฏ)^n}{n!}.$$
(A.1)
We obtain the expansion up to $`O(g^4)`$ as
$$Z_๐ฏ=1+\left(B_{2T}+B_{2S}\lambda ^2\right)g^2+B_{3T}g^3+\left(B_{4T}+B_{4TS}\lambda ^2+B_{4S}\lambda ^4\right)g^4+O(g^5),$$
(A.2)
where each coefficient is expressed as
$`B_{nT}`$ $`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle [dU]\left[\underset{q}{}c_\tau (V_q+\overline{V}_q)\right]^n},q(x,i,\tau ),`$
$`B_{nS}`$ $`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle [dU]\left[\underset{p}{}(U_p+\overline{U}_p)\right]^n},U_p\overline{U}_{x+\widehat{2},1}U_{x+\widehat{1},2}U_{x1}\overline{U}_{x2},`$
$`B_{4TS}`$ $`=`$ $`{\displaystyle \frac{{}_{4}{}^{}C_{2}^{}}{4!}}{\displaystyle [dU]\left[\underset{p}{}(U_p+\overline{U}_p)\underset{q}{}c_\tau (V_q+\overline{V}_q)\right]^2}.`$ (A.3)
To evaluate the above coefficients, we use the following formula for U(1) integral,
$`{\displaystyle ๐U_{x\mu }\overline{U}_{x\mu }^mU_{y\nu }^n}=\delta _{xy}\delta _{\mu \nu }\delta _{mn}.`$ (A.4)
Then we obtain the following result;
$`B_{2T}`$ $`=`$ $`{\displaystyle \underset{q}{}}c_\tau ^2=2VQ_2,Q_2{\displaystyle \underset{\tau }{}}c_\tau ^2,`$
$`B_{2S}`$ $`=`$ $`{\displaystyle \underset{p}{}}1=V,`$
$`B_{3T}`$ $`=`$ $`{\displaystyle \frac{32}{3!}}{\displaystyle \underset{x}{}}{\displaystyle \underset{i}{}}{\displaystyle \underset{\tau _1}{}}{\displaystyle \underset{\tau _2}{}}c_{\tau _1}c_{\tau _2}c_{\tau _3}V_{x,i,\tau _1}V_{x+\tau _1\widehat{0},i,\tau _2}\overline{V}_{x,i,\tau _3}\delta _{\tau _3,\tau _1+\tau _2}+\mathrm{c}.\mathrm{c}.`$
$`=`$ $`4V\stackrel{~}{Q}_3,\stackrel{~}{Q}_3{\displaystyle \underset{\tau _1=1}{\overset{N}{}}}{\displaystyle \underset{\tau _2=1}{\overset{N}{}}}{\displaystyle \underset{\tau _3=1}{\overset{N}{}}}c_{\tau _1}c_{\tau _2}c_{\tau _3}\delta _{\tau _3,\tau _1+\tau _2},`$
$`B_{4T}`$ $`=`$ $`B_{4T}^a+B_{4T}^b+B_{4T}^c,`$
$`B_{4T}^a`$ $`=`$ $`{\displaystyle \frac{1}{4!}}{\displaystyle \underset{q_1,q_2,q_3,q_4}{}}{\displaystyle \underset{\mathrm{}=1}{\overset{4}{}}}c_\tau _{\mathrm{}}{}_{4}{}^{}C_{2}^{}\left[\delta _{q_1q_2}\delta _{q_3q_4}+\delta _{q_1q_3}\delta _{q_2q_4}\delta _{q_1q_2}\delta _{q_2q_3}\delta _{q_3q_4}\right]`$
$`=`$ $`2V^2Q_2^2{\displaystyle \frac{1}{2}}VQ_4,Q_4{\displaystyle \underset{\tau }{}}c_\tau ^4,`$
$`B_{4T}^b`$ $`=`$ $`{\displaystyle \frac{432}{4!}}{\displaystyle \underset{x,i}{}}{\displaystyle \underset{\tau _1,\tau _2,\tau _3,\tau _4}{}}c_{\tau _1}c_{\tau _2}c_{\tau _3}c_{\tau _4}(1\delta _{\tau _1\tau _3})\delta _{\tau _4,\tau _1+\tau _2\tau _3}V_{x,i,\tau _1}V_{x+\tau _1\widehat{0},i,\tau _2}\overline{V}_{x,i,\tau _3}\overline{V}_{x+\tau _3\widehat{0},i,\tau _4}`$
$`=`$ $`2V(\stackrel{~}{Q}_4Q_2^2),\stackrel{~}{Q}_4{\displaystyle \underset{\tau _1,\tau _2,\tau _3,\tau _4}{}}c_{\tau _1}c_{\tau _1}c_{\tau _3}c_{\tau _4}\delta _{\tau _4,\tau _1+\tau _2\tau _3},`$
$`B_{4T}^c`$ $`=`$ $`{\displaystyle \frac{4!}{4!}}{\displaystyle \underset{x,i}{}}{\displaystyle \underset{\tau _1,\tau _2,\tau _3,\tau _4}{}}c_{\tau _1}c_{\tau _2}c_{\tau _3}c_{\tau _4}V_{x,i,\tau _1}V_{x+\tau _1\widehat{0},i,\tau _2}V_{x+(\tau _1+\tau _2)\widehat{0},i,\tau _3}\overline{V}_{x,i,\tau _4}\delta _{\tau _4,\tau _1+\tau _2+\tau _3}+\mathrm{c}.\mathrm{c}.`$
$`=`$ $`4V\stackrel{~}{Q}_{4+},\stackrel{~}{Q}_{4+}{\displaystyle \underset{\tau _1,\tau _2,\tau _3,\tau _4}{}}c_{\tau _1}c_{\tau _1}c_{\tau _3}c_{\tau _4}\delta _{\tau _4,\tau _1+\tau _2+\tau _3},`$
$`B_{4TS}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{p}{}}2{\displaystyle \underset{q}{}}2c_\tau ^2=2V^2Q_2,`$
$`B_{4S}`$ $`=`$ $`{\displaystyle \frac{1}{4!}}{\displaystyle \underset{p_1,p_2,p_3,p_4}{}}{}_{4}{}^{}C_{2}^{}\left[\delta _{p_1p_2}\delta _{p_3p_4}+\delta _{p_1p_3}\delta _{p_2p_4}\delta _{p_1p_2}\delta _{p_2p_3}\delta _{p_3p_4}\right]`$ (A.5)
$`=`$ $`{\displaystyle \frac{1}{2}}V^2{\displaystyle \frac{1}{4}}V.`$
Some of them are illustrated in Fig.14.
Finally, we obtain the expression of the partition function and the free energy as follows,
$`Z_๐ฏ`$ $`=`$ $`\mathrm{exp}(FV),`$
$`F`$ $`=`$ $`(2Q_2+\lambda ^2)g^24\stackrel{~}{Q}_3g^3+\left({\displaystyle \frac{Q_4}{2}}+{\displaystyle \frac{\lambda ^4}{4}}4\stackrel{~}{Q}_{4+}2\stackrel{~}{Q}_4\right)g^4+O(g^5).`$ (A.6)
In the following, we list up the values of $`Q_2`$ and $`Q_4`$ for $`N\mathrm{}`$,
(A.12)
By using $`Z_๐ฏ`$ of Eq.(A.6), we obtain the internal energy $`E`$ and the specific heat $`C`$ defined by Eq.(4.1) as
$`E`$ $`=`$ $`2(2Q_2+\lambda ^2)g^212\stackrel{~}{Q}_3g^3+4\left({\displaystyle \frac{Q_4}{2}}+{\displaystyle \frac{\lambda ^4}{4}}4\stackrel{~}{Q}_{4+}2\stackrel{~}{Q}_4\right)g^4+O(g^5),`$
$`C`$ $`=`$ $`2(2Q_2+\lambda ^2)g^2+24\stackrel{~}{Q}_3g^312\left({\displaystyle \frac{Q_4}{2}}+{\displaystyle \frac{\lambda ^4}{4}}4\stackrel{~}{Q}_{4+}2\stackrel{~}{Q}_4\right)g^4+O(g^5).`$ (A.13)
Fig.15 shows that, as one includes the higher-order terms, the HTE results approaches the MC result systematically. However, the approach is rather slow compared with the related models of local interactions like the 3D XY spin model or the 3D U(1) pure LGT. This is because the present nonlocal interactions generate various important higher-order terms in the HTE that are absent from the local models.
Let us comment on the convergence of the HTE. As usual, the HTE is an expansion in the disordered (confinement) phase in which $`U_{x\mu }`$ fluctuates wildly. Equation (A.6) shows that the convergence radius $`g_{\mathrm{HTE}}`$ of the expansion is finite $`g_{\mathrm{HTE}}0`$, because both the harmonic numbers $`Q_2`$ and $`Q_4`$ appearing in the coefficients are finite. This means that there exists certainly the finite region $`0g^2<g_{\mathrm{HTE}}^2`$ of the confinement phase. We notice that if the long-range interaction $`c_\tau `$ is very strong such that $`Q_2=\mathrm{}`$, the confinement phase may disappear.
### A.2 Low-Temperature Expansion (LTE) for large $`g`$
For large $`g`$, we evaluate $`Z_๐ฏ`$ by the LTE. The LTE is an expansion in powers of $`g^1`$ around a fixed โlowest-energy configurationโ of $`U_{x\mu }`$ like $`U_{x\mu }=1`$, which gives rise to the global maximum of $`A_๐ฏ`$. Let us expand $`U_{x\mu }`$ as
$`U_{x\mu }\mathrm{exp}(i\theta _{x\mu })=1+i\theta _{x\mu }{\displaystyle \frac{1}{2}}\theta _{x\mu }^2+O(\theta _{x\mu }^3),`$ (A.14)
where $`\theta _{x\mu }`$ is treated as $`O(g^{1/2})`$ as we shall see. $`A_๐ฏ`$ is expanded up to the second order in $`\theta _{x\mu }`$ in the following quadratic form;
$`A_๐ฏ=4gQ_1V+2g\lambda Vg{\displaystyle \underset{x,\mu }{}}{\displaystyle \underset{y,\nu }{}}\theta _{x\mu }G_{x\mu ,y\nu }(\lambda )\theta _{y\nu }+O(\theta ^4),`$ (A.15)
where the first two terms $`4gQ_1V+2g\lambda V`$ come from the first term, unity, of R.H.S. of Eq.(A.14). Due to the gauge invariance, one may extend the region of $`\theta _{x\mu }`$ from $`\theta _{x\mu }(\pi ,\pi )`$ to $`\theta _{x\mu }(\mathrm{},\mathrm{})`$ together with a gauge fixing. We take the temporal gauge $`\theta _{x0}=0`$ in the following calculation. Then we evaluate $`Z_๐ฏ`$ by rescaling $`\theta _{xi}^{}=g^{1/2}\theta _{xi}`$ and performing Gaussian integration as
$`Z_๐ฏ`$ $``$ $`e^{(4Q_1+2\lambda )gV}{\displaystyle \underset{x}{}}\left[{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\theta _{x1}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\theta _{x2}\right]\mathrm{exp}\left(g{\displaystyle \underset{x,\mu }{}}{\displaystyle \underset{y,\nu }{}}\theta _{x\mu }G_{x\mu ,y\nu }\theta _{y\nu }\right)`$
$`=`$ $`e^{(4Q_1+2\lambda )gV}{\displaystyle \underset{x}{}}\left[g^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\theta _{x1}^{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\theta _{x2}^{}\right]\mathrm{exp}\left({\displaystyle \underset{x,\mu }{}}{\displaystyle \underset{y,\nu }{}}\theta _{x\mu }^{}G_{x\mu ,y\nu }\theta _{y\nu }^{}\right)`$
$`=`$ $`\mathrm{exp}\left[(4Q_1g+2\lambda g\mathrm{ln}g)V{\displaystyle \frac{1}{2}}\mathrm{Tr}\mathrm{ln}G(\lambda )\right],`$
$`F`$ $`=`$ $`(4Q_1+2\lambda )g+\mathrm{ln}g+O(g^0).`$ (A.16)
This gives
$`E`$ $`=`$ $`(4Q_1+2\lambda )g+1+O(g^1),`$
$`C`$ $`=`$ $`1+O(g^1).`$ (A.17)
The higher-order terms in $`F`$ are $`O(g^n)(n0)`$ which may be calculated by the usual perturbation theory. |
warning/0506/math0506193.html | ar5iv | text | # Braid action on derived category of Nakayama algebras
## Introduction
<sup>0</sup><sup>0</sup>0 2000 Mathematics Subject Classification 16D90,16G20,18E30,20F36
The author was supported by EPSRC Grant number GR/S35387/01.
For a commutative ring $`R`$ we shall study the group of self-equivalences of the derived category of an $`R`$-algebra $`A`$ which is projective as an $`R`$-module. We denote such a group by $`TrPic(A)`$ and call it the Derived Picard group.
Rouquier and Zimmermann , motivated by Brouรฉโs conjecture, studied the $`TrPic(A)`$ of Brauer tree algebras $`A`$ (with no exceptional vertex) by defining an action of the Artin braid group $`B(A_n)`$ on $`D^b(A)`$, the bounded derived category of $`A.`$ Khovanov and Seidel , motivated by mirror symmetry, defined an action of $`B(A_n)`$ on bounded derived category of certain algebras similar to Brauer tree algebras, which turn out to be related to certain classes of representations of simple Lie algebras, and they proved that the action is faithful.
In this paper we will define an action of $`B(K_n),`$ the braid group whose associated Coxeter group is given by the complete graph on $`n`$ vertices, on $`D^b(N_n^n)`$, the bounded derived category of the Nakayama algebra $`N_n^n`$ which is a basic Brauer star algebra with no exceptional vertex.
We will also define a group homomorphism $`\eta `$ from $`B(K_n)`$ to $`B(A_n)`$ and give a categorification of $`\eta .`$ This categorification provides a link between the action of $`B(K_n)`$ on $`D^b(N_n^n)`$ and the faithful action of $`B(A_n)`$ on derived category of Brauer tree algebras defined by Rouquier and Zimmermann and Khovanov and Seidel .
In Schaps and Zakay-Illouz defined an action of affine braid group $`B(\stackrel{~}{A}_{n1})`$ on derived category of Brauer star algebras and left open the question whether this action is faithful.
We will provide a positive answer to this question for Brauer star algebras with no exceptional vertex. We will use the categorification of $`\eta `$ stated above, another categorification of braid group embeddings defined in , and and the faithfulness of the action in and .
The paper is organized as follows.
In Section 1, we review Rickardโs Morita theory for derived equivalence and we give the definition of derived Picard groups $`TrPic`$ and of Nakayama algebras.
In Section 2, we recall the Nakayama algebra $`N_n^n`$ which is also a basic Brauer star algebra with $`n`$ edges and no exceptional vertex. We define explicit elements of $`TrPic(N_n^n).`$
The action of the braid group $`B(K_n)`$ on $`D^b(N_n^n)`$ is defined in Section 3. And in Section 4 we show that this action is not faithful using a categorification of the group homomorphism $`\eta .`$
In the last section, we use the results of sections 2, 3 and 4 to show that the action defined in is faithful for Brauer star algebras with no exceptional vertex.
Acknowledgement I would like to thank Robert Marsh and Steffen Kรถnig for numerous discussions on studying the braid action on derived categories. I also thank Alexander Zimmermann for encouraging me to prove the faithfulness of the action defined in .
## 1 Derived Picard groups and Nakayama algebras
### 1.1 Derived equivalence
Let $`A`$ and $`B`$ be algebras over a commutative ring $`R`$ and assume that $`A`$ and $`B`$ are projective as modules over $`R.`$ The bounded derived category $`D^b(A)`$ is the category whose objects are complexes of finitely generated projective modules which are bounded to the right and which have nonzero homology only in finitely many degrees. Morphisms are morphisms of complexes up to homotopy. We refer to for the definition of derived categories.
Denote by $`Aper`$ the full subcategory of $`D^b(A)`$ consisting of the perfect complexes, i.e. bounded complexes of finitely generated projective $`A`$-modules.
Rickard and Keller have given a necessary and sufficient criterion for the existence of derived equivalences between two rings $`A`$ and $`B.`$
###### Theorem 1.1
(Rickard and Keller )
The following statements are equivalent :
1. $`D^b(A)D^b(B)`$ as triangulated categories.
2. There is a complex $`T`$ in $`Aper`$ such that
1. $`Hom(T,T[i])=0`$ for $`i0,`$
2. $`Aper`$ is generated by $`T,`$ and
3. $`End_{D^b(A)}(T)B.`$
3. There is a bounded complex $`X`$ in $`D^b(A_RB^{op})`$ whose restrictions to $`A`$ and to $`B^{op}`$ are perfect and a bounded complex $`Y`$ in $`D^b(B_RA^{op})`$ whose restrictions to $`B`$ and $`A^{op}`$ are perfect such that
$$X_B^LYA\text{in}D^b(AA^{op})\text{and}Y_A^LXB\text{in}D^b(BB^{op}).$$
A complex $`T`$ satisfying the condition in 2 is called a tilting complex for $`A`$ and the complexes $`X`$ and $`Y`$ in 3 are called two-sided tilting complexes inverse to each other. The image and the pre-image of an indecomposable projective $`B`$-module under derived equivalence $`D^b(A)D^b(B)`$ is a partial tilting complex in $`A`$, i.e. a complex satisfying the condition in 2(a).
We define the Derived Picard group of $`A`$ as
$$\begin{array}{ccc}TrPic(A)& =& \left\{\begin{array}{c}\text{isomorphism classes of two-sided tilting complexes}\\ \text{ in }D^b(A_RA^{op})\end{array}\right\}\end{array}$$
where the product of the classes of $`X`$ and $`Y`$ is given by the class of $`X_AY.`$
It is clear that if $`A`$ and $`B`$ are projective as $`R`$-modules then $`D^b(A)D^b(B)`$ implies
$$TrPic(A)TrPic(B).$$
The Picard group $`Pic(A)`$ (i.e. the group of isomorphism classes of invertible $`AA^{op}`$-modules $`M`$ where $`M_AM^{}A`$) can be embedded into $`TrPic(A)`$ by sending an invertible $`AA^{op}`$-module $`M`$ to a stalk complex concentrated in degree 0.
### 1.2 Nakayama algebras
The self-injective Nakayama algebra $`N_m^n`$ is the path algebra over a field $`K`$ of the following quiver
$$\begin{array}{ccccccc}& & 2& & 3& & \\ & & & & & & \\ 1& & & & & & 4\\ & & & & & & \\ & & m& & \mathrm{}& & \end{array}$$
modulo the ideal generated by all compositions of $`n+1`$ consecutive arrows. By $`N_m^n`$ is a representative of the class of standard self-injective algebras of finite representation type associated to the Dynkin diagram of type $`A_n.`$
The algebra $`N_m^n`$ is symmetric if and only if $`m`$ divides $`n.`$ In this case the Nakayama algebra $`N_m^n`$ is a basic Brauer star algebra with $`n`$ edges and multiplicity $`n/m`$, a particular Brauer tree algebra where all the edges are adjacent to the exceptional vertex. A Brauer tree algebra with $`n`$ edges and multiplicity $`n/m`$ is (up to Morita equivalence) a $`p`$-block with cyclic defect group of order $`p^d=m+1`$ where the number of its simple modules is $`n.`$ (See for the definition of Brauer tree algebras).
By \[13, Theorem 4.2\], up to derived equivalence, a Brauer tree algebra is determined by the number of edges of the Brauer tree and the multiplicity of the exceptional vertex. Hence, an arbitrary Brauer tree algebra is derived equivalent to a Brauer star algebra associated with a star having the same number of edges.
If $`m=n,`$ the Nakayama algebra $`N_n^n`$ is the trivial extension algebra of the path algebra of the quiver
$$123\mathrm{}n1n.$$
It is a Brauer star algebra with $`n`$ edges and multiplicity 1.
Notation : From now on, $`A`$ is the Nakayama algebra $`N_n^n`$ over $`K`$ and $`B`$ is the Brauer tree algebra associated to a line without exceptional vertex. For projective $`A`$-modules we use the letter $`P`$ and for projective $`B`$-modules we use the letter $`Q.`$ We use the letters $`F,`$ $`R,`$ and $`H`$ to denote two-sided tilting complexes (elements of $`TrPic`$) and the letters $`,`$ $`,`$ and $``$ for the corresponding auto-equivalences. The symbol $``$ means $`_K.`$
## 2 Two-sided tilting complexes for Nakayama algebra $`N_n^n`$
Denote by $`(i_1i_2\mathrm{}i_{l+1})`$ the path starting at $`i_1`$ and ending at $`i_{l+1}.`$ Note that we have $`i_{k+1}i_k=1(modn)`$ for all $`1kl.`$
Let $`A`$ be the Nakayama algebra $`N_n^n.`$ The algebra $`A`$ is the path algebra of the quiver
$$\begin{array}{ccccccc}& & 2& & 3& & \\ & & & & & & \\ 1& & & & & & 4\\ & & & & & & \\ & & n& & \mathrm{}& & \end{array}$$
modulo the relations : $`(ii+1\mathrm{}n\mathrm{\hspace{0.33em}\hspace{0.33em}1\hspace{0.33em}\hspace{0.33em}2}\mathrm{}ii+1)=0`$ for all $`1in.`$
The paths $`(i)`$ of length 0 are mutually orthogonal idempotents and their sum is the unit element. The Loewy series of the indecomposable projective modules $`P_i=A(i)`$ are as follows :
$$P_i=\begin{array}{c}S_i\\ S_{i+1}\\ \mathrm{}\\ S_n\\ S_1\\ S_2\\ \mathrm{}\\ S_i\end{array}$$
The dimension of homomorphisms between projective modules is given by :
$$dim_KHom_A(P_i,P_j)=\{\begin{array}{cc}2& \text{if }i=j,\\ 1& \text{if }ij.\end{array}$$
Define $`{}_{i}{}^{}P=(i)AHom_A(P_i,A).`$ Define the complexes
$$F_i:=0P_i_iP\stackrel{\alpha _i}{}A0$$
and
$$F_i^{}:=0A\stackrel{\beta _i}{}P_i_iP0$$
where $`A`$ is in degree 0 and
$$\begin{array}{ccccccc}\alpha _i& :& P_i& & {}_{i}{}^{}P& & A\\ & & (i)& & (i)& & (i)\end{array}$$
$`\beta _i:`$ $`A`$ $`P_i_iP`$
$`1`$ $`(i\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i)(i)+(i)(i\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i)`$
$`+(i+1\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i)(ii+1)`$
$`+(i+2\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i)(ii+1i+2)`$
$`+\mathrm{}+(n\mathrm{\hspace{0.33em}1}\mathrm{}i)(ii+1\mathrm{}n)`$
$`+(1\mathrm{}i)(i\mathrm{}n\mathrm{\hspace{0.33em}1})+(2\mathrm{}i)(i\mathrm{}n\mathrm{\hspace{0.33em}1\hspace{0.33em}2})`$
$`+\mathrm{}+(i1i)(i\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i1).`$
###### Theorem 2.1
The complexes $`F_i`$ and $`F_i^{}`$ are elements of $`TrPic(A).`$
Proof We will show that $`F_i`$ and $`F_i^{}`$ are two-sided tilting complexes and inverses to each other and to do that, we will show that $`F_i_AF_i^{}`$ is homotopy equivalent to $`A`$ as complex of $`AA`$bimodules. (See \[15, Theorem 6\]). We follow the method of \[16, Theorem 4.1\] and \[10, Proposition 2.4\]. We have
$$F_i_AF_i^{}=(0P_i_iP\stackrel{d^1}{}A(P_iU_iP)\stackrel{d^0}{}P_i_A{}_{i}{}^{}P0)$$
where $`U`$ is the 2-dimensional space $`{}_{i}{}^{}P_AP_i`$ with a basis $`u_1=(i)(i)`$ and $`u_2=(i\mathrm{}i)(i),`$ and the differentials are given by
$$d^1=\alpha _i+\tau \text{and}d^0=(\beta _i,\delta )\text{where}$$
$`\tau (xy)`$ $`=`$ $`xu_1(i\mathrm{}i)y+xu_2y,`$
$`\delta (xu_1y)`$ $`=`$ $`xy,`$
$`\delta (xu_2y)`$ $`=`$ $`x(i\mathrm{}i)y.`$
The map $`d^0`$ is surjective since $`\delta `$ is surjective. Therefore, since $`P_i_iP`$ is projective, $`d^0`$ is a split surjection.
Denote by $`_i`$ (resp. $`{}_{i}{}^{}`$) the basis of $`P_i`$ (resp. of $`{}_{i}{}^{}P`$) where
$$_i=\{(i),(i1i),(i2iii),\mathrm{},(n\mathrm{\hspace{0.33em}1\hspace{0.33em}2}\mathrm{}i),\mathrm{},(ii+1\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i)\}\text{and}$$
$${}_{i}{}^{}=\{(i),(ii+1),(ii+1,i+2),\mathrm{},(ii+1\mathrm{}n\mathrm{\hspace{0.33em}1}),\mathrm{},(ii+1\mathrm{}n\mathrm{\hspace{0.33em}1}\mathrm{}i)\}.$$
Let $`x_i`$ and $`y_i.`$ Then,
$`\tau (xy)`$ $`=`$ $`\{\begin{array}{ccc}xu_1(i\mathrm{}i)+xu_2(i)& \text{if}& y=(i),\\ xu_2y& \text{if}& y(i).\end{array}`$
Therefore, if $`xyP_i_iP,`$ then $`xy0`$ implies $`\tau (xy)0.`$ This means that $`\tau `$ is injective, therefore so is $`d^1.`$ This implies $`d^1`$ is a split injection.
Therefore, $`F_i_AF_i^{}`$ is homotopy equivalent to a module $`V`$ which satisfies
$$(P_i_iP)(P_i_iP)V(P_i_iP_AP_i_iP)A.$$
Since $`dim_K(P_i_A{}_{i}{}^{}P)=dim_KHom_A(P_i,P_i)=2,`$ we obtain $`VA`$ as a left module and this finishes the proof. $`\mathrm{}`$
Denote by
$$_i:=F_i_A.$$
It follows from above that :
$$_i^1=F_i^{}_A.$$
###### Lemma 2.2
We have
$$_i(P_j)=\{\begin{array}{cccccccc}0& & P_i& & 0& & & \text{if }i=j\\ 0& & P_i& & P_j& & 0& \text{if }ij\end{array}$$
$$_i^1(P_j)=\{\begin{array}{cccccccc}& & 0& & P_i& & 0& \text{if }i=j\\ 0& & P_j& & P_i& & 0& \text{if }ij\end{array}$$
where $`P_j`$ is in degree 0.
Proof The images of $`P_j`$ under $`_i`$ and $`_i^1`$ are
$$_i(P_j)=(0P_i_iP_AP_j\stackrel{\alpha _iid_{P_j}}{}A_AP_j0)\text{and}$$
$$_i^1(P_j)=(0A_AP_j\stackrel{\beta _iid_{P_j}}{}P_i_iP_AP_j0).$$
When $`ij,`$ $`{}_{i}{}^{}P_AP_j`$ is generated by the path $`(i\mathrm{}j),`$ hence
$$\alpha _iid_{P_j}(xy_Az)=x(i\mathrm{}j)\text{for all}xy_AzP_i_iP_AP_j,$$
$$\beta _iid_{P_j}(1_Ax)=x(j\mathrm{}i)(i\mathrm{}j)\text{for all}xP_j.$$
Since $`dim_K(_iP_AP_j)=dim_KHom_A(P_i,P_j)=1`$ we have $`P_i_iP_AP_jP_i.`$ The only map up to scalars between $`P_i`$ and $`P_j`$ is
$$\mu _{ij}:P_iP_j\text{where}\mu _{ij}(x)=x(i\mathrm{}j)\text{for all}xP_i.$$
Therefore, $`\alpha _iid_{P_j}\mu _{ij},`$ $`\beta _iid_{P_j}\mu _{ji}`$ and
$$_i(P_j)=(0P_i\stackrel{\mu _{ij}}{}P_j0)\text{and}_i^1(P_j)=(0P_j\stackrel{\mu _{ji}}{}P_i0).$$
When $`i=j,`$ $`P_i_iP_AP_i\stackrel{\alpha _iid_{P_i}}{}A_AP_i`$ is surjective since for all $`xP_i,`$
$$(\alpha _iid_{P_i})(x(i)(i))=x.$$
Now if $`xP_i`$, then one can check that
$$\beta _iid_{P_i}(1_Ax)=x(i\mathrm{}i)+x(i\mathrm{}i)(i).$$
Hence, if $`x0,`$ then $`\beta _iid_{P_i}(1_Ax)0;`$ therefore, $`A_AP_i\stackrel{\beta _iid_{P_i}}{}P_i_iP_AP_i`$ is injective.
Hence, $`_i(P_i)`$ (resp. $`_i^1(P_i)`$) has homology concentrated in degree 1 (resp. $`1`$). As $`dim_K(_iP_AP_i)=2,`$ $`_i(P_i)P_i[1]`$ and $`_i^1(P_i)P_i[1].`$ $`\mathrm{}`$
Denote by $`\mathrm{\Delta }`$ the subgroup of $`TrPic(A)`$ generated by $`F_1,\mathrm{},F_n.`$
###### Lemma 2.3
For $`F,F^{}\mathrm{\Delta },`$ $`FF^{}`$ if and only if
$$F_APF^{}_AP$$
for any indecomposable projective module $`P.`$
Proof Let us assume that $`F_APF^{}_AP`$ for any indecomposable projective $`A`$-module $`P.`$ Denote by $`B`$ the Brauer tree algebra associated to a line with $`n`$ edges and multiplicity 1. Denote by $`\mathrm{\Phi }`$ the isomorphism of groups
$$\mathrm{\Phi }:TrPic(A)TrPic(B)$$
and by $`G`$ and $`G^{}`$ the images of $`F`$ and $`F^{}`$ under $`\mathrm{\Phi }`$ ($`G:=\mathrm{\Phi }(F)`$ and $`G^{}:=\mathrm{\Phi }(F^{})).`$ Denote by $`X`$ the associated two-sided tilting complex in $`D^b(BA^{op})`$ giving the derived equivalence between $`A`$ and $`B,`$ so we have
$$G^1_BXX_AF^1\text{and}G^{}_BXX_AF^{}.$$
$$\begin{array}{ccc}& \stackrel{F_A}{}& \\ D^b(A)& \stackrel{F^1_A}{}& D^b(A)\\ & & \\ {}_{X_A}{}^{}& & _{X_A}\\ & & \\ D^b(B)& \stackrel{G_B}{}& D^b(B)\\ & \stackrel{G^1_B}{}& \end{array}$$
We will show that $`G^1_BQG^1_BQ,`$ or equivalently
$$QG^{}_BG^1_BQ$$
for any indecomposable projective $`B`$-module $`Q.`$
Let $`Q`$ be an indecomposable projective $`B`$-module and denote by $`T`$ the partial tilting complex of $`A`$ where $`Q=\mathrm{\Phi }(T)=X_AT`$ and denote by $`P`$ the indecomposable projective $`A`$-module where $`P:=F^1_AT.`$ Then we have
$`G^{}_BG^1_BQ`$ $``$ $`G^{}_BG^1_BX_AT`$
$``$ $`G^{}_BX_AF^1_AT`$
$``$ $`G^{}_BX_AP`$
$``$ $`X_AF^{}_AP`$
$``$ $`X_AF_AP\text{by assumption}`$
$``$ $`X_AF_AF^1_AT`$
$``$ $`X_ATQ.`$
Hence, by \[16, Remark 3\], using the identity component $`Out_0(B)`$ of the outer automorphism group of $`B`$ we obtain $`G^1G^1,`$ which implies $`GG^{}.`$ And since $`Out_0`$ is invariant under derived equivalence (see ) we obtain $`FF^{}.`$ $`\mathrm{}`$
## 3 Braid action
Denote by $`B(K_n)`$ the braid group whose associated Coxeter group has Coxeter graph given by the complete graph on $`n`$ vertices, generated by $`a_1,\mathrm{},a_n`$ with relations
$$a_ia_ja_i=a_ja_ia_j$$
for all $`1i,jn,ij.`$ We will define an action of $`B(K_n)`$ on $`D^b(A).`$
###### Theorem 3.1
There is a group homomorphism
$$\phi :B(K_n)TrPic(A)$$
where $`\phi (a_i)=F_i.`$
Proof By Lemma 2.3, it is enough to show that $`_i_j_i(P)_j_i_j(P)`$ for every indecomposable projective module $`P.`$ We will show that for every indecomposable projective module $`P`$
$$_i^1_j_i(P)_j_i_j^1(P).$$
Claim
$$\left(\begin{array}{c}P_1\\ P_2\\ \mathrm{}\\ P_{i1}\\ P_i\\ P_{i+1}\\ \mathrm{}\\ P_{j1}\\ P_j\\ P_{j+1}\\ \mathrm{}\\ P_n\end{array}\right)\stackrel{_i}{}\left(\begin{array}{ccc}P_i& & P_1\\ P_i& & P_2\\ & \mathrm{}& \\ P_i& & P_{i1}\\ P_i& & 0\\ P_i& & P_{i+1}\\ & \mathrm{}& \\ P_i& & P_{j1}\\ P_i& & P_j\\ P_i& & P_{j+1}\\ & \mathrm{}& \\ P_i& & P_n\end{array}\right)\stackrel{_j}{}\left(\begin{array}{ccccc}& & P_i& & P_1\\ & & P_i& & P_2\\ & & & \mathrm{}& \\ & & P_i& & P_{i1}\\ P_j& & P_i& & 0\\ & & P_i& & P_{i+1}\\ & & & \mathrm{}& \\ & & P_i& & P_{j1}\\ & & P_i& & 0\\ & & P_i& & P_{j+1}\\ & & & \mathrm{}& \\ & & P_i& & P_n\end{array}\right)\stackrel{_i^1}{}\left(\begin{array}{cc}& P_1\\ & P_2\\ & \mathrm{}\\ & P_{i1}\\ P_jP_i& P_i\\ & P_{i+1}\\ & \mathrm{}\\ & P_{j1}\\ & P_i\\ & P_{j+1}\\ & \mathrm{}\\ & P_n\end{array}\right)$$
and
$$\left(\begin{array}{c}P_1\\ P_2\\ \mathrm{}\\ P_{i1}\\ P_i\\ P_{i+1}\\ \mathrm{}\\ P_{j1}\\ P_j\\ P_{j+1}\\ \mathrm{}\\ P_n\end{array}\right)\stackrel{_j^1}{}\left(\begin{array}{ccc}P_1& & P_j\\ P_2& & P_j\\ & \mathrm{}& \\ P_{i1}& & P_j\\ P_i& & P_j\\ P_{i+1}& & P_j\\ & \mathrm{}& \\ P_{j1}& & P_j\\ & & P_j\\ P_{j+1}& & P_j\\ & \mathrm{}& \\ P_n& & P_j\end{array}\right)\stackrel{_i}{}\left(\begin{array}{ccccc}& & P_1& & P_j\\ & & P_2& & P_j\\ & & & \mathrm{}& \\ & & P_{i1}& & P_j\\ P_i& & P_i& & P_j\\ & & P_{i+1}& & P_j\\ & & & \mathrm{}& \\ & & P_{j1}& & P_j\\ & & P_i& & P_j\\ & & P_{j+1}& & P_j\\ & & & \mathrm{}& \\ & & P_n& & P_j\end{array}\right)\stackrel{_j}{}\left(\begin{array}{cc}& P_1\\ & P_2\\ & \mathrm{}\\ & P_{i1}\\ P_jP_i& P_i\\ & P_{i+1}\\ & \mathrm{}\\ & P_{j1}\\ & P_i\\ & P_{j+1}\\ & \mathrm{}\\ & P_n\end{array}\right)$$
(2)
Proof of Claim :
For all $`1in`$, denote by $`\delta _i:P_iP_i`$ where $`\delta _i((i))=(i),`$ and for all $`1i,jn`$ by $`\mu _{ij}:P_iP_j,`$ the only map up to a scalars between $`P_i`$ and $`P_j`$ where
$$\mu _{ij}(x)=x(i\mathrm{}j)\text{for all}xP_i.$$
Since $`_j`$ preserves cones, $`_j(P_iP_j)`$ is isomorphic to
$`Cone\left(_j\left(P_i\right)_j\left(P_j\right)\right)`$ $`=`$ $`Cone\left(\begin{array}{cccc}_j\left(P_i\right)=& P_j_jP_AP_i& & A_AP_i\\ & x\left(j\mathrm{}i\right)& & x\left(j\mathrm{}i\right)\\ & & & \\ & & & \\ & & & \\ _j\left(P_j\right)=& P_j_jP_AP_j& & A_AP_j\\ & x\left(j\mathrm{}j\right)& & x\left(j\mathrm{}j\right)\end{array}\right)`$
$``$ $`Cone\left(\begin{array}{ccc}P_j& \stackrel{\mu _{ji}}{}& P_i\\ x& & x\left(j\mathrm{}i\right)\\ _{\delta _j}& & \\ P_j& & \\ x& & \end{array}\right)`$
$``$ $`P_i\left[1\right].`$
For $`kj,`$ $`_j(P_iP_k)`$ is isomorphic to
$$Cone\left(\begin{array}{cccc}_j\left(P_i\right)=& P_j_jP_AP_i& & A_AP_i\\ & x\left(j\mathrm{}i\right)& & x\left(j\mathrm{}i\right)\\ & & & \\ & & & _{\mu _{ik}}\\ & & & \\ _j\left(P_k\right)=& P_j_jP_AP_k& & A_AP_k\\ & x\left(j\mathrm{}k\right)& & x\left(j\mathrm{}k\right)\end{array}\right)Cone\left(\begin{array}{ccc}P_j& \stackrel{\mu _{ji}}{}& P_i\\ x& & x\left(j\mathrm{}i\right)\\ _{\delta _j}& & _{\mu _{ik}}\\ P_j& \stackrel{\mu _{jk}}{}& P_k\\ x& & x\left(j\mathrm{}k\right)\end{array}\right)$$
$``$ $`\begin{array}{ccccc}P_j& \stackrel{\mu _{ji}}{}& P_i& & \\ & {}_{\delta _j}{}^{}& & _{\mu _{ik}}& \\ & & P_j& \stackrel{\mu _{jk}}{}& P_k\end{array}`$
$``$ $`Cone\left(Cone\left(P_jP_iP_j\right)P_k\right)`$
$``$ $`Cone\left(\begin{array}{c}Cone\left(\begin{array}{ccc}(P_j& \stackrel{\mu _{ji}}{}& P_i)\\ _{\delta _j}& & \\ P_j& & \end{array}\right)\left[1\right]\stackrel{(\mu _{ik},\mu _{jk})}{}P_k\end{array}\right)`$
$``$ $`P_i\stackrel{\mu _{ik}}{}P_k.`$
The complex $`_i^1(P_iP_k)`$ ($`ki`$) is isomorphic to
$$Cone\left(\begin{array}{cccc}_i^1\left(P_i\right)=& A_AP_i& & P_i_iP_AP_i\\ & x& & x\left(i\mathrm{}i\right)\left(i\right)+x\left(i\mathrm{}i\right)\\ & & & \\ & & & \\ & & & \\ _i^1\left(P_k\right)=& A_AP_k& & P_i_iP_AP_k\\ & x\left(i\mathrm{}k\right)& & x\left(i\mathrm{}i\right)\left(i\mathrm{}k\right)\end{array}\right)$$
$$Cone\left(\begin{array}{ccc}& & P_i\\ & & _{\delta _i}\\ (P_k& \stackrel{\mu _{ki}}{}& P_i)\end{array}\right)P_k\left(P_i\text{lies in degree}1\right).$$
The complex $`_i^1(P_jP_i0)`$ is isomorphic to
$$Cone\left(\begin{array}{cccc}_i^1\left(P_j\right)=& A_AP_j& & P_i_iP_AP_j\\ & x& & x\left(j\mathrm{}i\right)\left(i\mathrm{}j\right)\\ & & & \\ & & & \\ & & & \\ _i^1\left(P_i\right)=& A_AP_i& & P_i_iP_AP_i\\ & x\left(j\mathrm{}i\right)& & x\left(j\mathrm{}i\right)\left(i\mathrm{}i\right)\end{array}\right)$$
$$Cone\left(\begin{array}{ccc}(P_j& \stackrel{\mu _{ji}}{}& P_i)\\ x& & x\left(j\mathrm{}i\right)\\ & & _{\delta _i}\\ & & P_i\\ & & x\left(j\mathrm{}i\right)\end{array}\right)P_j\stackrel{\mu _{ji}}{}P_i\stackrel{\delta _i}{}P_i,$$
and the complex $`_i(P_iP_j)`$ is isomorphic to
$$Cone\left(\begin{array}{cccc}_i\left(P_i\right)=& P_i_iP_AP_i& & A_AP_i\\ & x\left(i\right)+x\left(i\mathrm{}i\right)& & x\left(i\right)\\ & & & \\ & & & _{\mu _{ij}}\\ & & & \\ _i\left(P_j\right)=& P_i_iP_AP_j& & A_AP_j\\ & x\left(i\mathrm{}j\right)& & x\left(i\mathrm{}j\right)\end{array}\right)$$
$$Cone\left(\begin{array}{ccc}P_i& & \\ x& & \\ {}_{\delta _i}{}^{}& & \\ (P_i& \stackrel{\mu _{ij}}{}& P_j)\\ x& & x\left(i\mathrm{}j\right)\end{array}\right)P_i\stackrel{\delta _i}{}P_i\stackrel{\mu _{ij}}{}P_j$$
And finally, in a similar way, the complex $`_j(P_iP_iP_j)`$ (where $`P_j`$ is in degree $`1`$) is isomorphic to
$`_j\left(Cone\left(P_i\left(Cone\left(P_iP_j\right)\right)\left[1\right]\right)\right)`$ $``$ $`Cone\left(_j\left(P_i\right)\left(_j\left(P_iP_j\right)\right)\left[1\right]\right)`$
$``$ $`Cone\left(\begin{array}{ccc}(P_j& \stackrel{\mu _{ji}}{}& P_i)\\ & & {}_{\delta _i}{}^{}\\ & & P_i\end{array}\right)`$
$``$ $`P_j\stackrel{\mu _{ij}}{}P_i\stackrel{\delta _i}{}P_i.`$
This finishes the proof. $`\mathrm{}`$
## 4 Categorification of a braid group homomorphism
Let $`B`$ be a Brauer tree algebra associated to a line with $`n`$ edges and no exceptional vertex. Denote by $`B(A_n)`$ the Artin braid group on $`n`$ strings generated by $`\sigma _1,\mathrm{},\sigma _n`$ satisfying the relations
$$\sigma _i\sigma _{i+1}\sigma _i=\sigma _{i+1}\sigma _i\sigma _{i+1}\text{and}\sigma _i\sigma _j=\sigma _j\sigma _i\text{if}|ij|>1.$$
In Rouquier and Zimmermann obtained an action of $`B(A_n)`$ on derived category of $`B.`$
###### Theorem 4.1
(Rouquier and Zimmermann )
There is a group homomorphism
$`\psi _n:`$ $`B(A_n)`$ $`TrPic(B)`$
$`\sigma _i`$ $`R_i`$
where $`R_i`$ is the two-sided tilting complex $`0Q_iHom_K(Q_i,K)B0`$ and $`Q_1,\mathrm{}Q_n`$ are the indecomposable projective modules of $`B.`$ Moreover $`\psi _2`$ is injective.
In Khovanov and Seidel independently discovered this homomorphism $`\psi `$ for similar algebras in a very different context.
###### Theorem 4.2
(Khovanov and Seidel )
The map $`\psi _n`$ is injective for all $`n.`$
For simplicity, let us denote by $`\psi :=\psi _n.`$ Denote by $`_i`$ the corresponding functor $`_i:=R_i_B.`$ By \[16, Lemma 4.2\], the images of the indecomposable projective modules $`Q_j`$ of $`B`$ under $`_i`$ are
$$_i\left(Q_j\right)=\{\begin{array}{cccccccc}0& & Q_i& & 0& & & \text{if }i=j\\ 0& & Q_i& & Q_j& & 0& \text{if }\left|ij\right|=1\\ & & 0& & Q_j& & 0& \text{if}\left|ij\right|>1\end{array}$$
where $`Q_j`$ is in degree 0.
In the next subsections we will provide a link between the $`B(K_n)`$, $`B(A_n)`$, the algebras $`A`$ and $`B`$ which will show that the action defined in Theorem 3.1 is unfaithful.
### 4.1 Braid group homomorphism
Define
$$c_n:=\sigma _n\text{and}c_k:=\sigma _k^1c_{k+1}\sigma _k\text{for all}\mathrm{\hspace{0.33em}\hspace{0.33em}1}kn1.$$
Define a mapping
$`\eta :`$ $`B(K_n)`$ $`B(A_n)`$
$`a_i`$ $`c_i`$
###### Proposition 4.3
The mapping $`\eta `$ is a surjective group homomorphism.
First we will need the following Lemma :
###### Lemma 4.4
$`c_j\sigma _i^1=\sigma _i^1c_j`$ for all $`1jn1,\mathrm{\hspace{0.33em}1}in,ij1,ij`$
Proof If $`i>j,`$ we have
$`c_j\sigma _i^1`$ $`=`$ $`\sigma _j^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _j\sigma _i^1`$
$`=`$ $`\sigma _j^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _{i+1}\sigma _i\sigma _{i1}\sigma _i^1\sigma _{i2}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _j^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _{i+1}\sigma _{i1}^1\sigma _i\sigma _{i1}\sigma _{i2}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _j^1\mathrm{}\sigma _{i1}^1\sigma _i^1\sigma _{i1}^1\sigma _{i+1}^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _{i+1}\sigma _i\sigma _{i1}\sigma _{i2}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _j^1\mathrm{}\sigma _i^1\sigma _{i1}^1\sigma _i^1\sigma _{i+1}^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _i^1c_j`$
If $`i<j,`$ $`ij1,`$
$`c_j\sigma _i^1`$ $`=`$ $`\sigma _j^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _j\sigma _i^1`$
$`=`$ $`\sigma _i^1\sigma _j^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _i^1c_j`$
$`\mathrm{}`$
Proof of Proposition 4.3 We need to show
1. $`\sigma _nc_k\sigma _n=c_k\sigma _nc_k`$ for all $`1kn1.`$
2. $`c_ic_jc_i=c_jc_ic_j`$ for all $`1i,jn1,`$ $`ij.`$
3. $`\sigma _i`$ is in the image of $`\eta `$ for all $`1in.`$
1. We will show this using induction. For $`k=n1,`$
$`\sigma _nc_{n1}\sigma _n`$ $`=`$ $`\sigma _n\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _n`$
$`=`$ $`\sigma _n\sigma _n\sigma _{n1}\sigma _n^1\sigma _n`$
$`=`$ $`\sigma _n\sigma _n\sigma _{n1}.`$
$`c_{n1}\sigma _nc_{n1}`$ $`=`$ $`\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _n\sigma _{n1}^1\sigma _n\sigma _{n1}`$
$`=`$ $`\sigma _{n1}^1\sigma _n\sigma _n^1\sigma _{n1}\sigma _n\sigma _n\sigma _{n1}`$
$`=`$ $`\sigma _n\sigma _n\sigma _{n1}.`$
Assume that $`\sigma _nc_j\sigma _n=c_j\sigma _nc_j;`$ we will show that $`\sigma _nc_{j1}\sigma _n=c_{j1}\sigma _nc_{j1}.`$
$`\sigma _nc_{j1}\sigma _n`$ $`=`$ $`\sigma _n\sigma _{j1}^1c_j\sigma _{j1}\sigma _n`$
$`=`$ $`\sigma _{j1}^1\sigma _nc_j\sigma _n\sigma _{j1}`$
$`=`$ $`\sigma _{j1}^1c_j\sigma _nc_j\sigma _{j1}`$
$`=`$ $`\sigma _{j1}^1c_j\sigma _{j1}\sigma _n\sigma _{j1}^1c_j\sigma _{j1}`$
$`=`$ $`c_{j1}\sigma _nc_{j1}.`$
2. For $`i=n1,`$ if $`j=n2,`$
$`c_{n1}c_{n2}c_{n1}`$ $`=`$ $`\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}\sigma _{n1}^1\sigma _n\sigma _{n1}`$
$`=`$ $`\sigma _n\sigma _{n1}\sigma _n^1\sigma _{n2}^1\sigma _n\sigma _{n1}\sigma _n^1\sigma _{n2}\sigma _n\sigma _{n1}\sigma _n^1`$
$`=`$ $`\sigma _n\sigma _{n1}\sigma _{n2}^1\sigma _{n1}\sigma _{n2}\sigma _{n1}\sigma _n^1`$
$`=`$ $`\sigma _n\sigma _{n1}\sigma _{n2}^1\sigma _{n2}\sigma _{n1}\sigma _{n2}\sigma _n^1`$
$`=`$ $`\sigma _n\sigma _{n1}\sigma _{n1}\sigma _{n2}\sigma _n^1.`$
$`c_{n2}c_{n1}c_{n2}`$ $`=`$ $`\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}`$
$`=`$ $`\sigma _{n2}^1\sigma _n\sigma _{n1}\sigma _n^1\sigma _{n2}\sigma _n\sigma _{n1}\sigma _n^1\sigma _{n2}^1\sigma _n\sigma _{n1}\sigma _n^1\sigma _{n2}`$
$`=`$ $`\sigma _{n2}^1\sigma _n\sigma _{n1}\sigma _{n2}\sigma _{n1}\sigma _{n2}^1\sigma _{n1}\sigma _n^1\sigma _{n2}`$
$`=`$ $`\sigma _{n2}^1\sigma _n\sigma _{n2}\sigma _{n1}\sigma _{n2}\sigma _{n2}^1\sigma _{n1}\sigma _{n2}\sigma _n^1`$
$`=`$ $`\sigma _n\sigma _{n1}\sigma _{n1}\sigma _{n2}\sigma _n^1.`$
If $`jn2,`$ $`c_{n1}c_jc_{n1}=`$
$`=`$ $`\sigma _{n1}^1\sigma _n\sigma _{n1}\left(\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n3}^1\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}\sigma _{n3}\mathrm{}\sigma _{j+1}\sigma _j\right)\sigma _{n1}^1\sigma _n\sigma _{n1}`$
$`=`$ $`\left(\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n3}^1\right)\sigma _{n1}^1\sigma _n\sigma _{n1}\left(\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}\right)\sigma _{n1}^1\sigma _n\sigma _{n1}\left(\sigma _{n3}\mathrm{}\sigma _j\right)`$
$`=`$ $`\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n3}^1c_{n1}c_{n2}c_{n1}\sigma _{n3}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n3}^1c_{n2}c_{n1}c_{n2}\sigma _{n3}\mathrm{}\sigma _j,`$
$`c_jc_{n1}c_j=`$
$`=`$ $`\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _j\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n1}^1\sigma _n\sigma _{n1}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n3}^1\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}^1\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _{n2}\sigma _{n3}\mathrm{}\sigma _j`$
$`=`$ $`\sigma _j^1\sigma _{j+1}^1\mathrm{}\sigma _{n3}^1c_{n2}c_{n1}c_{n2}\sigma _{n3}\mathrm{}\sigma _j.`$
Now assume that $`c_jc_kc_j=c_kc_jc_k`$ for all $`k>i.`$ We will show that
$$c_jc_ic_j=c_ic_jc_i.$$
If $`ji+1`$
$`c_jc_ic_j`$ $`=`$ $`c_j\sigma _i^1c_{i+1}\sigma _ic_j`$
$`=`$ $`\sigma _i^1c_jc_{i+1}c_j\sigma _i\text{by Lemma }\text{4.4}`$
$`=`$ $`\sigma _i^1c_{i+1}c_jc_{i+1}\sigma _i`$
$`=`$ $`\sigma _i^1c_{i+1}\sigma _ic_j\sigma _i^1c_{i+1}\sigma _i`$
$`=`$ $`c_ic_jc_i.`$
If $`j=i+1`$
$`c_{i+1}c_ic_{i+1}`$ $`=`$ $`c_{i+1}\sigma _i^1c_{i+1}\sigma _ic_{i+1}`$
$`=`$ $`c_{i+1}\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _ic_{i+1}`$
$`=`$ $`\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _i\sigma _{i+1}^1c_{i+2}\sigma _{i+1}`$
$`=`$ $`\sigma _{i+1}^1c_{i+2}\sigma _i^1\sigma _{i+1}^1\sigma _ic_{i+2}\sigma _i^1\sigma _{i+1}\sigma _ic_{i+2}\sigma _{i+1}`$
$`=`$ $`\sigma _{i+1}^1\sigma _i^1c_{i+2}\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _ic_{i+2}\sigma _{i+1}`$
$`=`$ $`\sigma _{i+1}^1\sigma _i^1c_{i+2}c_{i+1}c_{i+2}\sigma _i\sigma _{i+1}`$
$`=`$ $`\sigma _{i+1}^1\sigma _i^1c_{i+1}c_{i+2}c_{i+1}\sigma _i\sigma _{i+1}`$
$`=`$ $`\sigma _{i+1}^1\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _{i+1}c_{i+2}\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _i\sigma _{i+1}`$
$`=`$ $`\sigma _i^1\sigma _{i+1}^1\sigma _i^1c_{i+2}\sigma _{i+1}c_{i+2}\sigma _{i+1}^1c_{i+2}\sigma _i\sigma _{i+1}\sigma _i`$
$`=`$ $`\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _i^1\sigma _{i+1}c_{i+2}\sigma _{i+1}^1\sigma _ic_{i+2}\sigma _{i+1}\sigma _i`$
$`=`$ $`\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _i^1\sigma _{i+1}\sigma _ic_{i+2}\sigma _i^1\sigma _{i+1}^1\sigma _ic_{i+2}\sigma _{i+1}\sigma _i`$
$`=`$ $`\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _i\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _i^1\sigma _{i+1}^1c_{i+2}\sigma _{i+1}\sigma _i`$
$`=`$ $`c_ic_{i+1}c_i.`$
3. We will show that
$$c_kc_{k+1}c_k^1=c_{k+1}^1c_kc_{k+1}=\sigma _k$$
for all $`1kn1.`$ For $`k=n1,`$
$`c_{n1}c_nc_{n1}^1`$ $`=`$ $`\sigma _{n1}^1\sigma _n\sigma _{n1}\sigma _n\sigma _{n1}^1\sigma _n^1\sigma _{n1}`$
$`=`$ $`\sigma _n\sigma _{n1}\sigma _n^1\sigma _n\sigma _{n1}^1\sigma _n^1\sigma _{n1}`$
$`=`$ $`\sigma _{n1}.`$
Assume $`c_{k+1}c_{k+2}=c_{k+2}\sigma _{k+1}.`$ We will show that $`c_kc_{k+1}=c_{k+1}\sigma _k.`$
$`c_kc_{k+1}`$ $`=`$ $`\sigma _k^1c_{k+1}\sigma _kc_{k+1}`$
$`=`$ $`\sigma _k^1c_{k+2}\sigma _{k+1}c_{k+2}^1\sigma _kc_{k+2}\sigma _{k+1}c_{k+2}^1`$
$`=`$ $`\sigma _k^1c_{k+2}\sigma _{k+1}\sigma _k\sigma _{k+1}c_{k+2}^1`$
$`=`$ $`\sigma _k^1c_{k+2}\sigma _k\sigma _{k+1}\sigma _kc_{k+2}^1`$
$`=`$ $`c_{k+2}\sigma _{k+1}\sigma _kc_{k+2}^1`$
$`=`$ $`c_{k+1}c_{k+2}\sigma _kc_{k+2}^1`$
$`=`$ $`c_{k+1}\sigma _k.`$
$`\mathrm{}`$
###### Proposition 4.5
The homomorphism $`\eta `$ is not injective.
Proof We have $`\sigma _1\sigma _n=\sigma _n\sigma _1`$ but $`a_1a_2a_1^1a_na_na_1a_2a_1^1`$ as we will show using the geometric representation of the Coxeter group of $`B(K_n).`$
Denote by $`W(K_n)`$ the Coxeter group of $`B(K_n).`$ Denote by $`d_i`$ the image of $`a_i`$ in $`W(K_n).`$ By \[7, section 5.3\], $`W(K_n)`$ acts on an $`n`$-dimensional vector space $`V`$ with basis $`\{v_1,\mathrm{},v_n\}.`$
The action of $`W(K_n)`$ is defined by $`d_if_i:VV`$ where
$$f_i(v_i)=v_i\text{and}$$
$$f_i(v_j)=v_i+v_j.$$
Now if $`a_1a_2a_1^1a_na_1a_2^1a_1^1a_n^1=e,`$ then
$$d_1d_2d_1d_nd_1d_2d_1d_n=e\text{and}f_1f_2f_1f_nf_1f_2f_1f_n=id_V;$$
but
$`f_1f_2f_1f_nf_1f_2f_1f_n(v_n)`$ $`=`$ $`f_1f_2f_1f_nf_1f_2f_1(v_n)`$
$`=`$ $`f_1f_2f_1f_nf_1f_2(v_1v_n)`$
$`=`$ $`f_1f_2f_1f_nf_1(2v_2v_1v_n)`$
$`=`$ $`f_1f_2f_1f_n(2v_12v_2v_n)`$
$`=`$ $`f_1f_2f_1(3v_n2v_12v_2)`$
$`=`$ $`f_1f_2(3v_13v_n2v_2)`$
$`=`$ $`f_1(4v_23v_13v_n)`$
$`=`$ $`(4v_14v_23v_n)v_n,\text{a contradiction}.`$
$`\mathrm{}`$
### 4.2 Categorification
Let $`T`$ be the direct sum of the following complexes of projective $`B`$-modules.
$$\begin{array}{ccccccccccccccc}T_1& :& 0& & Q_n& & Q_{n1}& & \mathrm{}& & Q_2& & Q_1& & 0\\ T_2& :& 0& & Q_n& & Q_{n1}& & \mathrm{}& & Q_2& & 0& & \\ \mathrm{}& & & & & & \mathrm{}& & & & & & & & \\ T_{n1}& :& 0& & Q_n& & Q_{n1}& & 0& & & & & & \\ T_n& :& 0& & Q_n& & 0& & & & & & & & \end{array}$$
where $`Q_i`$ is in degree $`i1.`$
By \[13, Theorem 4.2\], using the stable equivalence between $`End_{D^b(B)}(T)`$ and $`A`$ and the result of Gabriel and Riedtmann , $`T`$ is a tilting complex giving a derived equivalence between $`B`$ and $`A.`$ (We note that by computing the endomorphism ring of $`T`$ directly with the method in , we can show that $`End_{D^b(B)}(T)A`$).
Denote by $`G`$ the unique two-sided tilting complex of $`(AB^{op})`$-modules associated to $`T`$ (such a complex is unique up to isomorphism in $`D^b(AB^{op})`$ by a theorem of Keller ). Denote by $`๐ข`$ the corresponding functor
$$๐ข=G_B.$$
It is easy to see that
$$๐ข(T_i)=P_i\text{for all}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in.$$
(9)
We will compute the image of $`R_i`$ under this functor.
$$\begin{array}{ccc}D^b(B)& \stackrel{๐ข}{}& D^b(A)\\ & & \\ {}_{R_i_B}{}^{}& & _{๐ข(R_i)_A}\\ & & \\ D^b(B)& \stackrel{๐ข}{}& D^b(A)\end{array}$$
This will give the image of $`R_i`$ in $`TrPic(A)`$ under the group isomorphism
$$\mathrm{\Gamma }:TrPic(B)\stackrel{}{}TrPic(A)$$
induced by $`๐ข.`$
From now on, when we multiply two elements of $`TrPic,`$ we omit the tensor product symbol $`_A`$ ($`FF^{}:=F_AF^{}`$).
###### Proposition 4.6
We have
$$\mathrm{\Gamma }(R_i)=\{\begin{array}{cc}F_iF_{i+1}F_i^1& \text{if}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,\\ F_n& \text{if}i=n.\end{array}$$
We need the following lemmas
###### Lemma 4.7
For all $`1in1,`$
$$_i\left(T_j\right)=\{\begin{array}{cccc}& T_j& \text{if}& \left|ij\right|>1\\ & T_{i1}& \text{if}& j=i1\\ & T_{i+1}& \text{if}& j=i\\ Cone(Cone(T_iT_{i+1})& T_{i+1})& \text{if}& j=i+1\end{array}$$
Proof For all $`1jn1`$ denote by
$$\begin{array}{ccccccccc}\nu _{j,j+1}:& Q_j& & Q_{j+1}& \text{and}& \nu _{j+1,j}:& Q_{j+1}& & Q_j\\ & x& & x(jj+1)& & & x& & x(j+1j)\end{array}$$
and for all $`1in`$ denote by $`\rho _i:Q_iQ_i`$ where $`\rho _i((i))=(i).`$
For all $`1jn1`$ we have a triangle
$$T_j\stackrel{f_j}{}T_{j+1}\stackrel{g_j}{}Q_j[j]\stackrel{h_j}{}T_j[1]$$
where $`f_j=(\rho _n,\rho _{n1},\mathrm{},\rho _{j+1},0,\mathrm{},0),`$ $`\rho _k`$ in degree $`k1;`$ $`g_j=\nu _{j+1,j}`$ in degree $`j`$ and 0 otherwise; and $`h_j=\rho _j`$ in degree $`j`$ and 0 otherwise.
Therefore,
$$Cone(T_{j+1}[1]Q_j[j1])T_j\text{and}$$
$$Cone(T_jT_{j+1})Q_j[j].$$
For $`j>i+1,`$ we will prove the lemma using induction : if $`j=n,`$
$$_i(T_n)=_i(Q_n[n1])=Q_n[n1]=T_n.$$
Suppose that $`_i(T_{j+1})=T_{j+1}.`$ We will show that $`_i(T_j)=T_j.`$
$`_i\left(T_j\right)`$ $``$ $`_i\left(Cone\left(T_{j+1}\left[1\right]Q_j\left[j1\right]\right)\right)`$ (10)
$``$ $`Cone\left(_i\left(T_{j+1}\right)\left[1\right]_i\left(Q_j\left[j1\right]\right)\right)`$
$``$ $`Cone\left(T_{j+1}\left[1\right]Q_j\left[j1\right]\right)T_j.`$
If $`j=i+1,`$
$`_i\left(T_{i+1}\right)`$ $`=`$ $`_i\left(Cone\left(T_{i+2}\left[1\right]Q_{i+1}\left[i\right]\right)\right)`$
$``$ $`Cone\left(_i\left(T_{i+2}\right)\left[1\right]_i\left(Q_{i+1}\left[i\right]\right)\right)`$
$``$ $`Cone\left(T_{i+2}\left[1\right]\left(Q_iQ_{i+1}\right)\left[i\right]\right)`$
$``$ $`Cone\left(\begin{array}{cccc}(Q_n\mathrm{}& Q_{i+3}& \stackrel{\nu _{i+3,i+2}}{}& Q_{i+2})\\ & _0& & _{\nu _{i+2,i+1}}\\ & (Q_i& \stackrel{\nu _{i,i+1}}{}& Q_{i+1})\end{array}\right)`$
$`\text{where}Q_{i+1}\text{lies in degree}i`$
$``$ $`\left(\begin{array}{ccccc}Q_n\mathrm{}Q_{i+3}& & Q_{i+2}& & \\ & _0& & & \\ & & Q_i& & Q_{i+1}\end{array}\right)`$
$``$ $`Cone\left(\begin{array}{cc}& Q_i\\ & _{\nu _{i,i+1}}\\ (Q_n\mathrm{}Q_{i+2}& Q_{i+1})\end{array}\right)`$
$`\left(\text{since}Hom_B(Q_{i+3},Q_i)=0\right)`$
$``$ $`Cone\left(Q_i\left[i\right]T_{i+1}\right)`$
$``$ $`Cone\left(Cone\left(T_iT_{i+1}\right)T_{i+1}\right)`$
If $`j=i,`$
$`_i\left(T_i\right)`$ $`=`$ $`_i\left(Cone\left(T_{i+1}\left[1\right]Q_i\left[i1\right]\right)\right)`$
$``$ $`Cone\left(_i\left(T_{i+1}\left[1\right]\right)_i\left(Q_i\left[i1\right]\right)\right)`$
$``$ $`Cone\left(\left(\begin{array}{ccccc}Q_n\mathrm{}Q_{i+3}& & Q_{i+2}& & \\ & & & & \\ & & Q_i& & Q_{i+1}\end{array}\right)Q_i\left[i\right]\right)`$
$`\text{where}Q_{i+1}\text{lies in degree}i1`$
$``$ $`\left(\begin{array}{ccccc}Q_n\mathrm{}Q_{i+3}& \stackrel{\nu _{i+3,i+2}}{}& Q_{i+2}& & \\ & {}_{0}{}^{}& & _{\nu _{i+2,i+1}}& \\ & & Q_i& \stackrel{\nu _{i,i+1}}{}& Q_{i+1}\\ & & & _{\rho _i}& \\ & & & & Q_i\end{array}\right)`$
$`\text{where}Q_{i+1}\text{lies in degree}i`$
$``$ $`Cone\left(\left(Q_n\mathrm{}Q_{i+3}Q_{i+2}\right)\stackrel{\nu _{i+2,i+1}}{}\left(\begin{array}{ccc}Q_i& \stackrel{\nu _{i,i+1}}{}& Q_{i+1}\\ & {}_{\rho _i}{}^{}& \\ & & Q_i\end{array}\right)\left[i\right]\right)`$
$``$ $`Cone\left(T_{i+2}\left[1\right]Cone\left(\begin{array}{ccc}Q_i& & Q_{i+1}\\ & & \\ Q_i& & \end{array}\right)\left[i1\right]\right)`$
$``$ $`Cone\left(T_{i+2}\left[1\right]Q_{i+1}\left[i\right]\right)T_{i+1}.`$
For $`j<i,`$ using induction and the same argument as (10) we only need to show that
$$_i(T_{i1})=T_{i1}.$$
We have :
$`_i\left(T_{i1}\right)`$ $`=`$ $`_i\left(Cone\left(T_i\left[1\right]Q_{i1}\left[i2\right]\right)\right)`$
$``$ $`Cone\left(_i\left(T_i\left[1\right]\right)_i\left(Q_{i1}\left[i2\right]\right)\right)`$
$``$ $`Cone\left(T_{i+1}\left[1\right]\left(Q_iQ_{i1}\right)\left[i2\right]\right)`$
$``$ $`Cone\left(\begin{array}{ccc}(Q_n\mathrm{}& Q_{i+1})& \\ & _{\nu _{i+1,i}}& \\ & (Q_i& Q_{i1})\end{array}\right)`$
$`\text{where}Q_{i1}\text{in degree}i2`$
$``$ $`T_{i1}.`$
$`\mathrm{}`$
###### Lemma 4.8
$$_n(T_j)=\{\begin{array}{cc}Cone(T_nT_j)& \text{if}j<n,\\ T_n[1]& \text{if}j=n.\end{array}$$
Proof For $`j=n,`$ $`R_n(T_n)=R_n(Q_n[n1])=Q_n[n]=T_n[1].`$
For $`j<n,`$ we will show the result using induction.
$`_n\left(T_{n1}\right)`$ $`=`$ $`_n\left(Cone\left(Q_nQ_{n1}\right)\right)`$
$``$ $`Cone\left(_n\left(Q_n\right)_n\left(Q_{n1}\right)\right)`$
$``$ $`Cone\left(\begin{array}{ccc}Q_n& & \\ _{\rho _n}& & \\ (Q_n& & Q_{n1})\end{array}\right)`$
$``$ $`Cone\left(T_nT_{n1}\right)`$
Assume that $`_n(T_{j+1})=Cone(T_nT_{j+1}).`$ We will show that
$$_n(T_j)=Cone(T_nT_j).$$
$`_n\left(T_j\right)`$ $``$ $`_n\left(Cone\left(T_{j+1}\left[1\right]Q_j\left[j1\right]\right)\right)`$
$``$ $`Cone\left(_n\left(T_{j+1}\left[1\right]\right)_n\left(Q_j\left[j1\right]\right)\right)`$
$``$ $`Cone\left(Cone\left(T_nT_{j+1}\right)\left[1\right]Q_j\left[j1\right]\right)`$
$``$ $`Cone\left(T_nCone\left(T_{j+1}\left[1\right]Q_j\left[j1\right]\right)\right)`$
$``$ $`Cone\left(T_nT_j\right).`$
$`\mathrm{}`$
Proof of Proposition 4.6 We have a commutative diagram
$`D^b(B)`$ $`\stackrel{๐ข}{}`$ $`D^b(A)`$
$`__i`$ $`_{๐ฒ_i}`$
$`D^b(B)`$ $`\stackrel{๐ข}{}`$ $`D^b(A)`$ (20)
where we denote
$$๐ฒ_i:=๐ข_i๐ข^1.$$
We will show that for all indecomposable projective $`A`$-modules $`P,`$
1. $`๐ฒ_i(P)=_i_{i+1}_i^1(P)`$ for all $`1in1,`$ and
2. $`๐ฒ_n(P)=_n(P).`$
1. Fix $`1in1.`$ Recall that $`๐ข(T_j)=P_j`$ for all $`1jn.`$ By Lemma 4.7 and the above commutative diagram (4.2),
$`๐ฒ_i\left(P_j\right)`$ $`=`$ $`๐ฒ_i\left(๐ข\left(T_j\right)\right)`$
$`=`$ $`๐ข_i\left(T_j\right)`$
$`=`$ $`\{\begin{array}{ccccc}& P_j& \text{if}& \left|ij\right|>1\text{or}& j=i1,\\ & P_{i+1}& \text{if}& j=i,& \\ P_iP_{i+1}& P_{i+1}& \text{if}& j=i+1.& \end{array}`$
From (2) (see Theorem 3.1) we get
$$_i_{i+1}_i^1\left(P_j\right)=\{\begin{array}{ccccc}& P_j& \text{if}& \left|ij\right|>1\text{or}& j=i1,\\ & P_{i+1}& \text{if}& j=i,& \\ P_iP_{i+1}& P_{i+1}& \text{if}& j=i+1.& \end{array}$$
(22)
Therefore, we get $`๐ฒ_i=_i_{i+1}_i^1.`$
2. By (9), Lemma 4.8 and Diagram (4.2),
$`๐ฒ_n\left(P_j\right)`$ $`=`$ $`๐ฒ_n๐ข\left(T_j\right)`$
$`=`$ $`๐ข_n\left(T_j\right)`$
$`=`$ $`๐ข\left(Cone\left(T_nT_j\right)\right)`$
$`=`$ $`P_nP_j\text{for all}\mathrm{\hspace{0.33em}\hspace{0.33em}1}j<n,`$
and $`๐ฒ_n(P_n)=๐ฒ_n๐ข_n(T_n)=๐ข_n(T_n)=๐ข(T_n[1])=P_n[1].`$ Hence, by Lemma 2.2, $`๐ฒ_n=_n.`$ $`\mathrm{}`$
###### Proposition 4.9
$`\mathrm{\Gamma }^1(F_n)=R_n`$ and
$`\mathrm{\Gamma }^1(F_k)`$ $`=`$ $`R_k^1R_{k+1}^1\mathrm{}R_{n1}^1R_nR_{n1}\mathrm{}R_{k+1}R_k`$
$`=`$ $`R_nR_{n1}\mathrm{}R_{k+1}R_kR_{k+1}^1\mathrm{}R_{n1}^1R_n^1`$
Proof Define $`C_n:=R_n`$ and
$$C_k:=R_k^1C_{k+1}R_k,\text{for all}\mathrm{\hspace{0.33em}\hspace{0.33em}1}kn1.$$
By Proposition 4.6, $`\mathrm{\Gamma }(C_n)=F_n.`$ Now assume that $`\mathrm{\Gamma }(C_{k+1})=F_{k+1}.`$ We will show that $`\mathrm{\Gamma }(C_k)=F_k.`$
$`\mathrm{\Gamma }(C_k)`$ $`=`$ $`\mathrm{\Gamma }(R_k^1C_{k+1}R_k)`$
$`=`$ $`(F_kF_{k+1}^1F_k^1)F_{k+1}(F_kF_{k+1}F_k^1)`$
$`=`$ $`F_kF_{k+1}^1F_k^1F_kF_{k+1}F_kF_k^1=F_k.`$
$`\mathrm{}`$
###### Corollary 4.10
The following diagram is commutative :
$$\begin{array}{ccc}B(K_n)& \stackrel{\eta }{}& B(A_n)\\ & & \\ _\phi & & _\psi \\ & & \\ TrPic(A)& \stackrel{\mathrm{\Gamma }^1}{}& TrPic(B)\end{array}$$
This defines a categorification of the group homomorphism $`\eta .`$
###### Corollary 4.11
Since $`\eta `$ is not injective, the action of $`B(K_n)`$ on $`D^b(A)`$ is not faithful.
## 5 A faithful action
Denote by $`B(\stackrel{~}{A}_{n1})`$ the affine braid group generated by $`h_1,\mathrm{},h_n`$ subject to relations
$`h_ih_{i+1}h_i`$ $`=`$ $`h_{i+1}h_ih_{i+1}\text{if}\mathrm{\hspace{0.33em}1}in1,`$
$`h_ih_j`$ $`=`$ $`h_jh_i\text{if}|ij|1,n1\text{and}`$
$`h_nh_1h_n`$ $`=`$ $`h_1h_nh_1.`$
In Schaps and Zakay-Illouz obtained a group homomorphism between $`B(\stackrel{~}{A}_{n1})`$ and the subgroup of $`TrPic(A)`$ generated by $`H_1,\mathrm{},H_n`$ where $`H_i`$ is the refolded tilting complex associated to the transposition $`(ii+1).`$ (See for the definition).
###### Theorem 5.1
(Schaps and Zakay-Illouz )
There is a group homomorphism
$`\rho :`$ $`B(\stackrel{~}{A}_{n1})`$ $`TrPic(A)`$
$`h_i`$ $`H_i.`$
Remark In the action is defined for Brauer star algebras in general.
Denote by $`_i`$ the corresponding functor $`_i:=H_i_A.`$
Using \[17, Proposition 1\] we get the images of indecomposable projective modules of $`A`$ under $`_i`$ :
$$_i(P_j)=\{\begin{array}{ccccccc}& & & & P_j& \text{if}& |ij|>1,\\ & & & & P_{i1}& \text{if}& j=i1,\\ & & & & P_{i+1}& \text{if}& j=i,\\ P_i& & P_{i+1}& & P_{i+1}& \text{if}& j=i+1.\end{array}$$
Comparing this with (22) we get
$$_i=_i_{i+1}_i^1\text{for all}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,$$
(23)
and it is easy to see from (2) that
$$_n=_n_1_n^1.$$
(24)
### 5.1 Injective braid group homomorphisms
Let $`B(B_n)`$ be the braid group associated to Dynkin diagram of type $`B_n`$ with generators $`b_1,b_2,\mathrm{},b_n`$ subject to relations :
$`b_ib_{i+1}b_i`$ $`=`$ $`b_{i+1}b_ib_{i+1}\text{if}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`b_ib_j`$ $`=`$ $`b_jb_i\text{if}|ij|>1,`$
$`b_{n1}b_nb_{n1}b_n`$ $`=`$ $`b_nb_{n1}b_nb_{n1}.`$
By , $`B(B_n)`$ is also generated by $`\tau ,s_1,s_2,\mathrm{},s_n`$ subject to the relations :
$`s_1s_ns_1`$ $`=`$ $`s_ns_1s_n,`$
$`s_is_{i+1}s_i`$ $`=`$ $`s_{i+1}s_is_{i+1}\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`s_is_j`$ $`=`$ $`s_js_i\text{for}|ij|1,n1,`$
$`\tau s_i\tau ^1`$ $`=`$ $`s_{i1}\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`\tau s_1\tau ^1`$ $`=`$ $`s_n,`$
and the subgroup of $`B(B_n)`$ generated by $`s_1,\mathrm{},s_n`$ is the affine braid group $`B(\stackrel{~}{A}_{n1}).`$ By , this shows that $`B(B_n)`$ is a semidirect product of the infinite cyclic group generated by $`\tau `$ and $`B(\stackrel{~}{A}_{n1}).`$ Therefore $`B(\stackrel{~}{A}_{n1})`$ injects into $`B(B_n)`$ (see for more details).
We can reformulate \[5, Proposition 6\] to get
$`\tau `$ $`=`$ $`b_nb_{n1}\mathrm{}b_2,`$
$`s_i`$ $`=`$ $`b_i\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`s_n`$ $`=`$ $`\tau b_1\tau ^1=b_nb_{n1}\mathrm{}b_2b_1b_2^1\mathrm{}b_{n1}^1b_n^1,`$
such that the embedding of $`B(\stackrel{~}{A}_{n1})`$ into $`B(B_n)`$ is defined by the following map :
$`\mu :`$ $`B(\stackrel{~}{A}_{n1})`$ $`B(B_n)`$
$`h_i`$ $`b_i\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`h_n`$ $`b_nb_{n1}\mathrm{}b_2b_1b_2^1\mathrm{}b_{n1}^1b_n^1.`$
Now define a group homomorphism
$`\chi :`$ $`B(B_n)`$ $`B(A_n)`$
$`b_i`$ $`\sigma _i\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`b_n`$ $`\sigma _n^2.`$
By , $`\chi `$ is injective. Therefore we obtain an injective composition of injective group homomorphisms
$`\chi \mu :`$ $`B(\stackrel{~}{A}_{n1})`$ $`B(A_n)`$
$`h_i`$ $`\sigma _i,\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`h_n`$ $`\sigma _n\sigma _n\sigma _{n1}\mathrm{}\sigma _2\sigma _1\sigma _2^1\mathrm{}\sigma _{n1}^1\sigma _n^1\sigma _n^1.`$
### 5.2 Another Categorification
By Proposition 4.6, Proposition 4.9, (23) and (24) we obtain
$`\mathrm{\Gamma }^1(H_i)`$ $`=`$ $`R_i\text{for all}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`\mathrm{\Gamma }^1(H_n)`$ $`=`$ $`\mathrm{\Gamma }^1(F_nF_1F_n^1)`$
$`=`$ $`R_nR_nR_{n1}\mathrm{}R_2R_1R_2^1\mathrm{}R_{n1}^1R_n^1R_n^1.`$
Therefore, we get the following commutative diagram :
$$\begin{array}{ccc}B(\stackrel{~}{A}_{n1})& \stackrel{\chi \mu }{}& B(A_n)\\ _\rho & & _\psi \\ TrPic(A)& \stackrel{\mathrm{\Gamma }^1}{}& TrPic(B)\end{array}$$
And since $`\chi \mu `$ and $`\psi `$ are injective, we obtain the following result :
###### Theorem 5.2
The action of $`B(\stackrel{~}{A}_{n1})`$ on $`D^b(A)`$ defined in Theorem 5.1 is faithful.
Remark We leave as an open question whether the action defined in is faithful for Brauer star algebras with exceptional vertex. We also get a faithful action of $`B(B_n)`$ on $`D^b(A)`$ by defining
$`B(B_n)`$ $``$ $`TrPic(A),`$
$`b_i`$ $``$ $`F_iF_{i+1}F_i\text{if}\mathrm{\hspace{0.33em}\hspace{0.33em}1}in1,`$
$`b_n`$ $``$ $`F_nF_n.`$
Intan Muchtadi-Alamsyah
Department of Mathematics,
University of Leicester,
University Road,
LE1 7RH, Leicester, UK,
E-mail address : idma1@le.ac.uk
Current address :
Institut Teknologi Bandung,
Faculty of Mathematics and Natural Sciences,
Algebra Research Group,
Jl.Ganesha no.10,
Bandung 40132, Indonesia,
E-mail address : ntan@math.itb.ac.id |
warning/0506/math0506546.html | ar5iv | text | # Yang-Baxter bases of 0-Hecke algebras and representation theory of 0-Ariki-Koike-Shoji algebras
## 1. Introduction
Given an *inductive tower of algebras*, that is, a sequence of algebras
(1)
$$A_0A_1\mathrm{}A_n\mathrm{},$$
with embeddings $`A_mA_nA_{m+n}`$ satisfying appropriate compatibility conditions, one can introduce two *Grothendieck rings*
(2)
$$๐ข:=\underset{n0}{}G_0(A_n),๐ฆ:=\underset{n0}{}K_0(A_n),$$
where $`G_0(A)`$ and $`K_0(A)`$ are the (complexified) Grothendieck groups of the categories of finite-dimensional $`A`$-modules and projective $`A`$-modules respectively, with multiplication of the classes of an $`A_m`$-module $`M`$ and an $`A_n`$-module $`N`$ defined by
(3)
$$[M][N]=[M\widehat{}N]=[MN_{A_mA_n}^{A_{m+n}}].$$
On each of these Grothendieck rings, one can define a coproduct by means of restriction of representations. Under favorable circumstances, this turns these two rings into mutually dual Hopf algebras.
The basic example of this situation is the character ring of symmetric groups (over $``$), due to Frobenius. Here the $`A_n=๐_n`$ are semi-simple algebras, so that
(4)
$$G_0(A_n)=K_0(A_n)=R(A_n),$$
where $`R(A)`$ denotes the vector space spanned by isomorphism classes of indecomposable modules which in this case are all simple and projective. The irreducible representations $`[\lambda ]`$ of $`A_n`$ are parametrized by partitions $`\lambda `$ of $`n`$, and the Grothendieck ring is isomorphic to the algebra $`Sym`$ of symmetric functions under the correspondence $`[\lambda ]s_\lambda `$, where $`s_\lambda `$ denotes the Schur function associated with $`\lambda `$.
Other known examples with towers of group algebras over the complex numbers, $`A_n=G_n`$, include the cases of wreath products $`G_n=\mathrm{\Gamma }๐_n`$ (Specht), finite linear groups $`G_n=GL(n,๐
_q)`$ (Green), *etc.*, all related to symmetric functions (see ).
Examples involving modular representations of finite groups and non-semisimple specializations of Hecke algebras have also been worked out (see ). For example, finite Hecke algebras of type $`A`$ at roots of unity ($`A_n=H_n(\zeta )`$, $`\zeta ^k=1`$) yield quotients and subalgebras of $`Sym`$
(5)
$$๐ข=Sym/(p_{km}=0),๐ฆ=\left[p_i|i0(modk)\right],$$
where the $`p_i`$ are the power sum symmetric functions, supporting level $`1`$ irreducible representations of the affine Lie algebra $`\widehat{๐ฐ๐ฉ}_k`$, while Ariki-Koike algebras at roots of unity give rise to higher level representations of the same Lie algebras . The $`0`$-Hecke algebras $`A_n=H_n(0)`$ correspond to the pair Quasi-symmetric functions - Noncommutative symmetric functions, $`๐ข=\mathrm{๐๐๐ฆ๐}`$, $`๐ฆ=\mathrm{๐๐ฒ๐ฆ}`$ . Affine Hecke algebras at roots of unity lead to $`U^+(\widehat{๐ฐ๐ฉ}_k)`$ and $`U^+(\widehat{๐ฐ๐ฉ}_k)^{}`$ , and the cases of generic affine Hecke algebras can be reduced to a subcategory admitting as Grothendieck rings $`U^+(\widehat{๐ค๐ฉ}_{\mathrm{}})`$ and $`U^+(\widehat{๐ค๐ฉ}_{\mathrm{}})^{}`$ .
A further interesting example is the tower of $`0`$-Hecke-Clifford algebras , giving rise to the peak algebras .
Here, we shall show that appropriate versions at $`q=0`$ of the Ariki-Koike algebras derived from the presentation of Shoji admit as Grothendieck rings two known combinatorial Hopf algebras, the Mantaci-Reutenauer descent algebras (associated with the corresponding wreath products) , and their duals, a generalization of quasi-symmetric functions, introduced by Poirier in and more recently considered in . This work can be understood as the solution of an *inverse problem* in representation theory: given the Grothendieck rings, reconstruct the algebras. The main point is that the standard presentations of Ariki-Koike algebras at $`q=0`$ do not give the required Grothendieck rings.
This article is structured as follows. After recalling some classical definitions, we introduce new combinatorial structures (cycloribbons and anticycloribbons) needed in the sequel (Section 2). Next, we study a limit case of the Yang-Baxter bases of Hecke algebras introduced in , which will be one of our main tools (Section 3). In Section 4, we finally introduce the fundamental objects of our study, the $`0`$-Ariki-Koike-Shoji algebras $`_{n,r}(0)`$, and special bases, well suited for analyzing representations. Next, we obtain the classification of simple $`_{n,r}(0)`$-modules, which turn out to be all one-dimensional, labelled by the $`r(r+1)^{n1}`$ cycloribbons. We then describe induction products and restrictions of these simple modules, which allows us to identify the first Grothendieck ring $`๐ข`$ with a Hopf subalgebra of Poirierโs Quasi-symmetric functions, dual to the Mantaci-Reutenauer Hopf algebra (Section 5). Since $`_{n,r}(0)`$ is self-injective, this duality gives immediately the Grothendieck ring $`๐ฆ`$ associated with projective modules. An alternative labelling of the indecomposable projective modules leads then to a simple description of basic operations such as induction products, restriction to $`H_n(0)`$, or induction from a projective $`H_n(0)`$-module. Summarizing, we obtain an explicit description of the Cartan-Brauer triangle, in particular of the Cartan invariants and of the decomposition matrices (Section 6). We conclude with a description of the Ext-quiver of $`_{n,r}(0)`$. The last section contains tables of $`q`$-Cartan invariants (recording the radical filtrations of projective modules) and decomposition matrices.
## 2. Notations and preliminaries
### 2.1. Words and permutations
Let $`w`$ be a word on a totally ordered alphabet $`A`$. We denote by $`\overline{w}`$ the *mirror image* (reading from right to left) of $`w`$. If the alphabet is $`^{}`$, the *evaluation* of $`w`$ is the sequence of number of occurrences of $`1`$, $`2`$, and so on until all letters of $`w`$ have been seen. For example, the evaluation of $`15423341511457`$ is $`(4,1,2,3,3,0,1)`$.
We will generally use the operation of concatenation on words but we will also use the *shifted concatenation* $`uv`$ of two words $`u`$ and $`v`$ over the positive integers consisting in concatenating $`u`$ with $`v[k]:=(v_1+k,\mathrm{},v_l+k)`$, where $`k=|u|`$ is the length of $`u`$. For example $`1431232=1431676`$.
We shall also make use of another classical operation on words, known as *standardization*. The standardized word $`\mathrm{Std}(w)`$ of $`w`$ is the permutation having the same inversions as $`w`$, and its right standardized word is $`\mathrm{RStd}(w):=\overline{\mathrm{Std}(\overline{w})}`$. For example, $`\mathrm{Std}(412461415)=514692738`$, $`\mathrm{RStd}(412461415)=734692518`$.
We will represent permutations as words of size $`n`$ over the alphabet $`\{1,\mathrm{},n\}`$. Denote by $`s_i`$ the transposition exchanging $`i`$ and $`i+1`$, and by $`\omega _n`$ the maximal permutation $`(n,n1,\mathrm{},1)`$. For $`\sigma =(\sigma _1,\mathrm{},\sigma _n)`$, we set
(6)
$$\overline{\sigma }:=(\sigma _n,\mathrm{},\sigma _1)=\sigma \omega _n,\sigma ^\mathrm{\#}:=(n+1\sigma _1,\mathrm{},n+1\sigma _n)=\omega _n\sigma .$$
### 2.2. Compositions
Recall that a *composition* of an integer $`n`$ is any finite sequence of positive integers $`I=(i_1,\mathrm{},i_k)`$ of sum $`|I|:=i_1+\mathrm{}+i_k=n`$. It can be pictured as a ribbon diagram, that is, a set of rows composed of square cells of respective lengths $`i_j`$, the first cell of each row being attached under the last cell of the previous one. $`I`$ is called the *shape* of the ribbon diagram. The *conjugate* $`I^{}`$ of a composition $`I`$ is the composition built with the number of cells of the columns of its ribbon diagram read from *right to left*. Its *mirror image* $`\overline{I}`$ is the reading of $`I`$ from right to left. For example, if $`I=(2,3,1,2)`$, $`I^{}`$ is $`(1,3,1,2,1)`$ and $`\overline{I}`$ is $`(2,1,3,2)`$.
Given a filling of a ribbon diagram, we define its *row reading* as the reading of its rows from left to right and from top to bottom. Finally, the *descent set* $`\mathrm{D}(I)`$ of a composition $`I=(i_1,\mathrm{},i_l)`$ is the set $`\{i_1,i_1+i_2,\mathrm{},i_1+\mathrm{}+i_{l1}\}`$ composed of the *descents* of $`I`$. Recall also that the *descent set* $`\mathrm{D}(\sigma )`$ of a permutation $`\sigma `$ is the set of $`i`$ such that $`\sigma (i)>\sigma (i+1)`$ (the *descents* of $`\sigma `$), and the *descent composition* $`\mathrm{C}(\sigma )`$ of $`\sigma `$ is the unique composition $`I`$ of $`n`$ such that $`\mathrm{D}(I)=\mathrm{D}(\sigma )`$, that is, the shape of a filled ribbon diagram whose row reading is $`\sigma `$ and whose rows are increasing and columns decreasing. For example, Figure 2 shows that the descent composition of $`(3,5,4,1,2,7,6)`$ is $`I=(2,1,3,1)`$.
Conversely, with a composition $`I`$, associate its *maximal permutation* $`\sigma =\omega (I)`$ as the permutation with descent composition $`I`$ and maximal inversion number. Similarly, the minimal permutation $`\alpha (I)`$ is the permutation with descent composition $`I`$ and minimal inversion number. For example, if $`I=(2,1,3,1)`$, $`\omega (I)=6752341`$ and $`\alpha (I)=1432576`$.
### 2.3. Shuffle, shifted shuffle, and convolution
The *shuffle product* $`\text{ }\text{ }`$ of two words $`u`$ and $`v`$ over an alphabet $`A`$ is inductively defined by:
(7)
$$\text{if }u=au^{}\text{ and }v=bv^{}\text{ with }a,bA\text{,\hspace{1em}then }u\text{ }\text{ }v=a(u^{}\text{ }\text{ }v)+b(u\text{ }\text{ }v^{}),$$
with the initial conditions $`u\text{ }\text{ }\epsilon =\epsilon \text{ }\text{ }u=u`$, $`\epsilon `$ being the empty word.
The *shifted shuffle* $``$ of two permutations $`\sigma ^{}๐_k`$ and $`\sigma ^{\prime \prime }๐_l`$ is the shuffle of $`\sigma ^{}`$ and $`\sigma ^{\prime \prime }[k]`$. For example,
(8)
$$2112=2134+2314+3214+2341+3241+3421.$$
The *convolution product* $``$ of two permutations $`\sigma ^{}๐_k`$ and $`\sigma ^{\prime \prime }๐_l`$ is the sum of permutations $`\sigma `$ in $`๐_{k+l}`$ such that $`\mathrm{Std}(\sigma _1\mathrm{}\sigma _k)=\sigma ^{}`$ and $`\mathrm{Std}(\sigma _{k+1}\mathrm{}\sigma _{k+l})=\sigma ^{\prime \prime }`$. Equivalently, it is the sum obtained by inverting the permutations occuring in the shifted shuffle of the inverses of $`\sigma ^{}`$ and $`\sigma ^{\prime \prime }`$. It is well-known that the convolution of two permutations is an interval of the (left) weak order. For example, one has
(9)
$$2112=2134+3124+3214+4123+4213+4312=[2134,4312].$$
In all the paper, we make use of a set $`C=\{1,\mathrm{},r\}`$, called the *color set*. The color words will be written in bold type.
The previous definitions can be extended to colored permutations as in , that is, pairs $`(\sigma ,๐)๐_n\times C^n`$. The *shifted shuffle* of two such elements $`(\sigma ^{},๐^{})(\sigma ^{\prime \prime },๐^{\prime \prime })`$ is the sum of colored permutations obtained by shuffling with shift the two permutations, the colors remaining attached to letters. For example, representing colors as exponents,
(10)
$$\begin{array}{cc}\hfill 2^21^11^12^2& =2^21^1\text{}\text{}\mathrm{\hspace{0.17em}3}^14^2\hfill \\ & =2^21^13^14^2+2^23^11^14^2+3^12^21^14^2+2^23^14^21^1+3^12^24^21^1+3^14^22^21^1.\hfill \end{array}$$
The *convolution* $`(\sigma ^{},๐^{})(\sigma ^{\prime \prime },๐^{\prime \prime })`$ of two colored permutations is the sum of colored permutations $`(\sigma ,๐^{}๐^{\prime \prime })`$, where $`\sigma `$ runs over $`\sigma ^{}\sigma ^{\prime \prime }`$. Equivalently, if one defines the *inverse* of a colored permutation $`(\sigma ,๐)`$ as $`(\sigma ^1,๐\sigma ^1)`$, the convolution of two colored permutations is the sum of the inverses of the elements occuring in the shifted shuffle of their inverses. For example,
(11)
$$2^21^11^12^2=2^21^13^14^2+3^21^12^14^2+3^22^11^14^2+4^21^12^13^2+4^22^11^13^2+4^23^11^12^2.$$
The sums being multiplicity free, we can regard these as sets, and write, *e.g.*, $`23142112`$, or $`3^21^12^14^22^21^11^12^2`$.
### 2.4. Cycloribbons and anticycloribbons
Let $`I`$ be a composition of $`n`$ and $`๐C^n`$ be a color word of length $`n`$. The pair $`[I,๐]`$ is called a *colored ribbon*, and depicted as the filling of $`I`$ whose row reading is $`๐`$. We say that this filling is a *cyclotomic ribbon* (*cycloribbon* for short) if it is *weakly increasing* in rows and *weakly decreasing* in columns. Notice that there are $`r(r+1)^{n1}`$ cycloribbons with at most $`r`$ colors since when building the ribbon cell by cell, one has $`r`$ possibilities for its first cell and then $`r+1`$ possibilities for the next ones: $`1`$ possibility for the $`r1`$ choices different from the previous one and $`2`$ possibilities for the same choice (right or down). Here are the five cycloribbons of shape $`(2,1)`$ with two colors:
(12) $`1`$ $`1`$ $`1`$ $`1`$ $`2`$ $`1`$ $`1`$ $`2`$ $`2`$ $`2`$ $`2`$ $`1`$ $`2`$ $`2`$ $`2`$
We say that a colored ribbon is an *anticyclotomic ribbon* (*anticycloribbon* for short) if it is *weakly decreasing* in rows and *weakly increasing* in columns. There are as many anticycloribbons as cycloribbons. The relevant bijection $`\varphi `$ from one set to the other is the restriction of an involution on all colored ribbons: read a ribbon $`R`$ row-wise and build its image $`\varphi (R)`$ cell by cell as follows:
* if the $`(i+1)`$-th cell has the same content as the $`i`$-th cell, glue it at the same position as in $`R`$ (right or down),
* if the $`(i+1)`$-th cell does not have the same content as the $`i`$-th cell, glue it at the other position (right or down).
For example, the colored ribbons of Figure 3 are exchanged by $`\varphi `$:
Let us now extend to colored permutations and anticycloribbons the correspondances between permutations and ribbons: we say that $`i`$ is an *anti-descent* of $`(\sigma ,๐)`$ iff $`c_i<c_{i+1}`$ or, $`c_i=c_{i+1}`$ and $`\sigma _i<\sigma _{i+1}`$. Given a colored permutation $`(\sigma ,๐)`$, define its associated anticycloribbon as the pair $`[I,๐]`$ where $`I`$ has a descent at the $`i`$-th position iff $`i`$ is an anti-descent of the colored permutation. Conversely, with a given anticycloribbon $`[I,๐]`$, associate its *maximal colored permutation* $`(\sigma ,๐)`$ where $`\sigma =\omega (J)`$, where $`\mathrm{D}(J)=\{i|c_ic_{i+1}\text{, or }c_i=c_{i+1}\text{ and }i\mathrm{D}(I)\}`$.
For example, the anti-descents of $`(81732564,11132221)`$ are $`\{2,3,5,6\}`$. Its associated anticycloribbon is
(13) $`1`$ $`1`$ $`1`$ $`2`$ $`2`$ 1 $`2`$
and the maximal colored permutation having this associated anti-cycloribbon is $`\omega (1,2,1,3,1)=(86752341,11132221)`$.
## 3. The Hecke algebra of type $`A`$ at $`q=0`$ and Yang-Baxter bases
In this section, we first recall some combinatorial aspects of the representation theory of $`H_n(0)`$, such as a realization of its simple and indecomposable projective modules in the left regular representation and the calculation of the composition factors of an induction product of two simple or indecomposable projective modules. We will show in the next Section how these constructions can be generalized to the $`0`$-Ariki-Koike-Shoji algebra $`_{n,r}(0)`$. To achieve this, we will need to investigate in some detail a special case of the so-called Yang-Baxter bases of $`H_n(0)`$.
### 3.1. Representation theory
Let us consider the Hecke algebra $`H_n(0)`$ of type $`A`$ at $`q=0`$ in the presentation of , that is, the $``$-algebra with generators $`T_i`$, $`1in1`$ and relations:
(14) $`T_i(1+T_i)`$ $`=0`$ $`(1in1),`$
(15) $`T_iT_{i+1}T_i`$ $`=T_{i+1}T_iT_{i+1}`$ $`(1in2),`$
(16) $`T_iT_j`$ $`=T_jT_i`$ $`(|ij|2).`$
Let $`\sigma =:\sigma _{i_1}\mathrm{}\sigma _{i_p}`$ be a reduced word for a permutation $`\sigma ๐_n`$. The defining relations of $`H_n(0)`$ ensure that the element $`T_\sigma :=T_{i_1}\mathrm{}T_{i_p}`$ is independent of the chosen reduced word for $`\sigma `$. Moreover, the well-defined family $`(T_\sigma )_{\sigma ๐_n}`$ is a basis of the Hecke algebra, which is consequently of dimension $`n!`$.
#### 3.1.1. Simple modules and induction
It is known that $`H_n(0)`$ has $`2^{n1}`$ simple modules, all one-dimensional, naturally labelled by compositions $`I`$ of $`n`$ : following the notation of , let $`\eta _I`$ be the generator of the simple $`H_n(0)`$-module $`S_I`$ associated with $`I`$ in the left regular representation. It satisfies
(17)
$$\{\begin{array}{ccc}\hfill T_i\eta _I& =\eta _I& \hfill \text{if }i\mathrm{D}(I)\text{,}\\ \hfill (1+T_i)\eta _I& =\eta _I& \hfill \text{otherwise.}\end{array}$$
The composition factors of the induction product $`S_I\widehat{}S_J`$ of two simple $`0`$-Hecke modules are easily described in terms of permutations.
Let us say that a family $`(b)`$ of elements of a $`H_n(0)`$-module $`M`$ is a *combinatorial basis* if it is a basis of $`M`$ such that $`T_ib`$ is either $`b`$, or $`0`$ or another basis element $`b^{}`$ for all $`i`$. Then $`M`$ is completely encoded by the edge-labelled directed graph having as vertices the basis elements, with a loop labelled $`T_i`$ iff $`T_ib=b`$, a loop labelled $`1+T_i`$ iff $`(1+T_i)b=b`$, and an edge from $`b`$ to $`b^{}`$ labelled $`T_i`$ iff $`T_ib_\tau `$ is $`b_{s_i\tau }`$ (if $`s_i\tau `$ occurs in $`\sigma \sigma ^{}`$, and $`s_i\tau `$ has one inversion more than $`\tau `$) or $`b_\tau `$ or $`0`$, depending on whether $`i`$ is a descent of $`\tau ^1`$ or not. The composition factors are then the simple modules indexed by the descent compositions of the inverses of the permutations occuring in the convolution (see Figure 4).
This can be interpreted as a computation in the algebra $`\mathrm{๐
๐๐๐ฒ๐ฆ}^{}`$ , where the product of two basis elements $`๐_\alpha ๐_\beta `$ is
(18)
$$๐_\alpha ๐_\beta =\underset{\gamma \alpha \beta }{}๐_\gamma .$$
Indeed, the composition factors of the induction product of two simple $`0`$-Hecke modules can be obtained by computing the product $`๐_\sigma ๐_\sigma ^{}`$ and taking the image of the result by the morphism sending $`๐_\tau `$ to the quasi-symmetric function $`F_{\mathrm{C}(\tau ^1)}\mathrm{๐๐๐ฆ๐}`$. It is known that this amounts to take the commutative image of the realization of $`\mathrm{๐
๐๐๐ฒ๐ฆ}^{}`$ by noncommutative polynomials . Since this is an epimorphism, the Grothendieck ring $`๐ข`$ of the tower of the $`0`$-Hecke algebras is isomorphic to $`\mathrm{๐๐๐ฆ๐}`$. We shall later use a similar argument to identify the Grothendieck ring of the $`0`$-Ariki-Koike-Shoji algebras.
#### 3.1.2. Indecomposable projective modules
The indecomposable projective $`H_n(0)`$-module $`P_I`$ labelled by a composition $`I`$ of $`n`$ is combinatorial in a certain basis $`(g_\sigma )`$. Its $`0`$-Hecke graph coincides with the restriction of the left weak order to the interval $`[\alpha (I),\omega (I)]`$, that is the set of permutations with descent composition $`I`$, an edge $`\sigma _i\tau `$ meaning that $`T_i(g_\sigma )=g_\tau `$. The basis $`g_\sigma `$ is defined in terms of Nortonโs generators:
(19)
$$g_{\alpha (I)}=\nu _I=T_{\alpha (I)}T_{\alpha (\overline{I}^{})}^{},$$
where $`T_i^{}=1+T_i`$ satisfies the braid relations, so that $`T_\sigma ^{}`$ is well defined for any permutation. Figure 5 shows the projective module $`P_{121}`$ of $`H_4(0)`$.
### 3.2. Yang-Baxter bases of $`H_n(0)`$
It is known that from a set of generators (depending on two parameters $`t`$ and $`u`$)
(20)
$$Y_i(t,u)=a(t,u)+b(t,u)T_ib0,$$
of $`H_n(0)`$ satisfying the quantum Yang-Baxter equation
(21)
$$Y_i(t,u)Y_{i+1}(t,v)Y_i(u,v)=Y_{i+1}(u,v)Y_i(t,v)Y_{i+1}(t,u),$$
one can associate with any vector $`๐ฑ=(x_1,\mathrm{},x_n)`$ of โspectral parametersโ, a Yang-Baxter basis $`(Y_\sigma (๐ฑ))_{\sigma ๐_n}`$ of $`H_n(0)`$, inductively defined by $`Y_{id}(๐ฑ)=1`$ and
(22)
$$Y_{\sigma _j\tau }(๐ฑ)=Y_j(x_{\tau ^1(j)},x_{\tau ^1(j+1)})Y_\tau (๐ฑ)=Y_{\sigma _j}(๐ฑ\tau ^1)Y_\tau (๐ฑ),$$
if $`\sigma _j\tau `$ has one inversion more than $`\tau `$. The special case $`Y_j(t,u)=1+(1u/t)T_j`$ is studied in . Here we need the solution
(23)
$$Y_j(t,u)=\{\begin{array}{cc}T_j\hfill & \text{if }t>u,\hfill \\ 1+T_j\hfill & \text{if }tu,\hfill \end{array}$$
Indeed, checking the six different possibilities of order among $`t`$, $`u`$, $`v`$, one has:
###### Lemma 3.1.
The $`Y_j(t,u)`$ defined by Equation (23) satisfy Equation (21).
The *Yang-Baxter graph* associated with a word $`๐ฑ`$ is the graph whose vertices are $`Y_\sigma (๐ฑ)`$ and whose edges go from $`Y_\tau `$ to $`Y_{\sigma _j\tau }`$ with the label $`Y_j(x_{\tau (j)},x_{\tau (j+1)})`$ if $`\sigma _j\tau `$ has one inversion more than $`\tau `$. For example, one can check on Figure 6 that $`Y_{3241}(2431)=(1+T_2)T_1T_2T_3=T_1T_2(1+T_1)T_3=T_1T_2T_3(1+T_1)`$.
###### Note 3.2.
The conditions of Equation (23) show that $`Y_\sigma (๐ฑ)`$ does not depend on $`๐ฑ`$ but only on its standardized word $`\mathrm{Std}(๐ฑ)`$:
(24)
$$Y_\sigma (๐ฑ)=Y_\sigma (\mathrm{Std}(๐ฑ)).$$
We can therefore assume that $`๐ฑ`$ is a color word with no repeated colors, hence a permutation.
### 3.3. Yang-Baxter graphs of semi-combinatorial modules
We shall now generalize the definition of a combinatorial module as follows: a *semi-combinatorial $`H_n(0)`$-module* is a $`H_n(0)`$-module with a basis $`(b)`$ satisfying: either $`T_ib`$, $`T_ib`$ or $`(1+T_i)b`$ is a basis element.
In particular, from any interval $`[Y_\sigma (๐ฑ),Y_\tau (๐ฑ)]`$ in a Yang-Baxter graph, one can build the graph of a semi-combinatorial module of $`H_n(0)`$, in the basis $`(Y_\rho )`$, by adding loops on each vertex $`Y_\rho `$ as follows: if $`Y_{s_i\rho }`$ does not belong to the interval, add a loop on $`Y_\rho `$ labelled $`T_i`$ or $`1+T_i`$ depending on whether $`๐ฑ_{\rho ^1(i)}<๐ฑ_{\rho ^1(i+1)}`$ or not. As in the case of combinatorial $`H_n(0)`$-modules, one can read on the graph $`G`$ of a semi-combinatorial module as many composition series of the module as linear extensions of $`G`$. One can see an example of a restricted Yang-Baxter graph and the corresponding semi-combinatorial module on Figure 7.
Note that the interval $`[Y_\sigma (๐ฑ),Y_{\omega _n}(๐ฑ)]`$ corresponds to the left $`H_n(0)`$-module generated by $`Y_\sigma (๐ฑ)`$. Consequently, any interval $`[Y_\sigma (๐ฑ),Y_\tau (๐ฑ)]`$ corresponds to a quotient of the previous $`H_n(0)`$-module, hence to an $`H_n(0)`$-module itself.
###### Lemma 3.3.
Let $`๐ฑ`$ be a word, $`\sigma `$ a permutation and consider the interval of the left weak order $`[Y_\sigma (๐ฑ),Y_\sigma (๐ฑ)]`$ as a one-dimensional quotient of the $`H_n(0)`$-module generated by $`Y_\sigma (๐ฑ)`$. It is isomorphic to the simple $`H_n(0)`$-module labelled by the shape of the anticycloribbon associated with $`(\sigma ,๐ฑ)^1`$.
In particular, $`Y_{\omega _n}(๐ฑ)`$ is the simple $`H_n(0)`$-module labelled by the shape of the anticycloribbon associated with $`(\omega _n,\overline{๐ฑ})`$.
###### Proof.
The action of $`T_i`$ on the quotient $`[Y_\sigma (๐ฑ),Y_\sigma (๐ฑ)]`$ is encoded by the Yang-Baxter graph associated with $`๐ฑ`$, and the lemma follows from the definitions.
For example, the one-dimensional quotient of the left $`H_n(0)`$-module generated by $`Y_{1342}(3213)`$ has as corresponding anticycloribbon
(25) $`1^3`$ $`4^3`$ $`3^1`$
so that it is the simple module $`S_{13}`$, as one can check on Figure 7 (the generator of the quotient module is annihilated by $`1+T_1`$, $`T_2`$ and $`T_3`$).
### 3.4. Representations in Yang-Baxter bases
#### 3.4.1. Simple modules and indecomposable projective modules
We can now rewrite the representation theory of $`H_n(0)`$ in terms of Yang-Baxter bases. First, we identify the simple and indecomposable projective modules of $`H_n(0)`$ as the semi-combinatorial modules associated with some intervals of Yang-Baxter graphs. Then, using Theorem 3.8, we will be able to describe in a very simple way the induction of those modules in terms of intervals in the Yang-Baxter graphs.
We first consider the case of simple $`H_n(0)`$-modules. The next lemma is a rewriting of a special case of Lemma 3.3. Recall that $`\eta _I`$ is the generator of the simple module $`S_I`$ (see Equation (17)).
###### Lemma 3.4.
Let $`I`$ be a composition. Then for any permutation $`\tau `$ such that $`\overline{\mathrm{C}(\tau )}=I`$, one has
(26)
$$\eta _I=Y_{\omega _n}(\tau ).$$
Let $`\phi `$ be the anti-automorphism of $`H_n(0)`$ such that $`\phi (T_i)=T_i`$ for all $`i`$. One easily proves by induction on the number of inversions of $`\sigma `$ that
###### Lemma 3.5.
For all color words $`\tau `$ and all permutations $`\sigma `$, one has
(27)
$$\phi (Y_\sigma (\tau ))=Y_{\sigma ^1}(\overline{\tau }^\mathrm{\#}(\sigma ^1)^\mathrm{\#}).$$
These lemmas give in particular the description of $`\eta _I`$ as a right $`H_n(0)`$-module:
###### Proposition 3.6.
For all color words $`\tau `$, one has
(28)
$$\phi (Y_{\omega _n}(\tau ))=Y_{\omega _n}(\overline{\tau }^\mathrm{\#})),$$
so that, for all compositions $`I`$,
(29)
$$\phi (\eta _I)=\eta _{\overline{I}}.$$
Let us now consider the case of indecomposable projective $`H_n(0)`$-modules. For those modules, from the definition of Nortonโs generators (see Equation (19)), one obtains:
###### Proposition 3.7.
Let $`I`$ be a composition. Then
(30)
$$\nu _I=Y_{\alpha (I)\alpha (I^{})}(\omega (I)).$$
#### 3.4.2. Induction product of modules
###### Theorem 3.8.
Let $`M^{}`$ and $`M^{\prime \prime }`$ be the semi-combinatorial modules of $`H_k(0)`$ and $`H_{nk}(0)`$ associated with the intervals $`[Y_\alpha ^{}(\tau ^{}),Y_\beta ^{}(\tau ^{})]`$ and $`[Y_{\alpha ^{\prime \prime }}(\tau ^{\prime \prime }),Y_{\beta ^{\prime \prime }}(\tau ^{\prime \prime })]`$. Then, $`M^{}\widehat{}M^{\prime \prime }`$ is the semi-combinatorial module associated with the interval $`[Y_\alpha (\tau ),Y_\beta (\tau )]`$, where $`\tau `$ is any element of $`\tau ^{}\tau ^{\prime \prime }`$, $`\alpha =\alpha ^{}\alpha ^{\prime \prime }[k]`$ and $`\beta =\beta ^{\prime \prime }[k]\beta ^{}`$.
###### Proof.
First, the semi-combinatorial module $`M`$ corresponding to the restriction of the Yang-Baxter graph determined by $`\tau `$ to the interval $`[\alpha ,\beta ]`$ and $`M^{}\widehat{}M^{\prime \prime }`$ have the same dimension and both are $`H_n(0)`$-modules. Moreover, the quotient of $`M`$ by the submodule $`H_n(0)T_kM`$ is isomorphic to $`M^{}M^{\prime \prime }`$ as an $`H_k(0)H_{nk}(0)`$-module. So $`M`$ is isomorphic to $`M^{}\widehat{}M^{\prime \prime }`$.
Since any simple or indecomposable projective module is an interval of a Yang-Baxter graph (seen as a semi-combinatorial module), the induction of such modules is also such a Yang-Baxter graph. The choice we have in color $`\tau `$ in Theorem 3.8 will be useful when computing the induced module of two indecomposable projective modules in the $`0`$-Ariki-Koike-Shoji algebras.
###### Corollary 3.9.
Let $`I_1`$ and $`I_2`$ be compositions of $`n_1`$ and $`n_2`$, and $`๐ฑ=๐ฑ_1๐ฑ_2`$ a color word with $`๐ฑ_1C^{n_1}`$ and $`๐ฑ_2C^{n_2}`$. Let $`\sigma _1=\omega (I_1)`$ and $`\sigma _2=\omega (I_2)`$.
Then, the semi-combinatorial module corresponding to the restriction of the Yang-Baxter graph determined by $`๐ฑ`$ to the convolution $`\sigma _1\sigma _2`$ is isomorphic to the induced module $`S_{I_1^{}}\widehat{}S_{I_2^{}}`$ where $`I_1^{}`$ (resp. $`I_2^{}`$) is the shape of the anticycloribbon associated with the colored permutation $`(\sigma _1,๐ฑ_1)^1`$ (resp. $`(\sigma _2,๐ฑ_2)^1)`$.
###### Corollary 3.10.
Let $`I^{}`$ and $`I^{\prime \prime }`$ be two compositions of integers $`n^{}`$ and $`n^{\prime \prime }`$. Let $`๐ฑ^{}`$ and $`๐ฑ^{\prime \prime }`$ be two permutations such that $`I^{}=\overline{C(๐ฑ^{})}`$ and $`I^{\prime \prime }=\overline{C(๐ฑ^{\prime \prime })}`$. Then the interval $`[Y_\sigma ^๐ฑ,Y_{\omega _{n^{}+n^{\prime \prime }}}^๐ฑ]`$ of the Yang-Baxter graph of color $`๐ฑ`$ where $`\sigma =\omega _n^{}\omega _{n^{\prime \prime }}`$ and $`๐ฑ=๐ฑ^{}๐ฑ^{\prime \prime }`$ is isomorphic the induced module $`S_I^{}\widehat{}S_{I^{\prime \prime }}`$.
For example, the semi-combinatorial module represented on Figure 7 is isomorphic to $`S_2\widehat{}S_2`$, as one can check on Figure 4. The edges are not labelled by the same elements but the composition factors are the same. Since the generators of both modules are the same, one can easily provide an explicit description of the basis of the semi-combinatorial module in the left regular representation: it is a product of factors $`T_i`$ and $`1+T_i`$ multiplied on their right by $`\eta _{I_1^{}}\eta _{I_2^{}}`$ seen as an element of $`H_{|I_1|+|I_2|}(0)`$ through the morphism $`H_{|I_1|}(0)H_{|I_2|}(0)H_{|I_1|+|I_2|}(0)`$.
Another example is given Figure 8. The graphs are the semi-combinatorial modules corresponding to the restriction of the Yang-Baxter graph associated with the color words $`3124`$ and $`4231`$ to the same interval $`[2314,4321]`$. Both are the graph of the induced module $`P_{(2,1)}\widehat{}P_{(1)}`$ where $`P_I`$ is the indecomposable projective module associated with $`I`$.
## 4. The $`0`$-Ariki-Koike-Shoji algebras
In , Shoji obtained a new presentation of the Ariki-Koike algebras defined in . We shall first give a presentation very close to his, put $`q=0`$ in the relations and then prove some simple results about another basis of the resulting algebra. To get our presentation from Shojiโs , one has to replace $`qa_i`$ by $`T_{i1}`$ and $`q^2`$ by $`q`$.
Let $`u_1,\mathrm{},u_r`$ be $`r`$ distinct complex numbers. We shall denote by $`P_k(X)`$ the Lagrange polynomial
(31)
$$P_k(X):=\underset{1lr,lk}{}\frac{Xu_l}{u_ku_l}.$$
The Ariki-Koike algebra $`_{n,r}(q)`$ is the associative $``$-algebra generated by elements $`T_1,\mathrm{},T_{n1}`$ and $`\xi _1,\mathrm{},\xi _n`$ subject to the following relations:
(32) $`(T_iq)(T_i+1)`$ $`=0`$ $`(1in1),`$
(33) $`T_iT_{i+1}T_i`$ $`=T_{i+1}T_iT_{i+1}`$ $`(1in2),`$
(34) $`T_iT_j`$ $`=T_jT_i`$ $`(|ij|2),`$
(35) $`(\xi _ju_1)\mathrm{}(\xi _ju_r)`$ $`=0`$ $`(1jn),`$
(36) $`\xi _i\xi _j`$ $`=\xi _j\xi _i`$ $`(1i,jn),`$
(37) $`T_i\xi _i=\xi _{i+1}T_i(q1)`$ $`{\displaystyle \underset{c_1<c_2}{}}(u_{c_2}u_{c_1})P_{c_1}(\xi _i)P_{c_2}(\xi _{i+1})`$ $`(1in1),`$
(38) $`T_i(\xi _{i+1}+\xi _i)`$ $`=(\xi _{i+1}+\xi _i)T_i`$ $`(1in1)`$
(39) $`T_i\xi _j`$ $`=\xi _jT_i`$ $`(ji+1,i).`$
As noticed in , it is obvious from this presentation that a generating set is given by the $`\xi _1^{c_1}\mathrm{}\xi _n^{c_n}T_\sigma `$ with $`\sigma ๐_n`$ and $`c_i`$ such that $`0c_ir1`$. Shoji proves that this is indeed a basis of $`_{n,r}(q)`$. A simple adaptation of his proof would enable us to conclude that this property still holds at $`q=0`$. We will prove it by introducing a new basis on which the product of generators has a simple expression (see Lemma 4.1.) This algebra $`_{n,r}(0)`$, which we call the $`0`$-Ariki-Koike-Shoji algebra, will be our main concern in the sequel.
If $`๐=(c_1,\mathrm{},c_n)`$ is a word on $`C=\{1,\mathrm{},r\}`$, we define
(40)
$$L_๐:=P_{c_1}(\xi _1)\mathrm{}P_{c_n}(\xi _n).$$
Since the Lagrange polynomials (31) associated with $`r`$ distincts complex numbers form a basis of $`_{r1}[X]`$ (polynomials of degree at most $`r1`$), the set of elements of $`_{n,r}(q)`$
(41)
$$B:=\{B_{๐,\sigma }(q):=L_๐T_\sigma \},$$
where $`\sigma `$ runs over $`๐_n`$ and $`๐=(c_1,\mathrm{},c_n)`$ runs over the color words of size $`n`$ span the same vector space as the $`\xi ^๐T_\sigma `$.
Let us now describe the action of a generator on the left of $`B_{๐,\sigma }`$: the generator $`\xi _i`$ acts diagonally by multiplication by $`u_{c_i}`$, so that it only remains to explicit the product $`T_iL_๐T_\sigma `$. The relevant expression comes from the following apparently unnoticed relation in $`_{n,r}(q)`$:
###### Lemma 4.1.
The following relation holds in $`_{n,r}(q)`$:
(42)
$$T_iL_๐=L_{๐s_i}T_i(q1)\{\begin{array}{ccc}L_๐\hfill & \text{if}& c_i<c_{i+1},\hfill \\ 0\hfill & \text{if}& c_i=c_{i+1},\hfill \\ L_{๐s_i}\hfill & \text{if}& c_i>c_{i+1}.\hfill \end{array}$$
where $`s_i`$ acts on the right of $`๐`$ by exchanging $`c_i`$ and $`c_{i+1}`$.
###### Proof.
It is enough to compute $`T_iP_a(\xi _i)P_b(\xi _{i+1})`$ since all the other terms commute with $`T_i`$. Using Relations (36), (37), and (38), one proves
(43)
$$T_i\xi _i\xi _{i+1}=\xi _i\xi _{i+1}T_i,$$
so that $`T_i`$ commutes with any symmetric function of $`\xi _i`$ and $`\xi _{i+1}`$. The computation then reduces to evaluate the expression
(44)
$$\underset{1lr,la,lb}{}\left(\frac{\xi _iu_l}{u_au_l}\frac{\xi _{i+1}u_l}{u_au_l}\right)T_i\frac{\xi _iu_b}{u_au_b}\frac{\xi _{i+1}u_a}{u_bu_a}.$$
Now, writing $`(\xi _iu_b)(\xi _{i+1}u_a)`$ as $`\xi _i\xi _{i+1}u_b(\xi _i+\xi _{i+1})+u_au_b+(u_bu_a)\xi _i`$, the first three terms commute with $`T_i`$ and the last one, thanks to Equation (37), simplifies into a sum of two terms since all but one Lagrange polynomial vanish when multiplied with the left factor of Equation (44), the terms depending on the order relation between $`a`$ and $`b`$.
The last lemma implies that the set $`B`$ defined in Equation (41) is a generating set of $`_{n,r}(q)`$. We now show that it is a basis of $`_{n,r}(q)`$ for all values of $`q`$ and in particular for $`q=0`$. One can argue as in . Let $`V=_{i=1}^rV_i`$, with $`V_i=^{m_i}`$, $`m_in`$, and let $`W=\{v_1,\mathrm{},v_m\}`$ with $`m=m_1+\mathrm{}+m_r`$, be a basis of $`V`$ such that $`v_1,\mathrm{},v_{m_1}V_1`$, $`v_{m_1+1},\mathrm{},v_{m_1+m_2}V_2`$, and so on. Let $`w=v_{k_1}\mathrm{}v_{k_n}V^n`$. There is a classical right action of $`H_n(q)`$ on $`V^n`$ given by
(45)
$$\{\begin{array}{cccc}wT_i& =& ws_i& \text{if}k_i<k_{i+1},\hfill \\ wT_i& =& qw& \text{if}k_i=k_{i+1},\hfill \\ wT_i& =& qws_i+(q1)w& \text{if}k_i>k_{i+1}.\hfill \end{array}$$
Let $`\xi `$ be the linear map whose restriction to $`V_i`$ is $`u_iid_{v_i}`$, and $`\xi _j=id^{(j1)}\xi id^{(nj)}`$ acting on the right of $`V^n`$. Then, as observed by Shoji, the $`T_i`$ and $`\xi _j`$ generate a right action of $`_{n,r}(q)`$ on $`V^n`$. It is known that the representation (45) of $`H_n(q)`$ is faithful for all $`q`$, provided that $`m_in`$. This easily implies the next proposition:
###### Proposition 4.2.
The set $`B`$ defined in Equation (41) is a basis of $`_{n,r}(q)`$, for all values of $`q`$.
With exactly the same arguments, one would also prove:
###### Proposition 4.3.
The set $`B^{}=\{B_{๐,\sigma }^{}:=T_\sigma L_๐\}`$ is a basis of $`_{n,r}(q)`$, for all values of $`q`$.
We now provide a new presentation of $`_{n,r}(0)`$. Since $`_{n,r}(0)`$ satisfies the relations of the next proposition, a simple comparison of dimensions proves that:
###### Proposition 4.4.
The algebra with generators $`L_๐`$, $`๐`$ being any color word of size $`r`$, and $`T_i`$ with $`1in1`$, and relations
(46) $`T_i(1+T_i)=0(1in1),`$
(47) $`T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1}(1in2),`$
(48) $`T_iT_j=T_jT_i(|ij|2),`$
(49) $`L_๐L_๐=\delta _{๐,๐}L_๐,`$
(53) $`\{\begin{array}{ccc}\hfill \text{if}& c_i<c_{i+1},& (1+T_i)L_๐=L_{๐s_i}T_i,\hfill \\ \hfill \text{if}& c_i=c_{i+1},& T_iL_๐=L_{๐s_i}T_i,\hfill \\ \hfill \text{if}& c_i>c_{i+1},& T_iL_๐=L_{๐s_i}(1+T_i).\hfill \end{array}`$
is isomorphic to $`_{n,r}(0)`$.
This description of the Ariki-Koike-Shoji algebra enables us to analyze the left regular representation of $`_{n,r}(0)`$ in terms of our basis elements. As a first application, we shall obtain a classification of the simple $`_{n,r}(0)`$-modules.
### 4.1. The associative bilinear form
Recall that a bilinear form $`(,)`$ on an algebra $`A`$ is said to be associative if $`(ab,c)=(a,bc)`$ for all $`a,b,cA`$, and that $`A`$ is called a Frobenius algebra whenever it has a nondegenerate associative bilinear form. Such a form induces an isomorphism of left $`A`$-modules between $`A`$ and the dual $`A^{}`$ of the right regular representation. Frobenius algebras are in particular self-injective, so that finitely generated projective and injective modules coincide (see ).
It is known that $`H_n(0)`$ is a Frobenius algebra for the following bilinear form: for a basis $`(Y_\sigma )`$, we denote by $`(Y_\sigma ^{})`$ the dual basis. We set $`\chi =T_{\omega _n}^{}`$. Then $`(f,g)=\chi (fg)`$ is a non-degenerate associative bilinear form . Moreover if we denote $`\eta _i=T_i`$ and $`\xi _i=1+T_i`$ (both satisfy the braid relations) then the following properties hold:
(54)
$$\zeta _\sigma =(1)^{\mathrm{}(\omega \sigma ^1)}\xi _{\omega \sigma ^1},\xi _\alpha =\underset{\beta \alpha }{}T_\beta \text{and}\eta _\alpha =\underset{\beta \alpha }{}(1)^{\mathrm{}(\beta )}\xi _\beta $$
shows that $`(\zeta _\sigma )`$ and $`(\eta _\tau )`$ are two adjoint bases of $`H_n(0)`$, that is
(55)
$$(\zeta _\sigma ,\eta _\tau )=\delta _{\sigma ,\tau }.$$
This bilinear form can be extended to $`_{n,r}(0)`$ by setting
(56)
$$\mathrm{X}:=\underset{๐}{}B_{๐,\omega }^{}$$
where $`B`$ is the basis defined by Equation (41).
###### Proposition 4.5.
The associative bilinear form defined by
(57)
$$(f,g):=\mathrm{X}(fg)$$
is non-degenerate on $`_{n,r}(0)`$. Therefore, $`_{n,r}(0)`$ is a Frobenius algebra.
###### Proof.
The form is obviously bilinear and associative. It remains to prove that it is non-degenerate. We have
(58)
$$L_๐T_\sigma =T_\sigma L_{๐\sigma ^1}+\text{smaller terms},$$
so that
(59)
$$(L_๐T_\sigma ,T_\tau L_๐)=\mathrm{X}(L_๐T_\sigma T_\tau L_๐)=\delta _{๐\overline{๐}}\mathrm{X}(T_\sigma T_\tau ).$$
This formula being linear on $`T_\sigma `$ and $`T_\tau `$, we have as well,
(60)
$$\mathrm{X}(L_๐\xi _\alpha \eta _\beta L_{\overline{๐}})=\delta _{\mathrm{๐๐}}\delta _{\alpha \beta },$$
so that $`B_{๐,\alpha }^{}:=L_๐\xi _\alpha `$ and $`B_{๐,\beta }^{\prime \prime }:=\eta _\beta L_{\overline{d}}`$ are two adjoint bases.
In particular, we have
###### Corollary 4.6.
$`_{n,r}(0)`$ is self-injective.
## 5. Simple modules of $`_{n,r}(0)`$
Let $`[I,๐]`$ be a cycloribbon. Let $`\eta _{[I,๐]}`$ be the element of $`_{n,r}(0)`$ defined as
(61)
$$\eta _{[I,๐]}:=L_๐\eta _I,$$
where $`\eta _I`$ is the generator of the simple $`H_n(0)`$-module associated with $`I`$ as in Section 3.1. This element generates a simple $`_{n,r}(0)`$-module:
###### Theorem 5.1.
Let $`[I,๐]`$ be a cycloribbon. Then
(62)
$$S_{[I,๐]}:=_{n,r}(0)\eta _{[I,๐]}$$
is a simple module of $`_{n,r}(0)`$, realized as a minimal left ideal in its left regular representation. The eigenvalue of $`L_๐`$ is $`1`$ if $`๐=๐`$ and $`0`$ otherwise, and that of $`T_i`$ is $`1`$ or $`0`$ according to whether $`i`$ is a descent of the shape of $`\varphi (R)`$ or not. All these simple modules are pairwise non isomorphic and of dimension $`1`$. Moreover, all one-dimensional $`_{n,r}(0)`$-modules are isomorphic to some $`S_{[I,๐]}`$.
###### Proof.
We prove simultaneously that the $`S_{[I,๐]}`$ are one-dimensional $`_{n,r}(0)`$-modules and that all one-dimensional modules are isomorphic to some $`S_{[I,๐]}`$. Let us first consider a simple module of dimension $`1`$. Since $`_{n,r}(0)`$ is self-injective (Corollary 4.6), it can be realized as a minimal left ideal in the left regular representation, that is, as $`_{n,r}(0)S`$ with $`S=_{๐,\sigma }C_{๐,\sigma }B_{๐,\sigma }`$ with $`C_{๐,\sigma }`$. Since all $`B_{๐,\sigma }`$ are eigenvectors of $`L_๐`$ with eigenvalue $`1`$ or $`0`$ depending on whether $`๐=๐`$ or not, $`S`$ is an eigenvector for all $`L_๐`$ iff the only non-zero coefficients correspond to the same $`๐`$. So $`S=L_๐S^{}`$ with $`S^{}=_\sigma C_\sigma T_\sigma `$. Since $`S`$ is also an eigenvector for all $`T_i`$, we are left with the following possibilities:
* If $`c_i<c_{i+1}`$, then thanks to Formula (53), first case, $`S`$ is an eigenvector of $`T_i`$ iff $`T_i`$ acts by $`0`$ on $`S^{}`$.
* If $`c_i=c_{i+1}`$, then thanks to Formula (53), second case, $`S`$ is an eigenvector of $`T_i`$ iff $`T_i`$ acts by $`0`$ or $`1`$ on $`S^{}`$.
* If $`c_i>c_{i+1}`$, then thanks to Formula (53), third case, $`S`$ is an eigenvector of $`T_i`$ iff $`T_i`$ acts by $`1`$ on $`S^{}`$.
So, in any case, $`S^{}`$ is an eigenvector for all $`T_i`$, so it is equal to some $`\eta _I`$. Moreover, $`I`$ has to be compatible with $`๐`$ in the following sense: if $`c_i<c_{i+1}`$ then $`i`$ cannot be a descent of $`I`$ and if $`c_i>c_{i+1}`$ then $`i`$ has to be a descent of $`I`$. This is precisely the definition of a cycloribbon. All those simple modules are non-isomorphic since the eigenvalues characterize them completely.
We will prove in Section 5.1 that all simple modules are of dimension one, thus isomorphic to some $`S_{[I,๐]}`$ by following the same pattern as in , proving that all composition factors of the modules induced from simple modules of the $`0`$-Hecke algebra are one-dimensional.
Notice that if $`[I,๐]`$ is a cycloribbon, then $`\eta _{[I,๐]}`$ is an eigenvector for all $`T_i`$, so that it can be written (since $`B^{}`$ is a basis of $`_{n,r}(0)`$, see Proposition 4.3) as $`\eta _I^{}_๐^{}\alpha _๐^{}L_๐^{}`$, where $`I^{}`$ is the shape of the anticycloribbon $`\varphi ([I,๐])`$ and $`\alpha _๐^{}`$ are unknown coefficients. The result is actually simpler:
###### Proposition 5.2.
Let $`[I,๐]`$ be a cycloribbon. Then,
(63)
$$\eta _{[I,๐]}=\eta _I^{}L_{\overline{๐}},$$
where $`I^{}`$ is the shape of the anticycloribbon $`\varphi ([I,๐])`$.
###### Proof.
Let $`n:=|I|`$ and let $`\rho ^{}`$ be a color word involving $`n`$ distinct colors such that
(64)
$$\{\begin{array}{ccc}\rho _i^{}<\rho _j^{}\hfill & \text{ if }& \hfill \overline{๐}_i<\overline{๐}_j,\\ \rho _i^{}>\rho _j^{}\hfill & \text{ if }& \hfill \overline{๐}_i>\overline{๐}_j.\end{array}$$
Let us define another color word $`\rho `$ as the permutation such that
(65)
$$\rho _i>\rho _j\text{ iff }\{\begin{array}{ccc}\rho _i^{}>\rho _j^{}\hfill & \text{and}& ๐_i=๐_j,\text{ or}\hfill \\ \rho _i^{}<\rho _j^{}\hfill & \text{and}& ๐_i๐_j.\hfill \end{array}$$
Thanks to Relation (53), one proves by induction that
(66)
$$Y_\sigma (\rho ^{})L_{\overline{c}}=L_{\overline{c}\sigma ^1}Y_\sigma (\rho ).$$
Thanks to Lemmas 3.5 and 3.4, one then gets
(67)
$$Y_{\omega _n}(\rho ^{})=\eta _{\overline{\mathrm{C}(\rho ^{})}}.$$
Now, it is easy to build a permutation $`\rho ^{}`$ satisfying these properties, with descent composition equal to $`\overline{I^{}}`$ iff $`I^{}`$ is a composition such that $`[I^{},๐]`$ is an anticycloribbon (the descent composition of $`\rho ^{}`$ is, by definition of $`\rho ^{}`$, already included in this set of compositions). It then comes that
(68)
$$\eta _I^{}L_{\overline{c}}=L_{\overline{c}\omega _n}\eta _{\overline{\mathrm{C}(\rho )}}=L_๐\eta _{\overline{\mathrm{C}(\rho )}}.$$
By definition of $`\rho `$, $`[\overline{\mathrm{C}(\rho )},๐]`$ is the cycloribbon $`\varphi ([I^{},๐])`$.
This result was not unexpected: it is clear that the generator of $`S_{[I,๐]}`$ has an expression as $`\eta _I^{}L_๐^{}`$, and from Lemma 3.4 and Equation (53), one can see that $`๐^{}`$ has to be $`\overline{๐}=๐\omega _n`$.
### 5.1. Induction of the simple $`0`$-Hecke modules
To describe the induction process, we need a partial order on colored ribbons. Let $`I`$ be a composition and $`๐=(c_1,\mathrm{},c_n)`$. The covering relation of the order $`_I`$ amounts to sort in increasing order any two adjacent elements in the rows of $`I`$ or to sort in decreasing order any two adjacent elements in the columns of $`I`$. For example, the elements smaller than or equal to
(69)
$$T:=\text{ โ 2 โ 1 โ 3 โ 4 โ 3 }$$
are
(70) $`2`$ $`1`$ $`3`$ $`4`$ $`3`$ $`1`$ $`2`$ $`3`$ $`4`$ $`3`$ $`2`$ $`1`$ $`4`$ $`3`$ $`3`$ $`1`$ $`2`$ $`4`$ $`3`$ $`3`$
If $`I`$ is a composition of $`n`$, let $`S_I:=H_n(0)\eta _I`$ be the corresponding simple module of $`H_n(0)`$ realized as a minimal left ideal in its left regular representation (see ) and
(71)
$$M_I:=S_I_{H_n(0)}^{_{n,r}(0)}.$$
Clearly, $`M_I`$ has dimension $`r^n`$ and admits $`L_๐\eta _I`$ as linear basis, when $`๐`$ runs over color words. For $`๐C^n`$, let $`M_{๐,I}`$ be the $`_{n,r}(0)`$-submodule of $`M_I`$ generated by $`L_๐\eta _I`$.
###### Lemma 5.3.
(72)
$$M_{๐,I}M_{๐^{},I}๐_I๐^{}.$$
###### Proof.
Let $`i\{1,\mathrm{},n1\}`$.
* If $`c_i<c_{i+1}`$, and $`T_i`$ acts by $`0`$ on $`\eta _I`$, we get $`(1+T_i)L_๐\eta _I=L_{๐\sigma _i}T_i\eta _I=0`$.
* If $`c_i<c_{i+1}`$, and $`T_i`$ acts by $`1`$ on $`\eta _I`$, we get $`(1+T_i)L_๐\eta _I=L_{๐\sigma _i}T_i\eta _I=L_{๐\sigma _i}\eta _I`$, and so $`(1+T_i)`$ sorts in decreasing order $`c_i`$ and $`c_{i+1}`$.
* If $`c_i=c_{i+1}`$, then $`T_iL_๐\eta _I=L_๐T_i\eta _I`$ so the result is either $`0`$ or $`1`$ times $`L_๐\eta _I`$.
* If $`c_i>c_{i+1}`$, and $`T_i`$ acts by $`1`$ on $`\eta _I`$, we get $`T_iL_๐\eta _I=L_{๐\sigma _i}(1+T_i)\eta _I=0`$.
* If $`c_i>c_{i+1}`$, and $`T_i`$ acts by $`0`$ on $`\eta _I`$, we get $`T_iL_๐\eta _I=L_{๐\sigma _i}(1+T_i)\eta _I=L_{๐\sigma _i}\eta _I`$, and so $`T_i`$ sorts in increasing order $`c_i`$ and $`c_{i+1}`$.
We can now complete Theorem 5.1 by proving that we have a complete set of simple $`_{n,r}(0)`$-modules:
###### Theorem 5.4.
The $`S_{[I,๐]}`$ form a complete family of simple modules of $`_{n,r}(0)`$.
###### Proof.
It is sufficient to prove that any simple $`_{n,r}(0)`$-module $`M`$ appears as a composition factor of some $`M_I`$ since by the previous lemma, $`M_I`$ admits a composition series involving only one-dimensional modules.
Now, if $`M`$ is any simple $`_{n,r}(0)`$-module, let $`S_I`$ be a simple $`H_n(0)`$-module occuring in the socle of $`MH_n(0)`$. This means that we have a monomorphism $`i:S_iM`$ of $`H_n(0)`$-modules. Inducing to $`_{n,r}(0)`$, we get a non-zero $`_{n,r}(0)`$-morphism $`f:M_IM`$, which has to be surjective since $`M`$ is simple. So $`M`$ is a composition factor of some $`M_I`$ and all simple $`_{n,r}(0)`$-modules are one-dimensional.
### 5.2. First Grothendieck ring
The induction product of two simple modules of $`_{n,r}(0)`$ can be worked out in a way very similar to the case of $`H_n(0)`$-modules such as recalled in Section 3.1. As in Proposition 5.2, we represent these modules by their anticycloribbons. In the context of $`_{n,r}(0)`$, one has to extend to colored permutations what was done with permutations in $`H_n(0)`$.
Let us consider two anticycloribbons $`[I,๐]`$ and $`[I^{},๐^{}]`$. The structure of the induction product $`S_{[I,๐]}\widehat{}S_{[I^{},๐^{}]}`$ of the corresponding simple modules is completely encoded in the graph of the convolution of the associated colored permutations: as in $`H_n(0)`$, there exists a basis $`(b)`$ of $`S_{[I,๐]}\widehat{}S_{[I^{},๐^{}]}`$ naturally labelled by the elements occuring in the convolution, such that, either $`T_i`$, $`T_i`$, or $`1+T_i`$ sends a given basis element $`b`$ to another one, and $`L_๐b`$ is $`b`$ or $`0`$. The unique one-dimensional quotient of the module generated by an element $`b`$ is the one labelled by the anticycloribbon associated with its colored permutation. Let us describe more precisely this graph.
###### Algorithm 5.5.
*Input*: Two anticycloribbons $`[I_1,๐_1]`$ and $`[I_2,๐_2]`$.
*Output*: A graph $`G`$.
Let $`(\sigma _1,๐_1)`$ and $`(\sigma _2,๐_2)`$ be the inverse of the corresponding maximal colored permutations and let $`๐=๐_1๐_2`$.
Let $`G`$ be the graph with vertices $`V_w`$ labelled by the elements of the convolution $`\sigma _1\sigma _2`$. The edges of $`G`$ are labelled by the generators of our presentation of $`_{n,r}(0)`$ (see Proposition 4.4), that is, all Lagrange polynomials $`L_๐`$ and, for all $`1in1`$, the Hecke generator $`T_i`$, $`T_i`$ or $`1+T_i`$.
The edge labelled by $`L_๐`$ on the vertex $`V_\tau `$ is a loop on $`V_\tau `$ iff $`๐=๐\tau ^1`$. Otherwise, $`V_\tau `$ is annihilated and the edge is not represented. The edges labelled by the Hecke generators depend on the following cases:
* if the letters $`\tau _i`$ and $`\tau _{i+1}`$ come from the same factor of the convolution, there is a loop labelled by $`T_i`$ (resp. $`1+T_i`$) if $`i`$ is (resp. not) an anti-descent of $`(\tau ^1,๐\tau ^1)`$.
* if $`\tau _i`$ comes from $`\sigma _1`$ and $`\tau _{i+1}`$ comes from $`\sigma _2`$, there is a vertex from $`V_\tau `$ to $`V_{s_i\tau }`$ labelled by $`1+T_i`$ (resp. $`T_i`$) if $`i`$ is (resp. not) an anti-descent of $`(\tau ,๐)^1`$.
* if $`\tau _i`$ comes from $`\sigma _2`$ and $`\tau _{i+1}`$ comes from $`\sigma _1`$, there is a loop labelled by $`T_i`$ (resp. $`1+T_i`$) if $`i`$ is (resp. not) an anti-descent of $`(\tau ,๐)^1`$.
For example, the second graph of Figure 9 is the graph associated with the anticycloribbons $`[(1,1),23]`$ and $`[(1,1),13]`$.
###### Theorem 5.6.
Let $`[I_1,๐_1]`$ and $`[I_2,๐_2]`$ be two anticycloribbons. Let $`V`$ the vector space with basis $`V`$ indexed by the vertices of the graph of Algorithm 5.5. Then the relations $`\rho V_i=V_j`$ whenever $`V_i\stackrel{๐}{}V_j`$ is a presentation of the induced module $`S_{[I_1,๐_1]}\widehat{}S_{[I_2,๐_2]}`$. Its composition factors are given by the shapes of the anticycloribbons corresponding to the colored permutations $`(\tau ,๐)^1`$, where $`\tau `$ runs over the labellings of the vertices $`V`$.
For example, Figure 9 presents the induction product of the two simple modules $`S_{[(1,1),23]}`$ and $`S_{[(1,1),13]}`$ (written as anticycloribbons) of $`_{2,2}(0)`$ to $`_{4,2}(0)`$. In particular, we get the following composition factors (written as anticycloribbons):
(73) $`3`$ $`2`$ $`1`$ $`3`$ $`3`$ $`1`$ $`3`$ $`1`$ $`3`$ $`3`$ $`1`$ $`2`$ $`3`$ $`1`$ $`2`$ $`3`$ $`1`$ $`2`$
Note that this construction gives an effective algorithm to compute the composition factors of the induction product of two simple modules. One recovers the same shapes as in Figure 4, that is, the induction of the simple modules $`(2)`$ and $`(2)`$ of $`H_2(0)`$ in $`H_4(0)`$. Moreover, one recovers these shapes in the same order as in Figure 7. This property comes from Proposition 5.2, since any basis element $`V_{(\sigma ,๐)}`$ can be written either in the form $`L_๐T`$ or $`T^{}L_๐^{}`$, where $`T`$ and $`T^{}`$ are in $`H_n(0)`$.
###### Proof.
Let us consider two simple modules $`M=S_{[I_1,๐_1]}`$, $`N=S_{[I_2,๐_2]}`$ of $`_{m,r}(0)`$ and $`_{n,r}(0)`$ and let $`\sigma _1`$ and $`\sigma _2`$ be the maximal colored permutations associated with their anticycloribbons. Let $`๐=๐_1๐_2`$ be the concatenation of the color words.
The induced module $`M\widehat{}N`$ of $`_{m+n,r}(0)`$ is by definition
(74)
$$_{m+n,r}(0)\underset{_{m,r}(0)_{n,r}(0)}{}\left(M_{}N\right).$$
Since every $`L_๐`$ belongs to $`_{m,r}(0)_{n,r}(0)`$, the induced module is isomorphic to $`H_{m+n}(0)_{H_m(0)H_n(0)}(MN)`$ as a vector space and even as a $`H_{m+n}(0)`$-module. In particular, its dimension is $`\left(\genfrac{}{}{0pt}{}{m+n}{n}\right)`$.
Now, as a $`H_{m+n}(0)`$-module, thanks to Lemma 3.1, $`M\widehat{}N`$ is described by the restriction $`G_1`$ of the graph of $`๐_{m+n}`$ to $`\sigma _1\sigma _2`$, with the Yang-Baxter elements determined by the color word $`๐`$. Let $`G`$ be the graph associated with $`M`$ and $`N`$ by Algorithm 5.5. The edges of $`G`$ and $`G_1`$ are the same since both are edges in $`\sigma _1\sigma _2`$ with the same constraints. The loops of $`G`$ and $`G_1`$ labelled by the Hecke generators are the same thanks to Algorithm 5.5 and Corollary 3.9. So $`G`$ satisfies the relations of the $`0`$-Hecke algebra, relations (46), (47), and (48). It trivially satisfies Relations (49) and it satisfies Relations (53), since each vertex of the graph can be written as some Lagrange polynomial $`L_๐`$ multiplied by an element of $`H_n(0)`$. So it satisfies all the relations of our presentation of the Ariki-Koike-Shoji algebra at $`q=0`$ (Proposition 4.4). Since it has the correct dimension and the correct generator, this graph encodes the structure of the induced module $`M\widehat{}N`$.
Finally, the composition factors are the anticycloribbons associated with the colored permutations $`(\tau ,๐)^1`$: the only non-zero $`L`$ is the right one and the action of $`T_i`$ is the one encoded by the shape of the anticycloribbon, thanks to Note 3.3.
### 5.3. Restriction of simple modules
###### Proposition 5.7.
Let $`[I,๐]`$ be a cycloribbon, regarded as a filling of the ribbon diagram of $`I`$. Let $`I^{}`$ and $`I^{\prime \prime }`$ be the compositions whose diagrams are formed respectively of the first $`k`$ cells and the last $`nk`$ cells of $`I`$ and $`๐^{}=(c_1,\mathrm{},c_k)`$, $`๐^{\prime \prime }=(c_{k+1},\mathrm{},c_n)`$. Then $`[I^{},๐^{}]`$ and $`[I^{\prime \prime },๐^{\prime \prime }]`$ are again cycloribbons and the restriction of $`S_{[I,๐]}`$ to $`_{k,r}(0)_{nk,r}(0)`$ is $`S_{[I^{},๐^{}]}S_{[I^{\prime \prime },๐^{\prime \prime }]}`$.
For example, Figure 10 shows that the restriction of the first module $`S_{[31122,122211323]}`$ to $`_{5,r}(0)_{4,r}(0)`$ is $`S_{[311,12221]}S_{[22,1323]}`$.
###### Proof.
This follows from the obvious graphical description of the restriction of simple $`H_n(0)`$-modules, which consists in cutting the ribbon diagram, as was done before on cycloribbons, and from the commutation relations between $`L`$ and $`T`$.
### 5.4. The quasi-symmetric Mantaci-Reutenauer algebra
A *colored composition* is a pair $`(I,u)=((i_1,\mathrm{},i_m),(u_1,\mathrm{},u_m))`$ formed of a composition $`I`$ and a color word of the same length.
There is a bijection between colored compositions and anticycloribbons: starting with a colored composition, one rebuilds an anticycloribbon by separating adjacent blocks of the colored composition with different colors and gluing them back in the only possible way according to the criterion of being an anticycloribbon. Conversely, one separates two adjacent blocks of different colors and glue them one below the other. For example, Figure 11 shows a colored composition represented as a filling of a ribbon ($`u_i`$ is written in row $`i`$) and its corresponding anticycloribbon.
Let $`X^{(i)}=\{x_k^{(i)}|k1\}`$ for $`1ir`$ be $`r`$ infinite linearly ordered sets of commuting indeterminates. The *monochromatic monomial quasi-symmetric functions* labelled by a colored composition $`(I,u)`$ is
(75)
$$M_{(I,u)}(X)=\underset{j_1<\mathrm{}<j_m}{}(x_{j_1}^{(u_1)})^{i_1}\mathrm{}(x_{j_m}^{(u_m)})^{i_m}.$$
It is known ( Proposition 5.2), that the $`M_{I,u}`$ span a subalgebra $`\mathrm{๐๐๐
}^{(r)}`$ of $`[X^{(1)},\mathrm{},X^{(r)}]`$. This algebra is also a subalgebra of the algebra of level $`r`$ quasi-symmetric functions defined in (see also ).
We will prove that the Grothendieck ring of the tower of $`_{n,r}(0)`$ algebras is isomorphic to $`\mathrm{๐๐๐
}^{(r)}`$.
Define an order on colored compositions as follows: let
$$(I,u)=((i_1,\mathrm{},i_m),(u_1,\mathrm{},u_m))\text{ and }(J,v)=((j_1,\mathrm{},j_p),(v_1,\mathrm{},v_p))$$
two colored compositions. Then $`(I,u)`$ is *finer* than $`(J,v)`$ if there exists a sequence $`(l_0=0,l_1,\mathrm{},l_p=m)`$ such that for any integer $`k`$,
(76)
$$j_k=i_{l_{k1}+1}+\mathrm{}+i_{l_k}\text{ and }v_k=u_{l_{k1}+1}=\mathrm{}=u_{l_k}.$$
For example, the compositions fatter than $`((1,1,3,2),(2,1,1,2))`$ are
(77)
$$((1,1,3,2),(2,1,1,2))\text{ and }((1,4,2),(2,1,2)).$$
This allows us to define the *monochromatic quasi-ribbon functions* $`F_{(I,u)}`$ by
(78)
$$F_{(I,u)}=\underset{(I^{},u^{})(I,u)}{}M_{(I^{},u^{})}.$$
Notice that this last description of the order $``$ is reminiscent of the order $`^{}`$ on descent sets used in the context of quasi-symmetric functions and non-commutative symmetric functions: more precisely, one gets the usual order when considering compositions colored with only one color.
###### Lemma 5.8 ().
The $`F_{(I,u)}`$, where $`(I,u)`$ runs over colored compositions, span $`\mathrm{๐๐๐
}^{(r)}`$. This algebra is isomorphic to the dual of the Mantaci-Reutenauer algebra $`\mathrm{๐๐}^{(r)}`$ defined in .
Let
(79)
$$๐ข^{(r)}:=\underset{n0}{}G_0(_{n,r}(0))$$
be the Grothendieck ring of the tower of $`_{n,r}(0)`$ algebras. Define a characteristic map
(80)
$$\begin{array}{ccc}\hfill ch:& ๐ข^{(r)}& \mathrm{๐๐๐
}^{(r)}\hfill \\ & [S_{[I,๐]}]& F_{[I,๐]}.\hfill \end{array}$$
Comparing the descriptions of $`S_{[I^{},๐^{}]}\widehat{}S_{[I^{\prime \prime },๐^{\prime \prime }]}`$ and of $`S_{[I,๐]}`$ obtained in Sections 5.2 and 5.3 with the formulas for product and coproduct of the $`F`$ basis in $`\mathrm{๐๐๐
}^{(r)}`$, we obtain:
###### Theorem 5.9.
$`ch`$ is an isomorphism of Hopf algebras.
## 6. Projective modules
### 6.1. The Grothendieck ring
The previous results imply, by duality, a description of the Grothendieck ring of the category of projective $`_{n,r}(0)`$-modules. Indeed, we already know that $`_{n,r}(0)`$ is self injective. Moreover, it is easy to see that $`_{n,r}(0)`$ is a free module and hence projective over any parabolic subalgebra
(81)
$$_{n_1,r}(0)\mathrm{}_{n_k,r}(0).$$
These conditions are sufficient to ensure that, as a Hopf algebra,
(82)
$$๐ฆ=\underset{n0}{}K_0(_{n,r}(0))\mathrm{๐๐}^{(r)}$$
is the graded dual of $`๐ข`$, and thus is isomorphic to the Mantaci-Reutenauer algebra. Under this isomorphism, the classes of indecomposable projective modules are mapped to the subfamily of the dual basis $`F_{[I,๐]}`$ of Poirier quasi-symmetric functions labelled by colored descent sets, or cycloribbons.
Recall that the Hopf algebra $`\mathrm{๐๐}^{(r)}`$ can be defined (see ) as the free associative algebra over symbols $`(S_j^{(i)})_{j1;1ir}`$, graded by $`\mathrm{deg}S_j^{(i)}=j`$, and with coproduct
(83)
$$\mathrm{\Delta }S_n^{(k)}=\underset{i=0}{\overset{n}{}}S_i^{(k)}S_j^{(k)}.$$
For a colored composition $`(I,u)`$ as above, one defines
(84)
$$S^{(I,u)}:=S_{i_1}^{(u_1)}\mathrm{}S_{i_p}^{(u_p)}.$$
Clearly, the $`S^{(I,u)}`$ form a linear basis of $`\mathrm{๐๐}^{(r)}`$.
The *monochromatic colored ribbon basis* $`R_{(I,u)}`$ of $`\mathrm{๐๐}^{(r)}`$ can now be defined by the condition
(85)
$$S^{(I,u)}=:\underset{(J,v)(I,u)}{}R_{(J,v)}.$$
Let $`(I,u)=(i_1,\mathrm{},i_p;u_1,\mathrm{},u_p)`$ and $`(J,v)=(j_1,\mathrm{},j_q;v_1,\mathrm{},v_q)`$ be two colored compositions. We set
(86)
$$(I,u)(J,v):=(IJ,uv),$$
where $`ab`$ denotes concatenation.
Moreover, if $`u_p=v_1`$, we set
(87)
$$(I,u)(J,v):=((i_1,\mathrm{},i_{p1},(i_p+j_1),j_2,\mathrm{},j_q);(u_1,\mathrm{},u_p,v_2,\mathrm{},v_q)).$$
The colored ribbons satisfy the very simple multiplication rule:
(88)
$$R_{(I,u)}R_{(J,v)}=R_{(I,u)(J,v)}+\{\begin{array}{cc}R_{(I,v)(J,v)}\hfill & \text{if}u_p=v_1,\hfill \\ 0\hfill & \text{if}u_pv_1.\hfill \end{array}$$
Let $`[K,๐]`$ be an anticycloribbon. Let $`P_{[K,๐]}`$ be the indecomposable projective module whose unique simple quotient is the simple module labelled by $`\varphi ([K,๐])`$ and let $`(I,u)`$ be the corresponding colored composition. Summarizing this discussion, we have
###### Theorem 6.1.
The map
(89)
$$\begin{array}{ccc}\hfill \mathrm{๐๐ก}:& ๐ฆ& \mathrm{๐๐}^{(r)}\hfill \\ & [P_{[K,๐]}]& R_{(I,u)}\hfill \end{array}$$
is an isomorphism of Hopf algebras.
The main interest of the labelling by colored compositions is that it allows immediate reading of some important information. For example, it follows from Theorem 6.1 that the products of complete functions $`S^{(I,u)}`$ are the characteristics of the projective $`_{n,r}(0)`$-modules obtained as induction products of the one-dimensional projective $`_{m,r}(0)`$-modules on which all $`T_i`$ act by $`0`$ and all $`\xi _j`$ by the same eigenvalue $`u_i`$.
We see that, as in the case of $`H_n(0)`$, each indecomposable projective $`_{n,r}(0)`$-module occurs as a direct summand of such an induced module, and that the direct sum decomposition is given by the anti-refinement order. For example, writing a bar over the parts of the composition corresponding to color $`2`$ (and nothing over the parts corresponding to color $`1`$), the identity
(90)
$$S^{(2\overline{1}\overline{2}13)}=R_{(2\overline{1}\overline{2}13)}+R_{(2\overline{3}13)}+R_{(2\overline{1}\overline{2}4)}+R_{(2\overline{3}4)}$$
indicates which indecomposable projective direct summands compose the projective $`_{9,2}(0)`$-module defined as the outer tensor product
(91)
$$S_2\widehat{}S_{\overline{1}}\widehat{}S_{\overline{2}}\widehat{}S_1\widehat{}S_3.$$
Since by multiplying a Lagrange element $`L_๐`$ of $`_{n,r}(0)`$ by any element of $`_{n,r}(0)`$, one can only obtain an element of the form $`L_๐T`$ where $`๐`$ has same evaluation as $`๐`$ and $`T`$ is any Hecke element, we have a direct sum decomposition of $`_{n,r}(0)`$ into two-sided ideals:
(92)
$$_{n,r}(0)=\underset{e}{}H_{n,r}^e,$$
where $`e`$ is any evaluation of $`C^n`$ and $`H_{n,r}^e:=_{n,r}(0)_{๐e}L_๐`$. Actually all such sums are central idempotents of $`_{n,r}(0)`$. The ideals $`H_{n,r}^e`$ where $`e=(i,i,\mathrm{},i)`$ are said to be *monochromatic* of color $`i`$. The monochromatic ideals are isomorphic to $`H_n(0)`$ as algebras, and their indecomposable projective modules $`P_{[I,๐]}`$ are obtained by defining the action of all $`L_๐`$ by zero except $`L_{(i^n)}`$ acting by one on the $`H_n(0)`$-modules $`P_I`$. Such modules are also called monochromatic. Note that one can realize $`P_{[I,i^k]}`$ in the left regular representation as
(93)
$$P_{[I,i^k]}=_{n,r}(0)L_{i^n}\nu _I.$$
It follows from Theorem 6.1 and Formula (88) that
###### Corollary 6.2.
Every indecomposable projective module $`P_{[I,๐]}`$ of $`_{n,r}(0)`$ is obtained as an induction product of monochromatic modules
(94)
$$P_{[I,๐]}=P_{[I_1,๐_1]}\widehat{}\mathrm{}\widehat{}P_{[I_k,๐_k]},$$
where $`I=(I_1,\mathrm{},I_k)`$ and $`I_j`$ is a composition, $`๐=(๐_1,\mathrm{},๐_k)`$ and $`๐_j`$ is a monochromatic word.
Corollary 6.2 and Theorem 3.8 implies that $`P_{[I,๐]}`$ is a semi-combinatorial module. Indeed, it is possible to choose for the underlying $`H_n(0)`$-module $`P_{I_1}\widehat{}\mathrm{}\widehat{}P_{I_k}`$ a Yang-Baxter graph such that all vertices are common eigenvectors of the Lagrange polynomials by making the right choice on the color vector $`๐`$: if $`c_i<c_j`$ then $`d_i<d_j`$ and if $`c_i>c_j`$ then $`d_i>d_j`$.
Figure 12 shows the two examples of induction product $`P_{[(2,1),(1)]}\widehat{}P_{[(1),(2)]}`$ and $`P_{[(2,1),(2)]}\widehat{}P_{[(1),(1)]}`$. On the first induced module, the color vector has to be in the convolution of $`312`$ and $`1`$ and have its last element greater than all the others, and so is $`3124`$. On the second induced module, the color vector has to be in the convolution of $`312`$ and $`1`$ and have its last element smaller than all the others, and so is $`4231`$. One can compare this figure with Figure 8.
Their restrictions to $`H_n(0)`$ (and hence, their dimensions) can be computed by means of the following result:
###### Corollary 6.3.
The homomorphism of Hopf algebras
(95)
$$\begin{array}{ccc}\hfill \pi :& \mathrm{๐๐}^{(r)}& \mathrm{๐๐ฒ๐ฆ}\hfill \\ & S_j^{(i)}& S_j\hfill \end{array}$$
maps the class of a projective $`_{n,r}(0)`$-module to the class of its restriction to $`H_n(0)`$.
###### Proof.
The previous considerations show in particular that the restriction of a monochromatic $`_{n,r}(0)`$-module $`P_{[I,i^n]}`$ to $`H_n(0)`$ is isomorphic to the indecomposable projective $`H_n(0)`$-module $`P_I`$. Hence, the restriction of any indecomposable projective module $`P_I`$ given as in Formula (94) is
(96)
$$P_{I_1}\widehat{}\mathrm{}\widehat{}P_{I_k},$$
so that the restriction map $`\pi `$ is given by
(97)
$$\pi (R_{[I,๐]})=R_{I_1}\mathrm{}R_{I_k},$$
which is equivalent to Formula (95).
Continuing the previous example, we see that the restriction of $`P_{2\overline{1}\overline{2}13}`$ to $`H_9(0)`$ is given by
(98)
$$\pi (R_{2\overline{1}\overline{2}13})=\pi (R_2R_{\overline{12}}R_{13})=R_2R_{12}R_{13}=R_{21213}+R_{2133}+R_{3213}+R_{333}.$$
Dually, one can describe the induction of projective $`H_n(0)`$-modules to $`_{n,r}(0)`$. The next result is a consequence of Corollary 6.3.
###### Corollary 6.4.
Let $`I`$ be a composition of $`n`$ and let $`N_I`$ be the $`_{n,r}(0)`$-module induced by the indecomposable projective $`H_n(0)`$-module $`P_I`$. Then
(99)
$$N_IP_{[I,๐]},$$
where the sum runs over all the anticycloribbons of shape $`I`$.
For example, let us complete the case $`I=(2,1)`$ with two colors. The following five anticycloribbons appear in the induction of $`P_I`$, with respective dimensions $`3`$, $`6`$, $`3`$, $`2`$, and $`2`$:
(100) $`2`$ $`1`$ $`1`$ $`2`$ $`1`$ $`2`$ $`1`$ $`1`$ $`2`$ $`2`$ $`2`$ $`2`$ $`1`$ $`1`$ $`1`$
### 6.2. Cartan invariants and decomposition numbers
Finally, we can describe the Cartan invariants and the decomposition matrices by means of the maps
(101)
$$\begin{array}{ccc}\hfill e:& \mathrm{๐๐}^{(r)}& Sym^{(r)}=(Sym)^rSym(X_1,\mathrm{},X_r),\hfill \\ & S_j^{(i)}& h_j(X_i)\hfill \end{array}$$
and
(102)
$$\begin{array}{ccc}\hfill d:& Sym^{(r)}& \mathrm{๐๐๐
}^{(r)}\hfill \\ & h_j(X_i)& F_{[(j),i^j]}.\hfill \end{array}$$
Then the Cartan map is $`c=de`$, and the entry $`c_{[I,๐],[J,๐]}`$ of the Cartan matrix giving the multiplicity of the simple module $`S_{[J,๐]}`$ as a composition factor of the indecomposable projective module $`P_{[I,๐]}`$ is equal to the coefficient of $`F_{[J,๐]}`$ in $`c(R_{[I,๐]})`$.
For an $`r`$-partition $`๐=(\lambda ^{(1)},\mathrm{},\lambda ^{(r)})`$, let $`Z_๐(q)`$ be the irreducible module of the generic algebra $`_{n,r}(q)`$ as constructed in , Theorem 4.1. With our normalization, this module is defined for arbitrary values of $`q`$, including $`q=0`$. For an $`r`$-colored composition $`(I,u)`$ of $`n`$, let
(103)
$$d_{๐,(I,u)}:=[S_{(I,u)}:Z_๐(0)]$$
be the decomposition numbers.
###### Theorem 6.5.
Let $`r_{(I,u)}=e(R_{(I,u)})`$ be the commutative image of $`R_{(I,u)}`$ in $`Sym^{(r)}`$. Let $`S_๐=s_{\lambda ^{(1)}}(X_1)\mathrm{}s_{\lambda ^{(r)}}(X_r)`$ as in . Then the multiplicity of the simple module $`S_{(I,u)}`$ as a composition factor of $`Z_๐(0)`$ is equal to the coefficient of $`S_๐`$ in $`r_{(I,u)}`$, that is
(104)
$$d_{๐,(I,u)}:=S_๐,r_{(I,u)}$$
where $`,`$ is the scalar product on $`Sym^{(r)}`$ for which the basis $`(S_๐)`$ is orthonormal.
###### Proof.
This follows from Shojiโs character formula (Theorem 6.14 of ). We first need to reformulate it on another set of elements. Instead of the $`a_\mu `$ of Shoji, we shall make use of elements $`b_{(J,v)}`$ labelled by colored compositions, and defined by
(105)
$$b_{(J,v)}:=L_{๐(J,v)}T_{w(J)}$$
where $`๐(J,v)=(v_1^{j_1}\mathrm{}v_s^{j_s})`$ and $`w(J)=\gamma _{j_1}\times \mathrm{}\times \gamma _{j_s}๐_{j_1}\times \mathrm{}\times ๐_{j_s}`$, each $`\gamma _k=s_{k1}\mathrm{}s_1`$ being a $`k`$-cycle. With each $`b_{(J,v)}`$, we associate a โnon-commutative cycle indexโ
(106)
$$๐(b_{(J,v)}):=(q1)^{l(J)}S^{(J,v)}((q1)A)\mathrm{๐๐}_n^{(r)}.$$
Then, Shojiโs formula (6.14.1) is equivalent to
(107)
$$\chi _q^๐(b_{(J,v)})=S_๐,๐(b_{(J,v)})$$
where $`S_๐`$ is interpreted as an element of $`QMR^{(r)}`$.
Specializing this at $`q=0`$ gives
(108)
$$\chi _0^๐(b_{(J,v)})=S_๐,(1)^{nl(I)}\mathrm{\Lambda }^{(J,v)}(A),$$
where $`\mathrm{\Lambda }^{(J,v)}(A)=\mathrm{\Lambda }_{j_1}^{(v_1)}\mathrm{}\mathrm{\Lambda }_{j_s}^{(v_s)}`$, and the colored elementary functions $`\mathrm{\Lambda }_j^{(i)}`$ are
(109)
$$\underset{j}{}(1)^j\mathrm{\Lambda }_j^{(i)}=\left(\underset{k}{}S_k^{(i)}\right)^1.$$
On another hand, using the two expressions
(110)
$$\eta _{(I,u)}=L_๐^{}\eta _I^{}=\eta _{I^{\prime \prime }}L_{๐^{\prime \prime }}$$
of the generator of $`S_{(I,u)}`$, one checks that the eigenvalue of $`b_{(J,v)}`$ on $`\eta _{(I,u)}`$ is
(111)
$$F_{(I,u)},๐(b_{(J,v)}).$$
Therefore, $`\chi _0^๐(b_{(J,v)})`$ is also given by
(112)
$$\chi _0^๐(b_{(J,v)})=\underset{(I,u)}{}d_{๐,(I,u)}F_{(I,u)},๐(b_{(J,v)}).$$
The $`๐(b_{(J,v)})`$ being linearly independent, one can conclude that
(113)
$$\underset{(I,u)}{}d_{๐,(I,u)}F_{(I,u)}=S_๐,$$
which is equivalent to Equation (103) since $`r_{(I,u)}`$ is the image by $`e`$ of $`R_{(I,u)}`$ which is itself the dual basis of $`F_{(I,u)}`$.
## 7. Quivers
Being a finite dimensional elementary $``$-algebra, $`_{n,r}(0)`$ can be presented in the form $`Q_{n,r}/`$, where $`Q_{n,r}`$ is a quiver, $`Q_{n,r}`$ its path algebra, and $``$ an ideal contained in $`๐ฅ^2`$ where $`๐ฅ`$ is the ideal generated by the edges of $`Q_{n,r}`$ . The vertices of $`Q_{n,r}`$ are naturally in bijection with the simple modules $`S_{[I,๐]}`$, and the number $`e_{[I,๐],[J,๐^{}]}`$ of edges $`S_{[I,๐]}S_{[J,๐^{}]}`$ is equal to
(114)
$$\mathrm{dim}\mathrm{Ext}^1(S_{[I,๐]},S_{[J,๐^{}]})=[S_{[J,๐^{}]}:\mathrm{rad}P_{[I,๐]}/\mathrm{rad}^2P_{[I,๐]}].$$
that is equal to $`[S_{[I,๐]}:\mathrm{rad}P_{[J,๐^{}]}/\mathrm{rad}^2P_{[J,๐^{}]}]`$ by auto-injectivity.
Therefore, $`e_{[I,๐][J,๐^{}]}=d_{[I,๐][J,๐^{}]}^{(1)}`$, where
(115)
$$d_{[I,๐][J,๐^{}]}^{(k)}:=[S_{[J,๐^{}]}:\mathrm{rad}^kP_{[I,๐]}/\mathrm{rad}^{k+1}P_{[I,๐]}]$$
are the coefficients of the $`q`$-Cartan invariants
(116)
$$d_{[I,๐][J,๐^{}]}(q)=\underset{k0}{}d_{[I,๐][J,๐^{}]}^{(k)}q^k$$
associated with the radical series. Let $`D_n(q)=(d_{[I,๐][J,๐^{}]}(q))`$ be the $`q`$-Cartan matrix. For $`n4`$, Section 8 provides examples of such matrices.
As one will see in the next paragraph, the quiver of $`_{n,r}(0)`$ splits into connected components spanned by evaluation classes of the colors words of the anti-cycloribbons. In other words, $`\mathrm{dim}\mathrm{Ext}^1(S_{[I,๐]},S_{[J,๐^{}]})0`$ implies that $`๐`$ and $`๐^{}`$ are permutations of one another.
The quivers corresponding to evaluations $`(1,3)`$ and $`(2,2)`$ of $`_{4,2}(0)`$ are given on Figures 13 and 14.
The quiver $`Q_{n,r}`$ can be described for all $`n`$: let $`[I,๐]`$ and $`[J,๐^{}]`$ be cycloribbons. Since the simple modules are one-dimensional, isomorphism classes of non trivial extensions
(117)
$$0S_{[J,๐^{}]}MS_{[I,๐]}0$$
are in one-to-one correspondence with isomorphism classes of indecomposable two-dimensional (left) modules $`M`$ of socle $`S_{[J,๐^{}]}`$ and such that $`M/\mathrm{rad}M=S_{[I,๐]}`$.
Let $`M`$ be such a module. Since the $`L_๐`$ are orthogonal idempotents summing to $`1`$, all but one or two of them act on $`M`$ by zero.
If there only is one $`L_๐`$ having a non-zero action on $`M`$, then it acts by $`1`$ and the other $`L_๐`$ act by $`0`$ on $`M`$. So all Lagrange polynomials act as scalars on $`M`$, and $`๐=๐^{}`$. Since $`M`$ is an indecomposable two-dimensional module of $`_{n,r}(0)`$, $`M`$ is also an indecomposable two-dimensional module of $`H_n(0)`$. Thanks to , there exist two elements $`b_1`$ and $`b_2`$ of $`M`$ and an integer $`i`$ such that $`T_ib_1=b_2`$ or $`(1+T_i)b_1=b_2`$, all $`T_j`$ for $`j\{1,\mathrm{},i2\}\{i+2,\mathrm{},n1\}`$ act as scalars, $`T_{i1}`$ does not act as a scalar, and $`T_{i+1}`$ can act as a scalar or not. Moreover, if $`T_ib_1=b_2`$, $`T_{i1}`$ acts by $`1`$ on $`b_1`$ and by $`0`$ on $`b_2`$ and in the opposite way if $`(1+T_i)b_1=b_2`$. Finally, if $`T_{i+1}`$ acts as a non-scalar, then it acts as $`T_{i1}`$.
Since $`M`$ is a module, we must have $`c_{i1}=c_i=c_{i+1}`$ and, it is also equal to $`c_{i+2}`$ if $`T_{i+1}`$ acts as $`T_{i1}`$. Moreover, if $`c_kc_{k+1}`$ then $`T_k`$ acts as a scalar, by $`1`$ if $`c_k<c_{k+1}`$ and by $`0`$ otherwise. One then checks that these relations are sufficient, so that there is an edge in $`Q_{n,r}`$ between $`[I,๐]`$ and $`[J,๐]`$ iff there is an edge between $`I`$ and $`J`$ in the quiver of $`H_n(0)`$.
Let us now consider the case where two elements, $`L_๐`$ and $`L_๐^{}`$ act as non-zero elements on $`M`$. Since they commute, there exist a basis where both are diagonal. Given that their sum is $`1`$ and that they are orthogonal, this proves that there exist two elements $`x`$ and $`y`$ of $`M`$ such that $`x=L_๐x`$ and $`y=L_๐^{}y`$. Since we assumed that $`M`$ is indecomposable, this means that there exists an integer $`i`$ such that, up to the exchange of $`x`$ and $`y`$, $`T_ix=\alpha x+\beta y`$ with $`\beta 0`$. Relations (37) then show that $`๐^{}=๐s_i`$, so that $`c_ic_{i+1}`$ and that $`\alpha =0`$ or $`1`$, so that $`T_ix=y`$, or $`(1+T_i)x=y`$ depending on whether $`c_i>c_{i+1}`$ or $`c_i<c_{i+1}`$. Then, as in , one proves that all $`T_j`$ except $`T_i`$ act diagonally on $`x`$ and $`y`$, and that $`T_j`$ acts as a scalar if $`j\{i1,i,i+1\}`$. Finally, since $`M`$ is a module, the action of any $`T_j`$ ($`ji`$) on $`x`$ (resp. $`y`$) is fixed if $`c_jc_{j+1}`$ (resp. $`c_j^{}c_{j+1}^{}`$). One then checks that these relations are sufficient, so that there is an edge in $`Q_{n,r}`$ between $`[I,๐]`$ and $`[J,๐^{}]`$ ($`๐^{}๐`$) iff there exist an integer $`i`$ such that $`๐^{}=๐s_i`$, and any integer $`k\{i1,i,i+1\}`$ is a descent of both or none of $`I`$ and $`J`$.
So we can conclude:
###### Proposition 7.1.
Let $`[I,๐]`$ and $`[J,๐^{}]`$ be two cycloribbons. There is an edge in $`Q_{n,r}`$ between $`[I,๐]`$ and $`[J,๐]`$ iff
* $`๐=๐^{}`$ and there is an edge in the quiver of $`H_n(0)`$,
* $`๐๐^{}`$ and there exists an integer $`i`$ such that $`๐^{}=๐s_i`$, and if $`k`$ is a descent of both or none of $`I`$, $`J`$ for $`k\{i1,i,i+1\}`$.
Note that there cannot be any edge between two cycloribbons if $`๐`$ and $`๐^{}`$ are not permutations of one another. So $`Q_{n,r}`$ splits into components generated by rearrangement classes of color vectors. These rearrangement classes will be called *evaluations* in the sequel.
Since the product of two basis elements $`B_{๐,\sigma }`$ and $`B_{๐^{},\sigma ^{}}`$ is non-zero only if $`๐`$ and $`๐^{}`$ are permutations of one another, the simple modules $`S_{[I,๐]}`$ occurring in any indecomposable module must be in the same evaluation class, so that the Cartan matrix of any $`_{n,r}(0)`$ is a block matrix, each block being indexed by anticycloribbons with the same evaluation. The monochromatic blocks coincide with the Cartan matrices of $`H_n(0)`$. The other one correspond to blocks (in the sense of indecomposable two-sided ideals) of $`_{n,r}(0)`$.
## 8. Tables
The first tables present the $`q`$-Cartan matrices of the blocks $`(2)`$ and $`(1,1)`$ of $`_{2,2}(0)`$ and the blocks $`(3)`$, $`(2,1)`$, $`(1,2)`$, and $`(1,1,1)`$ of $`_{3,3}(0)`$. Note that the block $`(1,2)`$ can be deduced from the block $`(2,1)`$ by changing any anticycloribbon $`(I,๐)`$ into $`(\overline{I},๐^{})`$ where $`๐^{}`$ is obtained from $`๐`$ by exchanging its smallest letter with its greatest, its second smallest with its second greatest, and so on. Note also that since each permutation is the color word of exactly one anticycloribbon, one can see a standard anticycloribbon as its color word. Then the matrix of the block $`(1,1,1)`$ is the matrix whose entry $`(\sigma ,\tau )`$ is $`q^{l(\sigma ^1\tau )}`$.
Next, we present the blocks $`(4)`$, $`(3,1)`$, $`(2,2)`$, $`(2,1,1)`$, and $`(1,2,1)`$ of $`_{4,4}(0)`$. The other blocks can be deduced from the previous ones: $`(1,1,2)`$ comes from $`(2,1,1)`$ and $`(1,1,1,1)`$ is again $`q^{l(\sigma ^1\tau )}`$.
These $`q`$-Cartan matrices are symmetric along both diagonals since the algebra is auto-injective so that the quiver is an unoriented graph (first main diagonal) and the morphism sending $`L_๐`$ to $`L_{\overline{๐}}`$ and $`T_i`$ to $`1+T_{ni}`$ is an automorphism of $`_{n,r}(0)`$ (second main diagonal). Moreover, for $`q=1`$, the Cartan matrix is the product of the decomposition matrix with its transpose. Finally, we give three examples of decomposition matrices. |
warning/0506/hep-ex0506049.html | ar5iv | text | # ๐-PHYSICS MEASUREMENTS AT THE TEVATRON: ๐ AND ฮโข๐
## 1 Introduction
This article reviews the recent results on $`b`$-physics from the experiments D0 and CDF, which are presently collecting data from the $`p\overline{p}`$ collisions produced at the TeVatron collider with a centre of mass energy of 1.96 TeV. The description of the two experiments and an update on the TeVatron performance are reported elsewhere in these proceedings $`^\mathrm{?}`$. In this paper we review the mass measurements of $`b`$-hadrons, the search for their rare decays and the search for the $`B_s`$ oscillations.
A precise measurement of the mass of $`b`$-hadrons allows testing of the methods used in Lattice QCD and in potential models. At present all but one ground state of $`b`$-mesons foreseen by the quark theory have been detected and for all but two the mass has been experimentally measured with uncertainties below the theoretical uncertainty. The mass of the $`B_c`$ meson was calculated using non-relativistic potential models $`^{\mathrm{?},\mathrm{?}}`$, Lattice QCD $`^\mathrm{?}`$ and perturbative QCD calculations $`^\mathrm{?}`$ with an error that is about two orders of magnitude below the experimental uncertainty $`^\mathrm{?}`$. The CDF collaboration reports here their recent and more precise measurements of the the mass of $`B^0`$,$`B^+`$, $`B_s`$ and $`\mathrm{\Lambda }_b`$. Also reported here is the first evidence of the $`B_c`$ meson decay in a fully reconstructed mode, thus allowing the measurement of the $`B_c`$ mass with a precision comparable to the mass of other $`b`$-mesons.
Both collaborations have updated their upper limit on the branching fractions of $`B_s`$ and $`B^0`$ to muon pairs, which would indicate the presence of new Physics $`^{\mathrm{?},\mathrm{?}}`$ if detected at a level above $`5\times 10^9`$. The D0 collaboration reports the first evidence of the decay $`B_sD_{s1}(2536)\mu +`$ anything. The mass and production yield of excited $`b`$-mesons have also been observed by both experiments.
The measurement of the $`B_s\overline{B}_s`$ oscillation parameters is one of the main physics goals of both experiments at the TeVatron. This paper will review the methods used and report on the recent limits achieved.
## 2 $`b`$ production cross section and trigger
The production cross section of $`b`$ quarks at the TeVatron is approximatively a factor 1000 times larger than at the $`B`$-factories. However, the signal due to $`b`$ physics has to be extracted from a background due to other QCD processes that is 1000 times larger than the signal. This is accomplished using three types of triggers: a trigger based on muon pairs from the decay of a $`J/\psi `$, a trigger based on the detection of a โsoft leptonโ and a trigger based on tracks that are originating from a secondary decay vertex $`^\mathrm{?}`$ that is presently implemented only in the CDF experiment. The data from the $`J/\psi `$ trigger has been used for mass and lifetime measurements, the data from the semileptonic and the secondary vertex triggers have been used for lifetime$`^\mathrm{?}`$ and mixing measurement. Using the data collected with the $`J/\psi `$ trigger the CDF collaboration has measured $`^\mathrm{?}`$ the single $`b`$-quark production cross section integrated over one unit of rapidity:
$$\sigma (pp\overline{b}X,|y|<1)=(29.4\pm 0.6(stat.)\pm 6.2(syst.))\mu b$$
(1)
## 3 $`b`$-hadron masses
The $`b`$-hadron masses have been measured by reconstructing decays containing a muon pair from a $`J/\psi `$. The two experiments have complementary design features: CDF has a better mass resolution, while D0 has a larger angular acceptance for muons. The spectrometer of the CDF detector has been calibrated using the $`J/\psi `$ mass as a reference by correcting for passive material effect, to eliminate the $`p_T`$ dependence of the mass. The actual value of the magnetic field has been tuned to obtain the world average $`J/\psi `$ mass value $`^\mathrm{?}`$. The calibration has been checked against the $`\mathrm{{\rm Y}}`$ mass. The high statistics of $`J/\psi \mu ^+\mu ^{}`$ decays available has allowed the CDF experiment to reduce the systematic uncertainties to sub-MeV values. Using the following decay modes and a luminosity of 220 pb<sup>-1</sup> CDF has measured the following preliminary values for hadron masses, all in MeV/$`c^2`$:
## 4 Evidence of the decay $`B_c^\pm J/\psi \pi ^\pm `$ and $`B_c`$ mass measurement
The $`B_c`$ meson has been observed $`^\mathrm{?}`$ at the TeVatron Run I by the CDF collaboration in decays containing a neutrino, and therefore its mass has a large experimental uncertainty. The D0 collaboration has recently observed the semileptonic decays of the $`B_c`$ meson in the Run 2 data $`^\mathrm{?}`$. The LEP experiments searched for the fully reconstructed decays of the $`B_c`$ meson $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$, but the $`B_c`$ production cross section from $`Z`$ decays was too low for detecting these decays. The CDF collaboration has now evidence for the fully reconstructed decay mode $`B_c^\pm J/\psi \pi ^\pm `$, with $`J/\psi \mu \mu `$. This simple decay mode has a relatively large expected branching ratio $`^\mathrm{?}`$ when the daughter decays are also included. It is detected with a trigger that is relatively efficient and is not based on the decay vertex. The CDF collaboration has used an analysis technique that blinded a wide range of the mass distribution during cut optimization. The statistical test that was used to assess the significance of the signal was completely set before unblinding the mass distribution. The test was based on the score function $`\mathrm{\Sigma }=N_s/(1.5+\sqrt{N_b})`$ where $`N_s`$ and $`N_b`$ are the number of signal and backgorund candidates as obtained from a fit in a sliding mass window that was 300 MeV/c<sup>2</sup> wide. As the mass value was known with an uncertainty of $`\pm 400`$ MeV/c<sup>2</sup> the mass peak corresponding to this decay was searched for in the range $`5700M(B_c)7000`$ MeV/c<sup>2</sup>. The analysis required a good fit to a displaced vertex using well defined tracks with silicon hits, making use of the innermost silicon layer (L00). Upon unblinding the mass distribution, it was found that one region contained a peak that satisfied the predefined statistical test. The probability that the background generates a fluctuation equal or larger than the observed signal, anywhere in the search region, was estimated using Monte Carlo generated distributions that simulated only the background. This probability was found a posteriori to be about 0.27%. The experimental evidence of this decay mode allowed measuring the $`B_c`$ mass. The $`J/\psi \pi `$ mass distribution is shown in fig. 1. The unbinned likelihood fit returns $`18.9\pm 5.4`$ $`B_c`$ candidates on a background of $`10.1\pm 1.4`$ events. The main contribution to the systematic uncertainties comes from fitting with different background shapes. Other systematic uncertainties are derived from the mass measurements reported above. The experimental value of the $`B_c`$ mass is
$$m(B_c)=6287.4\pm 4.5(stat.)\pm 1.1(syst.)\text{ MeV/c}^2$$
(2)
The details of this analysis will be published in a forthcoming paper $`^\mathrm{?}`$. This result is in very good agreement with the theoretical expectations mentioned above $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$.
Checks on the detection of the partially reconstructed decays $`B_cJ/\psi +`$ track $`+`$ anything in the same sample also gave a positive result, with a significant excess only for $`m(J/\psi \pi )<m(B_c)`$, in agreement with the expectations.
## 5 Excited B mesons
The excited states of the $`b`$-mesons are of great interest to measure their mass, branching ratios and width, which are predicted with a good precision by various theoretical models. In addition, at the hadron collider the excited states of the $`b`$ mesons must be taken into account in mixing measurements, when using a flavour tagging method that is based on the information from tracks that are near to the candidate $`b`$ meson. The excited B states had been studied at LEP by the Aleph collaboration $`^\mathrm{?}`$. The D0 collaboration has measured the mass of the two narrow states corresponding to the orbital angular momentum L=1 $`B_1`$ and $`B_2^{}`$ in their decay to $`B^{()}\pi `$. Also the CDF collaboration confirms evidence of two separate peaks. In both experiments fully reconstructed decays of the $`B^+`$ and $`B^0`$ mesons are detected and are associated with a track from the primary vertex. The three contributions are from the decays $`B_1B^{}\pi `$, $`B_2^{}B^{}\pi `$ and $`B_2^{}B\pi `$, with $`B^{}B\gamma `$. The undetected photon of about 46 MeV causes the larger shift in the mass peak. The D0 preliminary results are $`^\mathrm{?}`$:
$$m(B_1)=5274\pm 4(stat.)\pm 7(syst.)MeV/c^2$$
(3)
$$m(B_2^{})m(B_1)=23.6\pm 7.7(stat)\pm 3.9(syst)MeV/c^2$$
(4)
$$\mathrm{\Gamma }_1=\mathrm{\Gamma }_2=23\pm 12(stat)\pm 9(syst)MeV$$
(5)
## 6 Search for rare decays
The decay of a neutral $`b`$-meson into two opposite charged muons is mediated in the standard model by flavour changing neutral currents due to ladder diagrams and the corresponding branching ratio is calculated to be $`^{\mathrm{?},\mathrm{?}}`$ $`BR(B_{d,s}\mu ^+\mu ^{})=(3.4\pm 0.54)\times 10^9`$. A branching ratio that is larger than this value would indicate the presence of new Physics processes, that can considerably enhance this decay channel. The D0 collaboration has searched for the decay of $`b`$-mesons to two muons using data from an integrated luminosity of 300 pb<sup>-1</sup>. Four candidates were found in the mass window, giving a limit on this branching ratio $`BR(B_{d,s}\mu ^+\mu ^{})3.7\times 10^7`$ at 95% conficence level.
The CDF collaboration uses its better detector resolution to put separate limits on the decay rate of the two mesons $`^\mathrm{?}`$: $`BR(B_s\mu ^+\mu ^{})5.8\times 10^7`$ and $`BR(B_d\mu ^+\mu ^{})1.5\times 10^7`$ at 90% confidence level using a collected luminosity of 171 pb<sup>-1</sup>. This result has been updated shortly after this conference. The D0 collaboration has also measured the sensitivity to detect a signal from the decay $`B_s\mu \mu \varphi `$. Using data from 300 pb<sup>-1</sup> they can put a limit on the branching fraction $`BR1.2\times 10^5`$ at 95% CL.
## 7 Flavour oscillation measurements
To perform time-dependent mixing measurements at the hadron collider it is necessary to select a flavourโspecific decay signal (e.g. $`B_sD_s\pi ^{}`$), establish the flavour of the $`b`$-meson at production and measure precisely the proper decay time, which involves measuring the decay length and the momentum of the reconstructed meson. The relevant parameters are the purity of the sample ($`f_{sig}=N_{signal}/N_{background}`$), which indicates the fraction of decay signal in the sample, the tagging efficiency $`ฯต`$, which indicates the fraction of the signal sample that has a flavour tag, $`ฯต=\frac{N_{tag}}{N_{cand}}`$, and the dilution $`D`$, which indicates the probability that the tag is correct: $`D=\frac{N_RN_W}{N_R+N_W}`$ The statistical power of the sample is diluted by a factor $`ฯตD^2`$: the significance of the oscillation signal is given by $`^\mathrm{?}`$:
$$Sf_{sig}\sqrt{1/2ฯตD^2N_{cand}}e^{\frac{1}{2}(\mathrm{\Delta }m_s\sigma _t)^2}$$
(6)
where $`\mathrm{\Delta }_m`$ is the mixing parameter that determines the frequency of the oscillations and $`\sigma _t`$ is the proper time resolution. The latter is the sum in quadrature of two terms: one related to the spatial resolution of the secondary vertex and the second related to the momentum resolution of the reconstructed meson.
For $`B_s`$ mixing we have a lower statistics compared to the $`B^0`$ channel and a larger oscillation parameter, which implies faster oscillations. Therefore, to have a large significance for large $`\mathrm{\Delta }m_s`$ we need a very good proper time resolution, which can be better achieved with fully reconstructed decays, at the price of smaller statistics.
Both collaborations have used, until now, only the โopposite sideโ tagging methods, which rely on tagging the flavour of a $`b`$-meson by looking at the characteristics of particles that are produced โawayโ from the reconstructed meson and that presumably contain hadrons from the $`b`$-quark that was produced in association with the one that originated the reconstructed $`b`$-meson. This tagging method is independent on the nature of the meson under study, so it has been tested on the $`B^0`$ mixing and applied to the search for $`B_s`$ mixing.
## 8 $`B^0`$ Mixing
The D0 collaboration has used only the semileptonic decays $`B^0D^\pm \mu ^{}X`$, (Here $`X`$ denotes any set of particles produced in the decay). The CDF collaboration has analysed some fully reconstructed decays from the vertex trigger, for which the lifetime bias is now well understood $`^\mathrm{?}`$. Both collaborations have results that are in very good agreement with the precise measurements, which have been performed at the $`B`$-factories, as shown in Table 1. The purpose of these measurement was mainly to test the fit algorithms and to obtain the dilution factors, to be used in the search for $`B_s`$ oscillations.
## 9 $`B_s`$ Mixing
The D0 collaboration has applied an updated version of the opposite side muon tagging algorithm to the signal sample enriched in $`B_s`$ semileptonic decays $`B_sD_s\mu X`$, with $`D_s^\pm \varphi \pi ^\pm `$. Using $`376\pm 31`$ reconstructed and tagged decays (on a total of 7037), and with a tagging dilution of $`D=0.552\pm 0.016`$ a very preliminary fit gives a null result on oscillations, with a limit to the parameter $`\mathrm{\Delta }m_s5.0`$ ps<sup>-1</sup> at 95% confidence level and a sensitivity of 4.6 ps<sup>-1</sup>, using both the statistic and the systematic uncertainties. This result is not as good as the world average, that indicates $`\mathrm{\Delta }m_s14.5`$ ps<sup>-1</sup> but it is only the first attempt to fit the data. A considerable improvement is expected, both on statistics and on tagging efficiency. In particular, using the same-side-tagging techniques the tagged sample can be increased using the same data collected to date.
By comparison, the CDF collaboration has obtained similar results in terms of $`ฯตD^2`$. The CDF sample has about 7500 candidates to $`B_s`$ semileptonic decays, in three reconstructed decay modes of the $`D_s`$: $`\varphi \pi ^+`$, $`K^0K^+`$ and $`\pi ^+\pi ^{}\pi ^+`$. For fullyโreconstructed hadronic decay candidates $`B_sD_s\pi ^{}`$, with the $`D_s`$ decaying to the same modes as above, the statistics is $`900`$ candidates; the expected tagging power $`ฯตD^2`$ factors are reported in Table 1.
## 10 Other results
The D0 collaboration has found evidence for the decay $`B_sD_{s1}(2536)\mu +X`$, heading to a measurement of this branching ratio and of the mass of the $`D_{s1}(2536)`$.
## 11 Conclusions and prospectives
The experiment at the Tevatron have considerably improved the precision on $`b`$-mesons mass measurements. In particular, the $`B_c`$ mass has been measured with good precision thanks to the detection of a signal in a fully reconstructed mode, which is reported here for the first time. Also the masses of two B-mesons excited states has been measured. The mixing parameter $`\mathrm{\Delta }m_d`$ is in good agreement with the world average, but an initial look at the $`B_s`$ oscillations has given a null result on the measurement of $`\mathrm{\Delta }m_s`$, with a limit that is still more than a factor of two lower than the world average. Large improvements are expected not only with more statistics, but also with improved techniques in tagging, improving vertex resolution and adding other decay channels.
## References |
warning/0506/hep-ph0506302.html | ar5iv | text | # Synthesis of DGLAP and total resummation of leading logarithms for the non-singlet spin structure function ๐โ
## I Introduction
The non-singlet component of the spin structure function $`g_1`$ have been investigated in great detail in deep inelastic scattering (DIS) experiments. The standard theoretical framework for studying the DIS structure functions is provided by DGLAPdglap . In this approach, $`g_1^{NS}(x,Q^2)`$ can be represented as a convolution of the coefficient functions and the evolved quark distributions. Combining these results with appropriate fits for the initial quark distributions, provides a good agreement with the available experimental data.
However, the DGLAP evolution eqs. were originally applied in a range of large $`x`$ values, where higher-loop contributions to the coefficient functions and the anomalous dimensions are small. Such corrections are becoming essential when $`x`$ is decreasing, so DGLAP should not work so well at $`x1`$. Nevertheless, DGLAP predictions are in a good agreement with available experimental data. It leads to the conclusion that the impact of the higher-order corrections is negligibly small for the available values of $`x`$. Below we use our results egt1 to show that the impact of the high-order corrections on the $`Q^2`$ and $`x`$ -evolutions of the non-singlet structure functions is quite sizable and bounds the region of strict applicability of DGLAP to $`x>10^2`$. We also show that the reason for the success of DGLAP at $`x<10^2`$ is related to the sharp $`x`$ -dependence assumed for the initial parton densities, which is able to mimic the role of high-order corrections.
The paper is organized as follows: In Sect. 2 we discuss the difference of our approach with DGLAP. Then we compare our and the DGLAP formulae for asymptotics of $`g_1`$. In Sect. 3 we suggest a method to combine DGLAP with our approach in order to obtain equally correct expressions for both large and small values of $`x`$. Sect. 4 contains our conclusions.
## II Comparison of DGLAP and our approach
As the DGLAP -expressions for the non-singlet structure functions are well-known, we discuss them briefly only. In this approach, $`g_{1DGLAP}^{NS}(x,Q^2)`$ can be represented as a convolution
$$g_{1DGLAP}^{NS}(x,Q^2)=_x^1\frac{dy}{y}C(x/y)\mathrm{\Delta }q(y,Q^2)$$
(1)
of the coefficient functions $`C(x)`$ and the evolved quark distributions $`\mathrm{\Delta }q(x,Q^2)`$. Similarly, $`\mathrm{\Delta }q(x,Q^2)`$ can be expressed through the convolution of the splitting functions and the initial quark densities $`\delta q(x1,Q^2\mu ^2)`$ where $`\mu ^2`$ is the starting point of the $`Q^2`$ -evolution. It is convenient to represent $`f(x,Q^2)`$ in the integral form, using the Mellin transform:
$$g_{1DGLAP}^{NS}(x,Q^2)=(e_q^2/2)_ฤฑ\mathrm{}^ฤฑ\mathrm{}\frac{d\omega }{2ฤฑ\pi }(1/x)^\omega C(\omega )\delta q(\omega )\mathrm{exp}\left[\gamma (\omega )_{\mu ^2}^{Q^2}\frac{dk_{}^2}{k_{}^2}\alpha _s(k_{}^2)\right]$$
(2)
where $`C(\omega )`$ are the non-singlet coefficient functions, $`\gamma (\omega )`$ the non-singlet anomalous dimensions and $`\delta q(\omega )`$ the Mellin transforms of the initial non-singlet quark densities. The standard DGLAP fits $`\delta q(x)`$ for the non-singlet parton densities (see e.g. Refs. a ; v ) consist of the terms singular when $`x0`$ and the regular in $`x`$ part. For example, the fit A of Ref.a is chosen as follows:
$`\delta q(x)=N\eta x^\alpha \varphi (x),`$ (3)
$`\varphi (x)(1x)^\beta (1+\gamma x^\delta ),`$
with $`N,\eta `$ being the normalization, $`\alpha =0.576`$, $`\beta =2.67`$, $`\gamma =34.36`$ and $`\delta =0.75`$. As the term $`x^\alpha `$ in the rhs of Eq. (3) is singular when $`x0`$ whereas the second one, $`\varphi (x)`$ is regular, we will address them as the singular and regular parts of the fit respectively. Obviously, in the $`\omega `$ -space Eq. (3) is a sum of the pole contributions:
$$\delta q(\omega )=N\eta \left[(\omega \alpha )^1+\underset{k=1}{\overset{\mathrm{}}{}}m_k\left((\omega +k\alpha )^1+\gamma (\omega +k+1\alpha )^1\right)\right],$$
(4)
with $`m_k=\beta (\beta 1)..(\beta k+1)/k!`$, so that the first term in Eq. (4) (the leading pole) corresponds to the singular term $`x^\alpha `$ of Eq. (3) and the second term, i.e. the sum of the poles, corresponds to the interference between the singular and regular terms. In contrast to the leading pole position $`\omega =\alpha `$, all other poles in Eq. (4) have negative values because $`k\alpha >0`$. An alternative approach was used in Refs. ber, by introducing and solving infrared evolution equations with fixed $`\alpha _s`$. This approach was improved in Refs. egt1 , where single-logarithmic contributions were also accounted for and the QCD coupling was regarded as running in all Feynman graphs contributing to the non-singlet structure functions. In contrast to the DGLAP parametrization $`\alpha _s=\alpha _s(k_{}^2)`$, we used in Refs. egt1 another parametrization where the argument of $`\alpha _s`$ in the quark ladders is given by the time-like virtualities of the intermediate gluons. Refs. egt1 suggest the following formulae for the non-singlet structure functions:
$$g_1^{NS}(x,Q^2)=(e_q^2/2)_ฤฑ\mathrm{}^ฤฑ\mathrm{}\frac{d\omega }{2\pi ฤฑ}(1/x)^\omega C_{NS}(\omega )\delta q(\omega )\mathrm{exp}\left(H_{NS}(\omega )y\right),$$
(5)
with $`y=\mathrm{ln}(Q^2/\mu ^2)`$ so that $`\mu ^2`$ is the starting point of the $`Q^2`$ -evolution. The new coefficient function $`C_{NS}`$ are expressed in terms of new anomalous dimensions $`H_{NS}`$ whereas $`H_{NS}`$ account for the total resummation of the double- and single- logarithmic contributions (see Ref. egt1 for details).
## III Comparison of DGLAP and our small-$`x`$ asymptotics
When $`x0`$, one can use the saddle point method in order to estimate the integrals in Eq. (5) and derive much simpler expressions for the non-singlet structure functions:
$$g_1^{NS}e_q^2\delta q(\omega _0)\xi ^{\omega _0},$$
(6)
with $`\xi =\sqrt{Q^2/(x^2\mu ^2)}`$ and with the intercept $`\omega _0=0.42`$. Eq. (6) predicts the asymptotic scaling for the non-singlet structure functions: Asymptotically, $`g_1^{NS}`$ depends on one argument $`\xi `$ instead of depending on $`x`$ and $`Q^2`$ separately.
When the standard DGLAP fits, e.g. the fit of Eq. (3), are used, the asymptotics of $`g_{1DGLAP}^{NS}(x,Q^2)`$ is also the Regge-like:
$$g_{1DGLAP}^{NS}(e_q^2/2)C(\alpha )(1/x)^\alpha \left((\mathrm{ln}(Q^2/\mathrm{\Lambda }^2))/(\mathrm{ln}(\mu ^2/\mathrm{\Lambda }^2))\right)^{\gamma (\alpha )/b},$$
(7)
with $`b=(332n_f)/12\pi `$.
Comparison of Eq. (6)and Eq. (7) demonstrates that both DGLAP and our approach lead to the Regge asymptotic behavior in $`x`$. However, it is important that our intercept $`\omega _0`$ is obtained by the total resummation of the leading logarithmic contributions and without any assumption about fits for $`\delta q`$ whereas the DGLAP intercept $`\alpha `$ in Eq. (7) is generated by the phenomenological factor $`x^{0.57}`$ of Eq. (3) which mimics the total resummation. In other words, the impact of the higher-loop radiative corrections on the small-$`x`$ behavior of the non-singlets is, actually, incorporated into DGLAP phenomenologically, through the fits. It means that the singular factors can be dropped from such fits when the coefficient function includes the total resummation of the leading logarithms and therefore in this case fits for $`\delta q`$ can be chosen as regular functions of $`x`$.
## IV Combining DGLAP with our higher-loop contributions
Eq. (5) accounts for the resummation of the double- and single logarithmic contributions to the non-singlet anomalous dimensions and the coefficient functions that are leading when $`x`$ is small. However, the method we have used does not allow us to account for other contributions which can be neglected for $`x`$ small but become quite important when $`x`$ is not far from 1. On the other hand, such contributions are naturally included in DGLAP, where the non-singlet coefficient function $`C_{DGLAP}`$ and anomalous dimension $`\gamma _{DGLAP}`$ are known with the two-loop accuracy:
$`C_{DGLAP}=1+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}C^{(1)},`$ (8)
$`\gamma _{DGLAP}={\displaystyle \frac{\alpha _s(Q^2)}{4\pi }}\gamma ^{(0)}+\left({\displaystyle \frac{\alpha _s(Q^2)}{4\pi }}\right)^2\gamma ^{(1)}`$
Therefore, we can borrow from the DGLAP formulae the contributions which are missing in Eq. (5) by adding $`C_{DGLAP}`$ and $`\gamma _{DGLAP}`$ to the coefficient function and anomalous dimension of Eq. (5). It is important to avoid a double counting DL and SL terms common for these expressions.
In order to do so, let us consider the region of $`x1`$ where the effective values of $`\omega `$ in Eqs. (2,5) are large. In this region we can expand $`H_{NS}`$ and $`C_{NS}`$ into a series in $`1/\omega `$. Retaining the first two terms in each series, we arrive at $`C_{NS}=\stackrel{~}{C}_{NS}+O(\alpha _s^2)`$, $`H_{NS}=\stackrel{~}{H}_{NS}+O(\alpha _s^3)`$, with (see Ref. egt1 for details)
$`\stackrel{~}{C}_{NS}=1+{\displaystyle \frac{A(\omega )C_F}{2\pi }}\left[1/\omega ^2+1/2\omega \right],`$ (9)
$`\stackrel{~}{H}_{NS}={\displaystyle \frac{A(\omega )C_F}{4\pi }}\left[2/\omega +1\right]+\left({\displaystyle \frac{A(\omega )C_F}{4\pi }}\right)^2(1/\omega )\left[2/\omega +1\right]^2+D[1/\omega +1/2].`$
Now let us define the new coefficient functions $`\widehat{C}_{NS}`$ and new anomalous dimensions $`\widehat{C}_{NS}`$ as follows:
$`\widehat{H}_{NS}=\left[H_{NS}\stackrel{~}{H}_{NS}\right]+{\displaystyle \frac{A(\omega )}{4\pi }}\gamma ^{(0)}+\left({\displaystyle \frac{A(\omega )}{4\pi }}\right)^2\gamma ^{(1)},`$ (10)
$`\widehat{C}_{NS}=\left[C_{NS}^{(\pm )}\stackrel{~}{C}_{NS}\right]+1+{\displaystyle \frac{A(\omega )}{2\pi }}C^{(1)}.`$
These new, โsyntheticโ coefficient functions and anomalous dimensions of Eq. (10) include both the total resummation of the leading contributions and the DGLAP expressions in which $`\alpha _s(Q^2)`$ is replaced by $`A(\omega )`$ defined in Refs. egt1 because the factorization of the phase space into transverse and longitudinal spaces used in DGLAP to parametrize $`alpha_s`$ is a good approximation for large $`x`$ only. |
warning/0506/physics0506048.html | ar5iv | text | # Microsolvation of Li+ in small He Clusters. Li+Hen species from Classical and Quantum Calculations
## I Introduction
The study of the nanoscopic forces which act within small aggregates of weakly interacting particles, especially within assemblies containing helium atoms, has received a great deal of attention in the last few years 1 ; 2 ; 3 because of the broad range of phenomena that can be probed under the very special conditions provided by the He nanodroplets as containers of atoms or molecules. They are indeed, ultracold homogeneous matrices where the corresponding spectra often reach very high resolution due to the superfluid properties of the helium droplet 4 .
The possibility of causing electronic excitation and/or ionization of the dopant species can offer additional ways of probing the modified interaction between the new ionic species and the gentle matrix of the helium droplets 3 since the ensuing distribution surrounding the molecular impurity is usually markedly deformed as a consequence of the electrostriction effects on the solvent brought on by the induction forces between the charged dopant and the helium atoms in the droplet 5a ; 5 . The analysis of such effects in the case of lithium-containing impurities has been particularly intriguing because of the expected simplicity of the electronic structures involved and, at the same time, the unusual bonding behaviour of such systems. Furthermore, experiments in which the droplet was ionized after capture of Li atoms 7 have revealed a wealth of newly formed species like Li<sup>+</sup>, Li$`{}_{2}{}^{}{}_{}{}^{+}`$ and LiHe<sup>+</sup> which are produced during droplet fragmentation and evaporation thereby triggering the corresponding analysis of structure and bonding behaviour from theory and computations.
In a previous computational investigation 8 on the interaction of Li<sup>+</sup> with $`n`$ helium atoms, with $`n`$ varying from 1 to 6, we have shown, in fact, that the energy optimized structures of the clusters were largely determined by two-body (2B) forces, with the many-body effects being fairly negligible for determining the final geometry of those small aggregates. In the present study we have therefore decided to analyze in greater detail the quantum states and binding energies of the larger structures of Li<sup>+</sup> impurities within the helium droplets in order to extend our computational knowledge on such species.
The next Section briefly describes the LiHe<sup>+</sup> Potential Energy Curve (PEC) used here. The results for the trimer ground state are given and discussed in Section III, where we report the features obtained using the distributed Gaussian functions (DGF) method 11 ; 12 ; 13 ; 14 and the quantum Diffusion Monte Carlo (DMC) method 9 , in comparison with the values obtained via a classical minimization procedure. Section IV extends the discussion to the larger species, whose energetics and structural features, both obtained with DMC and classical methods, are analyzed.
## II The pairwise potentials
The actual potential energy curve (PEC) for Li<sup>+</sup>He has been studied before because of its interest in modeling low-energy plasmas 16 . Alrich and coworkers suggested earlier on a model potential as a variant of the Tang-Toennies model 17 while recent calculations of Soldรกn et al. 19 employed a CCSD(T) treatment and used the aug-cc-pV5Z basis set. Recently, we have also carried out a new set of calculations on these clusters 8 using the MP4 method and the cc-pV-5Z basis set. The results from the ab initio calculation of Ref. 19 , the potential obtained from a MP4/cc-pV-5Z calculation and the model potential of Ref. 17 are presented in pictorial form by the two panels of Fig. 1 where in the left panel we show the region of the potential minimum. One clearly sees there that the deepest well is presented by the CCSD(T) calculations of 19 , while the model potential of 17 has the most shallow well. The corresponding long-range part of the PECโs is reported on the right-hand side panel, where one can see that the calculations of 19 closely follow the long-range dipole polarisability tail, computed here with $`\alpha _{He}`$=1.3832 $`a_0^3`$ 20 . In Table 1 we report the bound states of the Li<sup>+</sup>-He diatom when $`J=0`$ for the three different potentials calculated using a DVR scheme dvr using 2000 grid points in the interval 1.5, 200 a.u.. In the case of the potential of Murrel et al. 19 our results are almost coincident with those provided by the same authors. As can be seen there, the three potentials yield rather similar energies, although in the well region the agreement is better between those from the MP4 potential and those from Soldan et al.. On the other hand, the latter has a long range behavior very similar to that from the semi-empirical ATT potential (see also Fig. 1 where we compare it to the โexperimentalโ long range behavior): both these potentials, indeed, support 8 bound states while the potential calculated by us with MP4 8 yields only 7 bound states for $`J=0`$. The MP4 potential therefore seems not to provide a good description of the long range, inductive tail because of the lack of augmented atomic functions in the basis set from which it has been obtained. Hence, we have decided to employ in the following the CCSD(T) PEC from Soldan et al..
For the He<sub>2</sub> dimer we have employed the semiempirical potential (LM2M2) of Aziz and Slaman 24 : in Fig. 2 we sketch the two pair potentials used in our calculations. The two potentials are completely different: the Li<sup>+</sup>-He one is dominated by electrostatic and induction forces while the He-He interaction is much weaker and purely van der Waals in nature.
As we mentioned above, the full interaction that acts in the M<sup>+</sup>Rg<sub>n</sub> clusters can be approximated by a sum of pairwise potentials. The main source of error in this model is due to the absence of the repulsive interactions between the induced dipoles on the helium atoms, especially those that are located near the ionic center. These interactions, however, can be taken to be very small for rare gas atoms and especially for helium partners. For example, in Ref. 29 for the analogous situation of the anionic dopant Cl<sup>-</sup> in Ar clusters it has been shown that the inclusion of the leading terms of such 3-body forces does not alter substantially the energetics and the structure of the clusters: in ref. 8 we have discussed and confirmed this feature for the system analyzed here. Furthermore, since the charge delocalization in Li<sup>+</sup>-He<sub>n</sub> is very small 28 and the interaction is mostly determined by ยจphysicalยจ forces (charge/induced multipole) rather than by ยจchemicalยจ effects, all the helium atoms that are attached directly to the ionic center can be considered to be equivalent because there is no detectable tendency for the latter to preferentially form chemically bound structures with only a subset of such adatoms, a fact which would therefore differentiate some of the binding effects with respect to those in the rest of the cluster.
## III An outline of the computational tools
### III.1 The classical optimization
The total potential in each cluster is described by the sum of pairwise potentials and searching for the global minimum in this hypersurface $`V_{TOT}`$ will give us the lowest energy structure for each aggregate from a classical picture viewpoint. All the classical minimizations were carried out with a modified version of the OPTIM code by Wales wales94 . This code is based on the eigenvector following method. The basis of this method is the introduction of an additional Lagrange parameter into a โtraditionalโ optimization framework baker86 , which seeks to simultaneously minimize in all directions. All searches were conducted in Cartesian coordinates using projection operators to remove overall translation and rotation following Baker and Hehre baker91 . Analytic first and second derivatives of the energy were employed at every step, and the resulting stationary point energies and geometries are essentially exact for the model potential in question. The details of the method as applied in our group have been given before and therefore will not be repeated here.
### III.2 The quantum stochastic calculation (DMC)
The Diffusion MonteCarlo Method has been extensively discussed in a number of papers (Refs. hammond ; ceperley ; suhm and references therein) and therefore will not be repeated here.
In our implementation a random walk technique is used to solve the diffusion equation where a large number of random walkers is propagated with time steps $`\mathrm{\Delta }\tau `$ starting from an arbitrarily chosen initial distribution. The ground state energy $`E_0`$ is obtained by averaging $`E_L(๐ซ)`$ over the final mixed distributions $`f(๐ซ,t)=\psi _T(๐ซ,t)\psi (๐ซ,t)`$:
$$<E_L>=\frac{E_L(๐ซ)f(๐ซ,\tau _f)๐๐ซ}{f(๐ซ,\tau _f)๐๐ซ}=\frac{\psi _0(๐ซ)\widehat{H}\psi _T(๐ซ)๐๐ซ}{\psi _0(๐ซ)\psi _T(๐ซ)๐๐ซ}=E_0$$
(1)
The energy is therefore affected by a bias due to the use of a mixed distribution. The bias is however minimized by using very long propagation time and very short timesteps. Expectation values of position operators $`\widehat{A}(๐ซ)`$ are also given by averaging over $`f(๐ซ)`$, a procedure that leads to biased distributions. However, we believe that the bias does not modify the qualitative picture of the present calculations especially because whenever our spatial distributions are compared with those from other quantum calculations (where possible) and with classically minimized structures, they are found to be in remarkable agreement with each other as we shall discuss further below.
The trial function used here for the $`\mathrm{He}\mathrm{He}`$ pairs is a product of trial wavefunction:
$$\mathrm{\Psi }_T=\underset{i,j\mathrm{He}}{}\mathrm{exp}\left(\frac{p_5}{R_{ij}^5}\frac{p_2}{R_{ij}^3}p_0\mathrm{log}R_{ij}p_1R_{ij}\right)$$
(2)
where the values of the coefficients have been taken by Ref lewerenz97 . The trial function for the $`\mathrm{Li}^+\mathrm{He}`$ pair has been chosen as a gaussian function centered around the energy minimum of the relative interaction. However, in order to adjust each trial wavefunction to the size of the larger clusters, its parameters have been chosen so that they make it increasingly more delocalized as the number of adatoms is increased. All prameters are available from direct requests to the corresponding author.
### III.3 The distributed Gaussian Functions (DGF) expansion
The DGF method 11 ; 12 is a variational approach which solves the trimer bound state problem written in terms of the atom-atom coordinates by employing a large set of distributed Gaussian functions as a basis set.
First of all, the Hamiltonian for a triatomic system is expressed in terms of atom pair coordinates R<sub>1</sub>, R<sub>2</sub> and R<sub>3</sub>, i.e. in terms of the distances between each pair of atoms along which the Gaussian functions are distributed (see 13 for the expression of the Hamiltonian for a triatomic system with an impurity).
The total potential for the trimer is assumed to be the sum of the three two-Body potentials as discussed previously and the calculations are carried out for a zero total angular momentum. The total wavefunction for the v-th vibrational states is then expanded in terms of symmetrized basis functions
$$\mathrm{\Phi }_v(R_1,R_2,R_3)=\underset{j}{}a_j^{(v)}\varphi _j(R_1,R_2,R_3)$$
(3)
with
$$\varphi _j(R_1,R_2,R_3)=N_{lmn}^{1/2}\underset{PS_2}{}P[\phi _l(R_1)\phi _m(R_2)]\phi _n(R_3)$$
(4)
for the two-identical-particle system. Here, $`j`$ denotes a collective index such as $`j=(lm;n)`$ and N<sub>lmn</sub> is a normalization constant expressed in term of the overlap integrals s$`{}_{pq}{}^{}=\phi _p|\phi _q`$. Each one-dimensional function $`\phi _p`$ is chosen to be a DGF (light86 ) centered at the $`R_p`$ position
$$\phi _p(R_i)=\sqrt[4]{\frac{2A_p}{\pi }}e^{A_p(R_iR_p)^2}.$$
(5)
With the DGF approach we can obtain several indicators on the spatial behaviour of the bound states of the systems (the root mean value of the square area, the average of the cosine value - and the various moments - of any angle etc.), along with several probability distribution functions like the pair distribution function
$`D^{(v)}(R_1)`$ $`=`$ $`{\displaystyle \mathrm{\Phi }_v(R_1,R_2,R_3)^2๐R_2๐R_3}.`$ (6)
The selection of a suitable set of Gaussian functions and their distribution within the physical space where the bound states are located is obviously of primary importance in order to finally obtain converged and stable results. We extensively experimented with different sets obtained by changing the number and location of the DGF depending on the features of the 2B potentials employed. The details of the basis set employed for the title system are specifically given in the following section.
## IV The $`\mathrm{Li}^+\mathrm{He}_2`$ trimer
Together with the classical optimization and the DMC calculations, we further carried out the analysis of the properties of the trimer Li<sup>+</sup>He<sub>2</sub> using the DGF method. We employed an optimized DGF basis set in order to obtain well-behaving total wave functions at the triangle inequality boundaries (for further details see Ref. isa04 ). The Gaussian functions are distributed equidistantly along the atom-atom coordinates starting from 2.55 a<sub>0</sub> and out to 9.18 a<sub>0</sub> with a step of 0.17 a<sub>0</sub>, a choice which ensures converged results within about 0.1 cm<sup>-1</sup> (0.01 % of the total ground state energy). In Table II we report the results obtained with DMC and DGF methods: they are seen to be in good agreement with each other. Due to the addition of the second light He atom, the Zero Point Energy (ZPE) of the trimer is slightly higher than the ZPE of the isolated dimer ion (see Table I). However, the two ZPE values are very similar in percentage values, showing that the Li<sup>+</sup> impurity strongly affects the features of the cluster, as expected from the involved potentials while the additional helium atom has little effect on the bonding features. The ionic system, as expected, does not present the typical high degree of delocalization shown in pure He clusters (the ZPE for He<sub>3</sub> is more than 99 % of its total well depth 10 ), or by doped He aggregates with weakly bound impurities as, e.g., H<sup>-</sup>, for which the impurity is clearly located outside the cluster (the ZPE for H<sup>-</sup>He<sub>2</sub> is 90.66 % nostroHeH ). Hence, we expect the Li<sup>+</sup>-He<sub>2</sub> trimer (and its larger clusters) to behave in a different way with respect to the more weakly bound neutral He clusters and thus presume that the classical description of the ionic structures should give us realistic indications on the structure and energetics of such systems (as it was the case in, e.g., the H<sup>-</sup>Ar<sub>n</sub> clusters that we studied earlier nostroArH ).
We therefore begin by looking at the average values of the radial distances and angles (see Table II) which, together with the corresponding standard deviations, describe the ground state of the trimer and clearly confirm its having a rather floppy structure: on the other hand, the Li<sup>+</sup> species still is undoubtedly seen by our calculations to coordinate the two He atoms at a distance determined within 0.3 a<sub>0</sub>. Information on the overall geometrical features of the trimer is further gained from the radial and angular distributions shown in Fig. 3, where the values obtained with the classical optimization are also reported as vertical lines. DMC and DGF results are substantially concident and the small differences are mainly due to the bias contained in our DMC distributions. We notice that the floppiness of the system is particularly evident when looking at the distribution function related to the He-He distance, whose standard deviation is more than twice larger than the one for the Li<sup>+</sup>-He bonds. The classical values obtained in the structure optimization are also closer to the quantum average values obtained for the Li<sup>+</sup>-He distances.
In the classical description we do find that the ground state of the trimer is depicted by an isosceles triangle, with the two shorter sides associated to the two Li<sup>+</sup>-He distances and a longer one corresponding to the He-He distance. On the other hand, the real quantum system cannot be described by one single structure only, and its distribution functions correctly show a delocalized triatom with a dominant contribution from the collinear arrangement. The image of a structure with the Li<sup>+</sup> coordinating the two He atoms at a rather well defined distance (notice the compact distribution function related to Li<sup>+</sup>-He distances) is not completely lost in the quantum description, meaning that the presence of the strong ionic forces from Li<sup>+</sup> is reducing the degree of delocalization which is always present in the pure He aggregates. This change determines the more rigid structure of the system in the sense that now the He atoms are more strongly coordinated directly to the Li<sup>+</sup> impurity (see next section on the larger clusters).
We carried out an additional analysis of the structural features of the trimer by taking advantage of the DGF pseudo-weights 11 , which allow us to pictorially describe the trimer in terms of types of triangular arrangements. In the upper panel of Fig. 4 we thus report all the employed basis set functions, grouped according to the triangular family to which they belong, and in the lower panel of the same figure we report the โweightโ of the dominant families when describing the ground state of the trimer. We can thus easily identify the most important arrangement to be given by the โflatโ isosceles, the collinear (with the the impurity in the middle) and the scalene triangles, while all other possibilities do not contribute in a significant way. Again, we find that we cannot associate the system to one unique structure. while the marked delocalization features are now mainly related to the He-He binding and less to the Li<sup>+</sup>-He ionic forces; hence we see that a conventional structure with the Li<sup>+</sup> coordinating the two He atoms can still be qualitatively identified. In the next section we shall further discuss how the situation evolves with the addition of more He atoms and to which extent the classical results can be still seen to qualitatively correspond to the quantum description of their structures surrounding the ionic dopant.
## V Energetics and structures of $`\mathrm{๐๐ข}^+\mathrm{๐๐}_๐ง`$ (3 $`๐ง`$ 30) clusters from quantum and classical calculations
We start now to discuss the energetics and the geometrical features of the larger $`\mathrm{Li}^+\mathrm{He}_n`$ clusters with $`n3`$. Up to $`n=10`$, infact, it was still possible to also carry out quantum DMC calculations which are not too demanding in terms of CPU-time. Hence, for clusters of such a size we can make a direct comparison between our quantum findings and the classical optimization results we obtained via the combination of the OPTIM procedure walesOPT with a random search for the minimum energy structures (see for details refs. nostroArH ; nostroNeH ). In Table III we report the results for the energetics. The left part of that table shows the minimum potential energies obtained by means of the classical optimizations (column labeled โclassicalโ) in comparison with the corresponding DMC ground state energies (column labelled โquantumโ). The differences between the two sets of values are due to the ZPE effects of the nuclear motions: in the third column we also display the ZPE value for each cluster as a percentage of the well depth. The amount of the ZPE effects increases as the cluster grows, since the incresing addition of He atoms brings the ZPE percentage from about 20% for $`\mathrm{Li}^+\mathrm{He}_2`$ to more than 40% for the last cluster studied here with the quantum DMC method ($`\mathrm{Li}^+\mathrm{He}_{20}`$). On the right part of Table III we report the total energies relative to the loss of one He atom between the pairs of $`\mathrm{Li}^+\mathrm{He}_n`$ and $`\mathrm{Li}^+\mathrm{He}_{n1}`$ clusters, calculated both with the classical and quantum methods. We notice that the evaporation energy is drastically reduced when passing from the cluster with n=6 to that with n=7 and correspondingly the ZPE percentage value increases most markedly (more than 3%). We can then surmise that the structure with n=6 is a particularly stable cluster as we shall further discuss below.
The data presented in Table III are pictorially reported in the two panels of Fig. 5 where the energies are plotted as functions of the number n of He atoms in each cluster. From the lower panel of Fig. 5 we can see the similar behaviour shown by the $`\mathrm{\Delta }E_{ev}`$ values up to $`n=6`$ while for $`7n10`$ both classical and quantum effects make the two curves show a marked drop in values. The single He atom evaporation energies, $`\mathrm{\Delta }E_{ev}=\left[E(\mathrm{Li}^+\mathrm{He}_n)E(\mathrm{Li}^+\mathrm{He}_{n1})\right]`$, (filled-in circles for the classical calculations and open square for those obtained with the DMC method) are plotted up to $`n=30`$ (DMC resul ts up to $`n=20`$): again the two curves show very similar behaviour, both presenting the same abrupt energy drop at $`n=7`$ and $`n=9`$. In contrast with what we found in our analysis of H<sup>-</sup>He<sub>n</sub> clusters nostroHeH , the step-like structure shown by the classical treatment is now also present in the quantum calculations. Given such a correspondence between the classical and quantum description for these smaller clusters it then becomes reasonable to try to explain the sudden energy jumps of Fig. 5 by looking at the lowest energy structures found with the classical minimizations (see Fig. 6). In that figure we also report the corresponding symmetry groups, the total energy (in cm<sup>-1</sup>), and the relevant distances between atoms (in a.u.). We therefore see that the $`\mathrm{Li}^+\mathrm{He}_6`$ has the very symmetrical octahedral geometry (see second panel in the upper part of the figure) where all the He atoms are equivalently coordinated to the central Li<sup>+</sup> at a distance very close to the $`\mathrm{R}_{\mathrm{eq}}`$ (3.58 a<sub>0</sub>) of the $`\mathrm{Li}^+\mathrm{He}`$ PEC. On the other hand, when moving to the $`\mathrm{Li}^+\mathrm{He}_7`$, the repulsive forces acting between the rare gas species do not allow any more for such a symmetrical arrangement around the $`\mathrm{Li}^+`$ ion and the net effect is that of decreasing the energy contributions from the interactions between each He atom and the positive ion, i.e. the $`\mathrm{Li}^+\mathrm{He}`$ distance becomes larger, as one can see in Fig. 6 from the reported $`\mathrm{R}_{\mathrm{LiHe}}`$ values. Similar reasoning can be applied to explain the second energy step between $`\mathrm{n}=8`$ and $`\mathrm{n}=9`$: the evaporation energy gives the mean value of the energy necessary to remove any of the He atoms and the presence of a non-equivalent rare gas atom in the apical position in $`\mathrm{Li}^+\mathrm{He}_9`$ (see second panel in the lower part of Fig. 6) causes a significant decrease of the evaporation energy with respect to its value for the $`\mathrm{Li}^+\mathrm{He}_8`$ cluster. Finally, for $`\mathrm{Li}^+\mathrm{He}_{10}`$ the lowest energy structure we obtained with the classical minimization corresponds to the symmetrical *bicapped square antiprism* geometry. Hence, by using the classical geometries and energy minimization procedures, we found three relatively more stable structures for clusters with $`n`$=6,8 and 10, in correspondence with symmetrically compact structures. However, we cannot associate the closing of a solvating โshellโ to the clusters with $`n`$=6 and $`n`$=8 because they do not constitute as yet a possible core around which the larger clusters grow.
The correspondence between classical and quantum structural pictures is clearly well reproduced if we now compare the quantum distribution functions with the classical results for the relative distances and angles. In Fig. 7 we report the DMC atom-atom distribution functions for the $`\mathrm{Li}^+\mathrm{He}`$ and $`\mathrm{He}\mathrm{He}`$ distances within each $`\mathrm{Li}^+\mathrm{He}_n`$ cluster, normalized to the total number n of possible โconnectionsโ between $`\mathrm{Li}^+`$ and He atoms (solid lines) and to the total number $`N=n(n1)/2`$ possible โconnectionsโ between the n He atoms (dashed lines). In this way we can make a direct comparison between the quantum calculations and the lowest energy geometrical structures obtained with the classical optimizations where the conventional picture of direct bonds existing between localized, point-like, partners can be used. In that figure we also report the classically optimized distances: we have grouped together sets of close values and have given as horizontal bars their statistical standard deviations: each set has a height proportional to the number of distances which have the same value. Finally, the displayed numbers are the values of the integration along r for each broad peak in the quantum distribution functions and represent the number of bonds within atoms which are in the given distance range under each integration. In Fig. 7 we report the results for four selected clusters ($`\mathrm{n}=4,6,7,10`$) which we shall use in our discussion: we also obtained similar results for all the other clusters. When we look at the panel showing the $`\mathrm{Li}^+\mathrm{He}_4`$ clusters in the upper left of Figure 7, we see that the DMC distribution function for the $`\mathrm{Li}^+\mathrm{He}`$ distance peaks at 3.88 a.u. while the classical distances are represented by one โstickโ at 3.59 a.u. whose height is 4, equal to the value of the area under the quantum distribution (solid line); we see also that the DMC distribution function for the $`\mathrm{He}\mathrm{He}`$ distance peaks at 6.13 a.u. while the classical distances are represented by one โstickโ at 5.85 a.u. whose height is 6 (which is the value of $`N=n(n1)/2`$, with $`n=4`$) a number which is indeed equal to the value of the area under the quantum distribution (dashed line). For the larger clusters (see the other panels of the same figure) the number of distinct sets of distances increases, but still the agreement between classical and quantum findings concerning the number of corresponding distances remains very good. Although the mean values are different as a consequence of the very diffuse behaviour of the wavefunctions in such weak interatomic potentials, we can clearly see that the classical results essentially provide the same qualitative structural picture. In order to compare even more in detail the classical and quantum results, we report in Fig 8 the DMC angular distribution functions P($`\theta `$) for the angles centered at the Li<sup>+</sup> ion and at any of the He atoms together with the results from classical optimization: we see again that the agreement between quantum and classical values (vertical lines) indeed remains very close, at least at the qualitative level.
These findings allow us to make further comments on the microsolvation process that occurs when the Li<sup>+</sup> is inserted in a small He cluster. Both the classical and the quantum treatments concur in locating the Li<sup>+</sup> impurity inside the He<sub>n</sub> moiety as one can clearly see from Figure 9 where we report the DMC distribution functions P(r) of the Li<sup>+</sup> (solid lines) and of the He atoms (dashed lines) from the geometrical center of each $`\mathrm{Li}^+\mathrm{He}_n`$ cluster. This quantity is defined as:
$$๐ซ_{gc}=\frac{1}{N}\underset{i=1}{\overset{N}{}}๐ซ_i,$$
(7)
where N runs over the total number of atoms in the cluster. The Li<sup>+</sup> is always closer to the geometrical center with respect to the He atoms, and when the cluster size increases the corresponding distribution functions have reduced overlap, thereby showing that the cluster growth is accompanied by the slow โdroppingโ of the Li<sup>+</sup> towards the geometrical center of the latter
Finally, we carried out classical optimizations for larger clusters $`\mathrm{Li}^+\mathrm{He}_\mathrm{n}`$ (n=11-15, 18,20,22,26,30) in order to better confirm what we have observed in the smaller ones; the lowest energy geometries for a selection of them are reported in Fig. 10. For all the clusters under inspection, the growth occurs around the highly symmetric $`\mathrm{Li}^+\mathrm{He}_{10}`$ core represented by the bicapped square antiprism polyhedron enclosing the Li<sup>+</sup> impurity and drawn with thicker lines in the figure. The additional He atoms are now being placed further away from the impurity without perturbing the structure of the $`\mathrm{Li}^+\mathrm{He}_{10}`$ moiety which therefore seems to constitute the first solvation shell of the atomic ion. We expect that the He atoms outside the shell will be characterized by a greater delocalization and weaker interactions with the central Li<sup>+</sup> that the less shielded inner core of ten atoms. From now on we expect, therefore, that the cluster will grow by adding more solvent atoms in a nearly isotropic fashion driven mainly by He-He interaction and we surmise that the binding energies of each atom will become increasingly similar to those of a pure He cluster. Correspondingly, the evaporation energy (see Fig. 6) now shows a markedly different behaviour with respect to what happens in the smaller clusters with n $``$ 10: the step-like feature disappears, to be substituted by a plateau giving us the average energy needed to remove one of the nearly equivalent external He atoms.
This expected result is also shown by the quantum calculations (see Figure 5, lower panel) and is also borne out by the corresponding quantum distributions of the helium adatoms given by the data of Figure 11, where our DMC calculations for the $`\mathrm{Li}^+(\mathrm{He})_{20}`$ clusters are reported: one clearly see there that the radial distributions associated with the He distances from the Li<sup>+</sup> moiety show a set of more compact values related to the first shell of about 10 adatoms and a further distribution at larger distances (and broader than the first one) associated with the outer atoms that are chiefly bound by dispersion and by strongly screened induction interactions.
## VI Conclusions
In the present work we have analized the solvation process of a Li<sup>+</sup> ion in pure bosonic Helium clusters in order to extend to a larger number of atoms the studies we had already carried out in previous work on the smaller aggregates 8 . Here a combination of classical energy minimization techniques and of โexactโ quantum Monte Carlo methods have been employed in order to describe the structure and the energetics of the Li<sup>+</sup>(He)<sub>n</sub> clusters with $`n`$ up to 30 (20 for DMC calculations), employing always a description based on the sum-of-potential approximation.
The basic approximation which this study relies on is that of calculating the full cluster interaction as a sum of pairwise potentials. Although the error introduced by this assumption is in general found to be (in absolute terms) larger in ionic systems than in neutral ones, its relative weight remains rather small 8 . Thus, we believe that the resulting stuctures and energies would not be substantially altered by the correct use of a full Many-Body potential (see for example, the discussion in Ref. 29 for the similar situations of H<sup>-</sup> and Cl<sup>-</sup> in Rg clusters and our calculations in Ref. 8 ).
We have therefore shown that:
1. the Li<sup>+</sup> is fully solvated inside larger <sup>4</sup>He clusters as already indicated by our earlier work on very small aggregates 8 ;
2. the ionic Li<sup>+</sup> core does not form preferential โmolecular coresโ with surrounding He atoms but rather that the helium adatoms remain equivalently bound within each solvation shell to the central, solvated lithium ion that persists in carrying the positive charge for more than 98% 8 .
3. the ZPE corrections play a role in such systems, albeit strongly reduced with respect to the one found in neutral aggregates: this means that, at least in the initial solvation shell, the quantum adatoms are less delocalized and that the lithium-helium direct โbondsโ are nearly rigid, classical structures;
4. the classical optimization procedures can provide structural details which are reasonably close to those given by the quantum DMC calculations and can yield for ionic moieties the same structural picture as that given by the quantum treatment.
Furthermore, our comparison of single particle evaporation energies given by classical and quantum results suggests in both cases the formation of an initial shell of about ten He atoms which are more strongly bound to the central ion. On the other hand, beyond that initial shell the cluster growth appears to be chiefly driven by He-He interactions, albeit at energies which are initially still kept larger than those of the neutral systems by the additional presence of the induction field due to the central ionic core. This bahevior may be compared to that of Na<sup>+</sup> and K<sup>+</sup> doped Helium clusters 5a where the first solvation shells were found to be of 9 and 12 atoms respectively.
###### Acknowledgements.
The financial support of the FIRB project, of the University of Rome โLa Sapienzaโ Scientific Committee and of the European Union โCold Moleculesโ Collaborative Research Project no. HPRN-CT-2002-00290 is gratefully acknowledged. One of us (M.Y.) thanks I. T. U. Research Fund for the financial support and I. T. U. High Performance Computing Center for the computer time provided. We are grateful to Prof. M. Morosi and Dr. D. Bressanini for their help in improving the choice of trial functions. We also acknowledge the support of the INTAS grant 03-51-6170. |
warning/0506/hep-th0506126.html | ar5iv | text | # Untitled Document
hep-th/0506126
Two-Loop Amplitudes of Gluons and Octa-Cuts
in $`๐ฉ=4`$ Super Yang-Mills
Evgeny I. Buchbinder<sup>a</sup> and Freddy Cachazo<sup>b</sup>
<sup>a</sup> School of Natural Sciences, Institute for Advanced Study, Princeton NJ 08540 USA
<sup>b</sup> Perimeter Institute for Theoretical Physics, Waterloo, Ontario N2J 2W9, Canada
After reduction techniques, two-loop amplitudes in $`๐ฉ=4`$ super Yang-Mills theory can be written in a basis of integrals containing scalar double-box integrals with rational coefficients, though the complete basis is unknown. Generically, at two loops, the leading singular behavior of a scalar double box integral with seven propagators is captured by a hepta-cut. However, it turns out that a certain class of such integrals has an additional propagator-like singularity. One can then formally cut the new propagator to obtain an octa-cut which localizes the cut integral just as a quadruple cut does at one-loop. This immediately gives the coefficient of the scalar double box integral as a product of six tree-level amplitudes. We compute, as examples, several coefficients of the five- and six-gluon non-MHV two-loop amplitudes. We also discuss possible generalizations to higher loops.
June 2005
1. Introduction
Recently there has been renewed interest in the perturbation expansion of $`๐ฉ=4`$ super Yang-Mills. This was motivated by the discovery of a twistor string theory that captures the perturbation theory of the maximally supersymmetric Yang-Mills theory (pMSYM). Twistor string theory has opened new avenues and has inspired new ideas for the computation of tree level amplitudes of gluons \[2,,3,,4,,5,,6,,7,,8,,9,,10,,11,,12,,13,,14,,15,,16,,17,,18,,19\] and one-loop amplitudes of gluons in QCD \[20,,21,,22\], $`๐ฉ=1`$ \[23,,24,,25,,26,,27,,28,,29\] and $`๐ฉ=4`$ \[30,,31,,32,,33,,34,,35\] super Yang-Mills. Before twistor string theory was introduced, the study of pMSYM at one-loop was mainly motivated by two facts: one is the decomposition of a QCD amplitude, $`A^{QCD}`$, with only a gluon running in the loop in terms of supersymmetric amplitudes and an amplitude with only a scalar running in the loop, $`A^{\mathrm{scalar}}`$, (see for a review),
$$A^{\mathrm{QCD}}=A^{๐ฉ=4}4A_{\mathrm{chiral}}^{๐ฉ=1}+A^{\mathrm{scalar}}$$
where $`A^{๐ฉ=4}`$ has the full $`๐ฉ=4`$ multiplet in the loop and $`A_{\mathrm{chiral}}^{๐ฉ=1}`$ only an $`๐ฉ=1`$ chiral multiplet. The other motivation is a surprising proposal of Anastasiou, Bern, Dixon, and Kosower (ABDK) that two- (and, perhaps, higher-) loop amplitudes in pMSYM can be completely determined in terms of one-loop amplitudes . This idea was inferred from studying collinear and IR singular behavior of the higher loop amplitudes. The conjecture is given in terms of normalized $`2`$-loop amplitudes $`M_n^{(2)}=A_n^{(2)}/A_n^{\mathrm{tree}}`$ and in dimensional regularization, as follows
$$M_n^{(2)}(ฯต)=\frac{1}{2}\left(M_n^{(1)}(ฯต)\right)^2+f(ฯต)M_n^{(1)}(2ฯต)\frac{5}{4}\zeta _4+๐ช(ฯต).$$
This relation was explicitly verified for four-gluon amplitudes in (see also section 7 of ). Also based on collinear limits , the schematic form of a relation analogous to (1.1) was proposed for higher loops . Very recently, an explicit formula, analogous to (1.1), for the three-loop four-gluon amplitude was obtained and successfully verified in . It is the aim of this paper to make some modest steps towards the calculation of higher loop amplitudes in pMSYM. The main motivation is to prepare the ground for future tests of the ABDK proposal. A proof of (1.1) would lead to the solution of pMSYM at two loops as a general solution to the one-loop problem can be obtained in terms of tree-level amplitudes by using quadruple cuts . This is possible thanks to the cut constructibility of one-loop amplitudes in pMSYM proven in and the decomposition in terms of scalar box integrals, with rational functions as coefficients, also given in . See also \[30,,31,,32,,33,,35\] for other techniques in pMSYM at one loop.
At two loops, a similar decomposition in terms of some given set of integrals is expected by using Passarino-Veltman or similar reduction procedures . Unfortunately, the complete basis of two-loop integrals is currently unknown. However, scalar double box integrals are a natural ingredient of such a basis<sup>1</sup> In fact, the four-gluon amplitude is given only in terms of scalar double boxes .. In this paper, we concentrate on the calculation of the coefficient of certain classes of planar scalar double box integrals. These are the integrals that arise in scalar field theory with a massless scalar running along internal lines and with the double-box structure depicted in fig. 1.
Fig. 1: The two possible different structures of planar scalar double box integrals. $`(a)`$ Double boxes. $`(b)`$ Split double boxes. Note that the momenta of the external lines is given by the sum of the momenta of external gluons.
The momenta of the external legs in fig. 1 are given by sums of momenta of external gluons.
We propose a method for computing the coefficient of any scalar double box integral given in fig. 1a when at least one of the two boxes has two adjacent massless three-particle vertices. We also give the form of the coefficient of any double box given in fig. 1b. In order to distinguish between the double boxes in fig. 1a and in fig. 1b we refer to the former simply as โdouble boxesโ and the latter as โsplit double boxesโ.
Our original motivation was the successful use of quadruple cuts in the calculation of one-loop $`๐ฉ=4`$ amplitudes . The basic idea is that at one-loop only scalar boxes contribute . A quadruple cut singles out the contribution of a given scalar box and localizes the integration over the loop momentum. The combination of these two facts allows one to calculate any coefficient in terms of the product of four tree-level amplitudes . Up to a numerical factor, every one-loop box coefficient is given by
$$B=A_{(1)}^{\mathrm{tree}}A_{(2)}^{\mathrm{tree}}A_{(3)}^{\mathrm{tree}}A_{(4)}^{\mathrm{tree}},$$
where the sum is over the solutions to the delta function equations and over all particles that can propagate in the loop. A straightforward application of this idea can be made for split double boxes (see fig. 1b). Again the idea is to cut all eight propagators, i.e., an octa-cut which localizes the two loop integrations and gives the coefficient as the product of seven tree-level amplitudes (up to a numerical factor)
$$B=\underset{i=1}{\overset{7}{}}A_{(i)}^{\mathrm{tree}}.$$
Naively, one might expect that the coefficient of double boxes in fig. 1a cannot be computed in a similar manner. The reason is that there are only seven propagators and a hepta-cut does not localize the integrals over the loop momenta.
A way to avoid the remaining integration arises in an unexpected manner. In studying singularities of Feynman integrals, one computes the discontinuity of an integral across a singularity by cutting propagators. When one cuts all propagators in a Feynman diagram one is computing the discontinuity across the leading singularity of the integral. However, at two (and higher) loops one finds a surprise when some of the external legs are massless. At two loops, if any of the two boxes in fig. 1a has at least two adjacent three-particle vertices (condition that is satisfied trivially for less than seven external gluons), then the integral has an extra propagator-like singularity beyond the naive leading singularity. The discontinuity across the new leading singularity is actually computed by an octa-cut<sup>2</sup> For more general double boxes, there is also an extra singularity, these are known as second-type singularities . They cannot easily be used to produce an octa-cut but they might give a generalization of it.. This octa-cut precisely localizes the loop integrations and allows a straightforward computation of the coefficient as the product of six tree-level amplitudes. Up to a numerical factor, it is given by
$$B=\underset{i=1}{\overset{6}{}}A_{(i)}^{\mathrm{tree}}.$$
The only two-loop amplitude in pMSYM known in the literature is the four-gluon amplitude . One reason is that very few double scalar box integrals are known explicitly . In particular, to our knowledge, not all double box integrals needed for a five-gluon amplitude are known. Nevertheless, we present the computation of several five-gluon and six-gluon non-MHV scalar double box integrals as illustrations of our technique.
This paper is organized as follows. In section 2, we review pMSYM at tree-, one-, two- and three-loop levels as well as the ABDK conjecture. In section 3, we show that the four-gluon amplitude of pMSYM can be found by using hepta-cuts. Even though the number of cut propagators is less than the number of integration variables, the integrand turns out not to depend on the loop momenta and can be pulled out of the integral. In section 4, we demonstrate that a certain class of double-box configurations admit an extra propagator-type singularity. Cutting this singularity allows us to write a universal formula for many double-box coefficients. In section 5, we illustrate our technique via various examples including non-MHV amplitudes. In section 6, we discuss applications to three- and higher-loop amplitudes. In particular, we show that by studying singularities, it is possible to realize that the basis of integrals has to contain integrals with some non-trivial factors in the numerator. This is in agreement with results of .
Throughout the paper, we use the following notation and conventions along with those of and the spinor helicity-formalism \[47,,48,,49\]. A external gluon labeled by $`i`$ carries momentum $`K_i`$. Since $`K_i^2=0`$, it can be written as a bispinor $`(K_i)_{a\dot{a}}=\lambda _{ia}\stackrel{~}{\lambda }_{i\dot{a}}`$. Inner product of null vectors $`p_{a\dot{a}}=\lambda _a\stackrel{~}{\lambda }_{\dot{a}}`$ and $`q_{a\dot{a}}=\lambda _a^{}\stackrel{~}{\lambda }_{\dot{a}}^{}`$ can be written as $`2pq=\lambda ,\lambda ^{}[\stackrel{~}{\lambda },\stackrel{~}{\lambda }^{}]`$, where $`\lambda ,\lambda ^{}=ฯต_{ab}\lambda ^a\lambda ^b`$ and $`[\stackrel{~}{\lambda },\stackrel{~}{\lambda }^{}]=ฯต_{\dot{a}\dot{b}}\stackrel{~}{\lambda }^{\dot{a}}\stackrel{~}{\lambda }^{\dot{b}}`$. Other useful definitions are:
$$\begin{array}{cc}\hfill K_{i,\mathrm{},j}& K_i+K_{i+1}+\mathrm{}+K_j\hfill \\ \hfill K_i^{[r]}& K_i+K_{i+1}+\mathrm{}+K_{i+r1}\hfill \\ \hfill i|\underset{r}{}K_r|j]& \underset{r}{}ir[rj]\hfill \\ \hfill i|(\underset{r}{}K_r)(\underset{s}{}K_s)|j& \underset{r}{}\underset{s}{}ir[rs]sj\hfill \\ \hfill [i|(\underset{r}{}K_r)(\underset{s}{}K_s)|j]& \underset{r}{}\underset{s}{}[ir]rs[sj]\hfill \\ \hfill i|(\underset{r}{}K_r)(\underset{s}{}K_s)(\underset{t}{}K_t)|j]& \underset{r}{}\underset{s}{}\underset{t}{}ir[rs]st[tj]\hfill \end{array}$$
where addition of indices is always done modulo $`n`$.
2. Review of $`๐ฉ=4`$ Amplitudes
In this paper we consider amplitudes of gluons in $`๐ฉ=4`$ super-Yang-Mills. Each gluon carries the following information: momentum $`p_{a\dot{a}}`$, polarization vector $`ฯต_{a\dot{a}}`$ and color index $`a`$. The color structure can be striped out by a color decomposition \[50,,51,,52,,53\]. Here we only consider the leading color or planar part of the amplitudes. The information in momentum and polarization vectors can be encoded in terms of spinors $`\lambda `$, $`\stackrel{~}{\lambda }`$ and the helicity of the gluon $`h`$.
2.1. Tree-Level $`๐ฉ=4`$ Amplitudes
At tree-level, the leading color approximation is exact. An amplitude is given by
$$A_{(\{p_i,ฯต_i,a_i\})}=g_{\mathrm{YM}}^{n2}\underset{\sigma S_n/Z_n}{}\mathrm{Tr}(T^{a_{\sigma (1)}}\mathrm{}T^{a_{\sigma (n)}})A_{(\{\lambda _{\sigma (1)},\stackrel{~}{\lambda }_{\sigma (1)},h_{\sigma (1)}\},\mathrm{},\{\lambda _{\sigma (n)},\stackrel{~}{\lambda }_{\sigma (n)},h_{\sigma (n)}\})}.$$
Here we are suppressing a delta function that imposes momentum conservation.
It is convenient to denote the set of data $`\{\lambda _i,\stackrel{~}{\lambda }_i,h_i\}`$ by $`i^{h_i}`$, where $`h_i=\pm `$ is the helicity of the $`i^{th}`$ gluon. The amplitudes on the right hand side of (2.1) are known as leading color partial amplitudes and are computed from color-ordered Feynman rules. One can study a given order $`A(1^{h_1},\mathrm{},n^{h_n})`$ and the rest can be obtained by application of permutations, $`\sigma `$.
The partial amplitude $`A(1^{h_1},\mathrm{},n^{h_n})`$ can be computed using a variety of methods (see for a nice review on many of the techniques developed in the 80โs and 90โs). More recently, two new techniques became available, namely, MHV diagrams and the BCFW recursion relations \[14,,15\]. The latter is a set of quadratic recursion relations for on-shell physical partial amplitudes of gluons. For a recent review see .
2.2. One-Loop $`๐ฉ=4`$ Amplitudes
Amplitudes of gluons at one-loop admit a color decomposition \[50,,51,,52,,53,,55\] with single and double trace contributions. As mentioned in the introduction we will only concentrate on the leading color partial amplitudes<sup>3</sup> It is interesting to note that since for $`๐ฉ=4`$ SYM all particles in the loop are in the adjoint representation, all sub-leading color amplitudes are given as linear combinations of the planar ones with permutations of the gluon labels (See section 7 of for a proof.).
One-loop amplitudes of gluons in supersymmetric theories are four-dimensional cut-constructible \[42,,56\]. This means that they can be completely determined by their finite branch cuts and discontinuities. $`๐ฉ=4`$ amplitudes are even more special. Reduction techniques can be used to express these amplitudes in terms of scalar box integrals . These are one-loop box Feynman integrals in a scalar field theory where a massless scalar runs in the loop,
$$=d^4\mathrm{}\frac{1}{(\mathrm{}^2+iฯต)((\mathrm{}k_1)^2+iฯต)((\mathrm{}k_1k_2)^2+iฯต)((\mathrm{}+k_4)^2+iฯต)}$$
where $`k_1,k_2,k_3,k_4`$ are the external momenta at each vertex. They are not independent since by momentum conservation $`k_3=(k_4+k_1+k_2)`$. Note that the integral (2.1) is singular when at least one $`k_i`$ is a null vector. Therefore, we should specify a regularization procedure, like dimensional regularization. However, we will be considering cuts that are finite and do not depend on the regularization procedure. Since $`A(1,\mathrm{},n)`$ is color-ordered, each $`k`$ can only be the sum of consecutive momenta of external gluons. Moreover, since we only consider the planar contributions, we can define a given contribution by specifying $`i,j,k,l`$ such that $`k_1=K_i+\mathrm{}+K_{j1}`$, $`k_2=K_j+\mathrm{}+K_{k1}`$ and $`k_3=K_k+\mathrm{}+K_{l1}`$. The reduction procedure then gives for the amplitude an expansion of the form
$$A(1,\mathrm{},n)=\underset{1<i<j<k<m<n}{}B_{ijkl}_{(K_i+\mathrm{}+K_{j1},K_j+\mathrm{}+K_{k1},K_k+\mathrm{}+K_{l1})},$$
where the coefficients $`B_{ijkm}`$ are rational functions of the spinor products. Since all scalar box integrals are known explicitly, the problem of computing $`A(1,\mathrm{},n)`$ is reduced to that of computing the coefficients $`B_{ijkl}`$.
A general formula for the coefficients $`B_{ijkl}`$ was found in in terms of products of tree level amplitudes. Let us review the derivation of the formula because the idea is useful in the analysis at higher loops. If we think of the scalar box integrals as an independent basis<sup>4</sup> The notion of independence is the equivalent of cut constructibility of the amplitude. of some vector space we can interpret $`A(1,\mathrm{},n)`$ as a general vector. All we need to do is to find a way to project $`A(1,\mathrm{},n)`$ along the space spanned by a given scalar box integral $``$. From the definition of $``$ in (2.1) we see that each integral is uniquely determined once its four propagators are given. It is natural to think that the way to determine the coefficient $`B`$ is by looking at the region of integration where all four propagators become singular. In fact, the integral obtained by cutting, i.e., by dropping the principal part of all four propagators computes the discontinuity of the given scalar box integral across a singularity which is unique to it.
The set of four equations that gives $`\mathrm{}`$ is the following
$$\mathrm{}^2=0,(\mathrm{}k_1)^2=0,(\mathrm{}k_1k_2)^2=0,(\mathrm{}+k_4)^2=0.$$
A little exercise shows that these equations do not have a solution if $`\mathrm{}`$ is a real vector in Minkowski space for general external momenta. The way out of this problem is to complexify all momenta and make a Wick rotation to $`(++)`$ signature. In the new signature the delta functions are still well defined and there are always solutions to (2.1).
Fig. 2: A quadruple cut diagram. Momenta in the cut propagators flows clockwise and external momenta are taken outgoing. The tree-level amplitude $`A_1^{\mathrm{tree}}`$, for example, has external momenta $`\{i+1,\mathrm{},j,\mathrm{}_2,\mathrm{}_1\}`$.
One can also look at the same regime on the left hand side of (2.1) by considering only Feynman diagrams that posses the four propagators entering in (2.1). Summing over them one finds the following equation
$$๐\mu \underset{J}{}n_JA_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}=B_{ijkm}๐\mu $$
where sum over $`J`$ represents a sum over all possible particles in the $`๐ฉ=4`$ multiplet. The measure $`d\mu `$ is the same one both sides of the integrals,
$$d\mu =d^4\mathrm{}\delta ^{(+)}(\mathrm{}^2)\delta ^{(+)}((\mathrm{}k_1)^2)\delta ^{(+)}((\mathrm{}k_1k_2)^2)\delta ^{(+)}((\mathrm{}+k_4)^2),$$
and the tree-level amplitudes are defined as follows (see fig. 2)
$$\begin{array}{cc}\hfill A_{(1)}^{tree}=& A(\mathrm{}_1,i+1,i+2,\mathrm{},j1,j,\mathrm{}_2),A_{(2)}^{tree}=A(\mathrm{}_2,j+1,j+2,\mathrm{},k1,k,\mathrm{}_3),\hfill \\ \hfill A_{(3)}^{tree}=& A(\mathrm{}_3,k+1,k+2,\mathrm{},m1,m,\mathrm{}_4),A_{(4)}^{tree}=A(\mathrm{}_4,m+1,m+2,\mathrm{},i1,i,\mathrm{}_1).\hfill \end{array}$$
where
$$\begin{array}{cc}& \mathrm{}_1=\mathrm{},\mathrm{}_2=\mathrm{}k_1,\mathrm{}_3=\mathrm{}k_1k_2,\mathrm{}_4=\mathrm{}+k_4,k_1=K_{i+1}+\mathrm{}+K_j,\hfill \\ & k_2=K_{j+1}+\mathrm{}+K_k,k_3=K_{k+1}+\mathrm{}+K_m,k_4=K_{m+1}+\mathrm{}+K_i.\hfill \end{array}$$
The integral $`๐\mu `$ is just given by a Jacobian $`1/\sqrt{\mathrm{\Delta }}`$. This Jacobian cancels on both sides since the integral is localized by the delta functions and the coefficient is given by
$$B_{ijkl}=\frac{1}{|๐ฎ|}\underset{๐ฎ,J}{}n_JA_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}.$$
Here $`๐ฎ`$ is the set of solutions to the conditions imposed by the delta functions, and $`|๐ฎ|`$ is the number of solutions. The sum also involves a sum over all possible particles that can propagate in the loop. For further details and many examples we refer to . Even though the Jacobian did not play an important role for the quadruple cut technique at one-loop, it is crucial for the two-loop analysis we carry out in section 4.1. For this reason let us write it down for future reference
$$\mathrm{\Delta }=s^2t^22st(k_1^2k_3^2+k_2^2k_4^2)+(k_1^2k_3^2k_2^2k_4^2)^2$$
with $`s=(k_1+k_2)^2`$ and $`t=(k_2+k_3)^2`$.
2.3. Two-Loop $`๐ฉ=4`$ Amplitudes
At two loops, only the four-gluon amplitude has been computed . The $`๐ฉ=4`$ calculation was the first full two-loop amplitude of gluons ever computed. The answer is given as a linear combination of double box scalar integrals with coefficients that are rational function of the spinor variables. A double box scalar integral is the analog of the one-loop scalar box integral introduced above, more explicitly,
$$\begin{array}{cc}\hfill (k_1,\mathrm{},k_6)=& \frac{d^4p}{(2\pi )^4}\frac{1}{(p^2+iฯต)((pk_1)^2+iฯต)((pk_1k_2)^2+iฯต)}\times \hfill \\ & \frac{d^4q}{(2\pi )^4}\frac{1}{((p+q+k_6)^2+iฯต)(q^2+iฯต)((qk_5)^2+iฯต)((qk_4k_5)^2+iฯต)}.\hfill \end{array}$$
This integral is UV finite but it might have IR divergences when some $`k`$โs are null vectors. Again, as in the one-loop case, one has to choose a regularization procedure but we do not do so because we only discuss finite cuts. The planar contribution to the four-gluon amplitude is
$$A_4^{2loop}(K_1,K_2,K_3,K_4)=A_4^{tree}st\left(s(K_1,K_2,0,K_3,K_4,0)+t(K_4,K_1,0,K_2,K_3,0)\right)$$
where $`s=(K_1+K_2)^2`$ and $`t=(K_2+K_3)^2`$. This was computed by using the unitarity-based method \[42,,56,,57,,58,,59\]. It is very important to mention that the double box scalar integral (2.1) is not known in general but explicitly formulas exists in dimensional regularization when $`k_3=k_6=0`$ and at least three of the other $`k_i`$โs are null vectors .
2.4. ABDK Conjecture
As mentioned in the introduction, one of the motivations of this work is to prepare the ground for a more extensive test of the ABDK conjecture. This conjecture asserts that the planar limit of $`L`$-loop amplitudes in $`๐ฉ=4`$ SYM is determined iteratively, i.e., as a function of $`l`$-loop amplitudes with $`l<L`$.
Let us make this more precise. Here we follow and where the original proposal was made. Consider the function
$$M_n^{(L)}(1,2,\mathrm{},n)=\frac{A^{Lloop}(1,2,\mathrm{},n)}{A^{tree}(1,2,\mathrm{},n)}$$
then the ABDK conjecture states that
$$M_n^{(L)}=P_L(M_n^{(1)},\mathrm{},M_n^{(L1)})$$
where $`P_L(x_1,\mathrm{},x_{L1})`$ is a certain polynomial of degree $`L`$ and independent of the helicity configuration. The explicit form of (2.1) at two loops was given in in terms of the function $`f(ฯต)=(\psi (1ฯต)\psi (1))/ฯต`$, where the digamma function is defined by $`\psi (x)=\mathrm{\Gamma }^{}(x)/\mathrm{\Gamma }(x)`$, as follows
$$M_n^{(2)}(ฯต)=\frac{1}{2}\left(M_n^{(1)}(ฯต)\right)^2+f(ฯต)M_n^{(1)}(2ฯต)\frac{5}{4}\zeta _4.$$
This conjecture was explicitly checked for four-gluon amplitudes. Very recently, the form of the polynomial in (2.1) was obtained for the three-loop four-gluon amplitude in . One of the impressive predictions of the conjecture is a relation between the finite remainders which are defined at $`ฯต=0`$. At two loops, one introduces the universal singular function $`C_n(ฯต)^{(2)}`$ \[60,,37\] which contains the infrared singularities and does not depend on the helicity configuration since it is normalized by the tree-level amplitude. Defining the finite remainder as
$$F_n^{(2)}(ฯต)=M_n^{(2)}(ฯต)C_n^{(2)}(ฯต),$$
one can write a finite (as $`ฯต0`$ ) analog of (2.1) as follows
$$F_n^{(2)}(0)=\frac{1}{2}\left(F_n^{(1)}(0)\right)^2\zeta _2F_n^{(1)}(0)\frac{1}{4}\left(\frac{11n}{8}+5\right)\zeta _4.$$
Recall that at one-loop $`F_n^{(1)}(0)`$ can have at most dilogarithms, while $`F_n^{(2)}(0)`$ can have higher polylogarithms. This means that very non-trivial cancelations must happen. These cancelations were found to occur for $`n=4`$ between terms coming from the two integrals in (2.1) and involved many polylogarithmic identities . In the recent paper , an impressive formula for the all loop finite remainder of MHV amplitudes was also presented. The formula is given in a kind of generating function structure
$$1+\underset{L=1}{\overset{\mathrm{}}{}}a^LF_n^{(L)}(0)=\mathrm{exp}\left[\frac{1}{4}\gamma _KF_n^{(1)}(0)+C\right]$$
where $`a`$ is basically the โt Hooft coupling, $`\gamma _K`$ is the universal soft anomalous dimension and $`C`$ is a function that admits a power series representation in $`a`$. $`\gamma _K`$ and $`C`$ are known up to the order needed to obtain the three-loop term<sup>5</sup> It is important to mention that there is no canonical definition of the finite remainder $`F_n`$. In fact, the definition of finite remainder used in (2.1) differs from that used in (2.1). For more details see . We thank Z. Bern and L. Dixon for useful discussions on this point..
3. Four-Gluon Two-Loop Amplitudes and Hepta-Cuts
In this section, we consider hepta-cuts of the two-loop four-gluon leading partial amplitude. This section can be viewed as a warm-up section where we introduce relevant notation and do some sample calculations which will be used in the rest of the paper. It is enough to consider $`A_4^{2loop}(1^{},2^{},3^+,4^+)`$ as all other $`A_4^{2loop}`$ with different helicity assignments can be obtained from this one by Ward identities. The leading partial amplitude was first computed in . As reviewed in section 2.3, the amplitude can be presented as a linear combination of two scalar double-box integrals (see fig. 1a)
$$=\frac{d^4p}{(2\pi )^4}\frac{d^4q}{(2\pi )^4}\frac{1}{p^2(pK_1)^2(pK_1K_2)^2(p+q)^2q^2(qK_4)^2(qK_3K_4)^2},$$
where $`K_i`$ are the four external gluon momenta, with rational coefficients. All external momenta are assumed to be outgoing. The integral (3.1) has seven propagators, hence it is natural to consider hepta-cuts. It turns out that the coefficients can easily be found from hepta-cuts when the loop momenta are analytically continued to signature $`(++)`$ or complexified. In the present case, there are two independent coefficients as well as two independent hepta-cuts. We refer to them as the $`s`$-channel cut and the $`t`$-channel cut. The corresponding coefficients will be denoted as $`c_s`$ and $`c_t`$. Let us start with the cut in the $`s`$-channel. In this case, there are six different helicity configurations. For all of them, only gluons can propagate in both loops. A sample helicity configuration is shown in fig. 3. In this paper, for simplicity, we depict tree level amplitudes as points. Since all propagator are cut, there is no need to indicate a cut by a dash line and we choose to omit them <sup>6</sup> Note that conventions in fig. 3 are different from those used in fig. 2 where all tree level amplitudes are denoted by blobs and cuts are indicated by dashed lines..
Fig. 3: A sample hepta-cut in the $`s`$-channel. Tree level amplitudes are depicted as points. All propagators are cut and therefore we omit the dashed lines used in fig. 2.
The rational coefficient $`c_s`$ is then given by
$$c_s=\frac{i^7\underset{I=1}{\overset{6}{}}๐\mu (A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree})_I}{๐\mu },$$
where by $`A_{(i)}^{tree}`$ we denote tree amplitudes at each of the six vertices, the integration measure $`d\mu `$ is given by
$$\begin{array}{cc}\hfill d\mu =& \frac{d^4p}{(2\pi )^4}\frac{d^4q}{(2\pi )^4}\delta (p^2)\delta ((pK_1)^2)\delta ((pK_1K_2)^2)\delta ((p+q)^2)\hfill \\ & \delta (q^2)\delta ((qK_4)^2)\delta ((qK_3K_4)^2)\hfill \end{array}$$
and in the numerator we sum over the product of the six tree level amplitudes corresponding to a given helicity configuration. The factor $`i^7`$ comes from the seven propagators. It seems that since the number of delta functions is less than the number of integration variables, the integral does not localize and one integration has to be performed. However, it turns out that the integrand can be simplified in such a way that the dependence on the loop momenta drops out and we are left with the integral of the measure which cancels out according to eq. (3.1).
In the discussion of one-loop amplitudes in section 2.2, we mentioned that the momentum $`\mathrm{}`$ has to be complexified in order to find solutions to the four equations from the cut propagators. Making $`\mathrm{}`$ complex also has as a byproduct the fact that three-particle vertices on-shell do not have to vanish. Tree-level three-gluon amplitudes with helicities $`(+)`$ or $`(++)`$ are given respectively by \[61,,62\]
$$A_3^{tree}(p^{},q^{},r^+)=i\frac{pq^3}{qrrp},A_3^{tree}(p^+,q^+,r^{})=i\frac{[pq]^3}{[qr][rp]}.$$
In Minkowski space, $`\lambda _p`$ and $`\stackrel{~}{\lambda }_p`$ are related to each other as $`\stackrel{~}{\lambda }_p=\pm \overline{\lambda }_p`$. This means that if $`pq=0`$, which follows from momentum conservation at the vertex, then both $`\lambda _p\lambda _q=0`$ and $`[\stackrel{~}{\lambda }_p\stackrel{~}{\lambda }_q]=0`$. This implies that both amplitudes in (3.1) vanish. If we complexify the momenta, then the equation $`pq=0`$ has two independent solutions. We have that either $`\lambda _p\lambda _q=0`$ or $`[\stackrel{~}{\lambda }_p\stackrel{~}{\lambda }_q]=0`$. That is either $`\lambda _p`$ and $`\lambda _q`$ are proportional or $`\stackrel{~}{\lambda }_p`$ and $`\stackrel{~}{\lambda }_q`$ are proportional. Also note that momentum conservation implies that $`pq=pr=qr=0`$. This means that either three $`\lambda `$โs are proportional or all three $`\stackrel{~}{\lambda }`$โs are proportional. Therefore, for every $`(++)`$ tree level amplitude we choose all $`\lambda `$โs to be proportional. Similarly, for every $`(+)`$ tree level amplitude we choose all three $`\stackrel{~}{\lambda }`$โs to be proportional.
Explicit calculations, considered for one of the helicity configurations in some detail below, show that every helicity configuration gives the same contribution equal to
$$A_4^{tree}s^2t๐\mu ,$$
where $`A_4^{tree}`$ is the tree-level four-gluon amplitude
$$A_4^{tree}(1^{},2^{},3^+,4^+)=i\frac{12^3}{233441},$$
and
$$s=(K_1+K_2)^2,t=(K_2+K_3)^2.$$
Note that the integral in (3.1) cancels against the denominator in (3.1). The coefficient $`6`$ in the numerator will also cancel. The reason is the following. In the denominator in (3.1), we have to sum over all different solutions to the delta-function conditions. It is easy to realize that in this particular case the number of different solutions equals the number of helicity configurations. Thus, each term in the numerator in (3.1) picks one of the six solutions whereas in the denominator we sum over all the six solutions. As a result, we obtain
$$c_s=A_4^{tree}s^2t.$$
This coincides with the corresponding coefficient found in .
Let us consider the helicity configuration shown in fig. 3 in some detail. The analysis of the remaining five configurations is completely analogous. Consider the integrand as the product of six tree-level amplitudes
$$\begin{array}{cc}& i(A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree})_1=i\left(\frac{[p_1,p]^3}{[p,1][1,p_1]}\right)\left(\frac{p_1,2^3}{2,p_2p_2,p_1}\right)\hfill \\ & \left(\frac{[q_2,l]^3}{[l,p_2][p_2,q_2]}\right)\left(\frac{p,l^3}{l,qq,p}\right)\left(\frac{[q_1,4]^3}{[4,q][q,q_1]}\right)\left(\frac{q_1,q_2^3}{q_2,33,q1}\right).\hfill \end{array}$$
Next, simplify this expression by using momentum conservation. For example, the product of $`[p_1,p]`$ and $`p_1,2`$ can be simplified as follows
$$[p_1,p]p_1,2=2|p_1|p]=2|K_1|p]=12[1p].$$
Then the product of the first four factors in (3.1) becomes
$$12^2[q,q_2]^2.$$
After using momentum conservation along the lines of eq. (3.1), one finds
$$i(A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree})_1=is^2\frac{12^2[34]}{34}=A_4^{tree}s^2t.$$
Note that this expression does not depend on the loop momenta and, thus, can be pulled out of the integration.
Now we consider the hepta-cut in the $`t`$-channel. Here we have ten helicity configurations. Note that in this case the number of helicity configurations does not equal the number of solutions of the delta-function equations. By a solution we mean a choice whether all $`\lambda `$โs are proportional or $`\stackrel{~}{\lambda }`$โs are proportional at each of six three gluon vertices. However, among the ten configurations, there are different configurations for which the choices of whether $`\lambda `$โs or $`\stackrel{~}{\lambda }`$โs are proportional are exactly the same. A solution then means summing up over such configurations. In this case, there are two helicity configurations corresponding to actual solutions and the remaining eight ones break up in pairs. The sum of the two helicity configurations in each pair corresponds to a solution of the delta-function conditions. Overall, we have six improved helicity configurations, each corresponding to an independent solution to the delta-function conditions. All paired up configurations involve fermions and scalars running in one of the loops and summation over the two configurations in a given pair provides a significant simplification. The coefficient $`c_t`$ is given by
$$c_t=\frac{i\underset{I=1}{\overset{10}{}}๐\mu (A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree})_I}{๐\mu },$$
where the sum is over all ten configurations, or over the six improved configurations, each corresponding to an actual solution of the delta-function equations. As before, all six improved configurations give the same contribution
$$A_4^{tree}st^2.$$
Since each improved configuration corresponds to a solution to the delta-function equations, the factor $`6`$ cancels out. As a result we obtain
$$c_t=A_4^{tree}st^2.$$
This coincides with the corresponding coefficient from . As an example, let us consider the two helicity configurations shown in fig. 4.
Fig. 4: Examples of hepta-cuts in the $`t`$-channel that correspond to the same solution of the delta function constraints.
Note that for both double boxes the choices whether all $`\lambda `$โs or all $`\stackrel{~}{\lambda }`$โs are proportional are exactly the same at every vertex. Therefore, it is the sum of these two diagrams that corresponds to one of the six improved helicity configurations. Both double boxes in fig. 4 involve gluons, fermions and scalars running in one of the loops and only gluons running in the remaining loop. The necessary tree level amplitudes are given by
$$A_3^{tree}(p^{},q^{},r^+)=i\frac{pq^3}{qrrp}\left(\frac{qr}{qp}\right)^a,A_3^{tree}(p^+,q^+,r^{})=i\frac{[pq]^3}{[qr][rp]}\left(\frac{[qr]}{[qp]}\right)^a,$$
where $`a=0`$ for gluons, $`a=1`$ for fermions and $`a=2`$ for scalars. After summation over the two configurations, we obtain
$$i\frac{(\alpha \beta )^4}{\gamma },$$
where
$$\begin{array}{cc}& \alpha =q_11q_2p_2[pq][4q_1],\beta =1q_2p_2l[lp][q4]\hfill \\ & \gamma =q_2q_1lq_2[ql][q_1q]1q_2p_2l[lp][q4]q1q_2p_2[4q_1][pq].\hfill \end{array}$$
By using momentum conservation along the lines of eq. (3.1), we can simplify (3.1) to obtain $`A_{tree}st^2`$. The integral of the measure factors out and cancels against the denominator in (3.1) to give (3.1).
Thus, we find that the coefficients of the four-gluon amplitude double boxes can be calculated by studying hepta-cuts. Of course, this is not enough to claim that this is the full answer. One still has either to prove that the answer has all the correct discontinuities across all branch cuts, which was done in , or to prove that the basis of integrals is given entirely by double boxes. Unfortunately, the basis of integrals is not known at two loops.
For four gluons, even though the number of integration variables is greater that the number of the delta-functions in a hepta-cut, no integration has to be performed. We find that this is not the case if the number of external gluons is greater than four. We will see in the next section that already in the case of five-gluon amplitude, the product of the corresponding tree-level amplitudes does depend on the loop momenta and cannot be pulled out of the integral.
4. Octa-Cuts of Two-Loop Amplitudes
In the introduction we distinguished between two different kinds of double box scalar integrals. In the first class, the two boxes share a propagator while in the second class they only share a vertex, see fig. 1a and fig. 1b respectively. In this section, we show how one can use octa-cuts to compute the coefficient of a certain subset of the first class and the coefficients of all integrals of the second class, which we called split double boxes.
4.1. Double-Box Scalar Integrals
Let us start with the double boxes that have seven propagators. We will show that when at least one of the two boxes has two adjacent three particle vertices then there is an extra propagator-like singularity that can be cut. This produces one more delta-function which together with the hepta-cut of the previous section completely localizes the cut integral. Even though we concentrate only on the planar configurations, exactly the same logic can be applied for non-planar configurations as well. Consider an arbitrary double-box configuration shown in fig. 5. The corresponding hepta-cut integral is
$$=\frac{d^4p}{(2\pi )^4}\frac{d^4q}{(2\pi )^4}\delta (p^2)\delta ((pk_1)^2)\delta ((pk_1k_2)^2)\delta ((p+q+k_6)^2)\delta (q^2)\delta ((qk_5)^2)\delta ((qk_4k_5)^2).$$
Fig. 5: An arbitrary double-box configuration.
Let us perform, for example, the $`p`$-integration. The integral over $`p`$,
$$_p=d^4p\delta (p^2)\delta ((pk_1)^2)\delta ((pk_1k_2)^2)\delta ((p+q+k_6)^2),$$
is localized and the answer is
$$_p=\frac{2}{(k_1+k_2)^2(k_1+k_6+q)^2\rho },$$
where
$$\begin{array}{cc}& \rho =\sqrt{12(\lambda _1+\lambda _2)+(\lambda _1\lambda _2)^2},\hfill \\ & \lambda _1=\frac{k_1^2(k_3+k_4+k_5q)^2}{(k_1+k_2)^2(k_1+k_6+q)^2},\lambda _2=\frac{k_2^2(k_6+q)^2}{(k_1+k_2)^2(k_1+k_6+q)^2}.\hfill \end{array}$$
The crucial observation is that when
$$\rho =1,$$
$``$ acquires an extra propagator-type singularity, i.e. $`1/(k_1+k_6+q)^2`$. We can formally cut the new propagator by replacing it with a delta-function creating an eighth cut. In other words, after performing the $`p`$-integration we end up with following integral over $`q`$ (we omit the overall $`q`$-independent factor)
$$_q=d^4q\delta (q^2)\delta ((qk_4)^2)\delta ((qk_3k_4)^2)\frac{1}{(k_1+k_6+q)^2}.$$
This integral looks like a triple cut of the following effective box
Fig. 6: Effective box that arises after a quadruple cut is used to localize the $`p`$ integral. The momentum flowing along the uncut line is $`q+k_1+k_6`$.
Note that the momentum flowing along the uncut line is exactly $`q+k_1+k_6`$. From this viewpoint it is natural to cut the remaining propagator. Note that this procedure localizes the $`q`$-integral. Then it is straightforward to write down the coefficients of such double-box integrals. They are given by
$$c_\alpha =\frac{i}{|๐ฎ|}\underset{h,J_1,J_2,๐ฎ}{}(n_{J_1}n_{J_2}A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree})_h,$$
where the sum over $`h`$ is the sum over all helicity configurations, the sums over $`J_1`$ and $`J_2`$ are the sums over all particles that can propagate in both loops, $`๐ฎ`$ is the set of all solutions for the internal lines of the following system of equations
$$\begin{array}{cc}& p^2=0,(pk_1)^2=0,(pk_1k_2)^2=0,(p+q+k_6)^2=0,\hfill \\ & q^2=0,(qk_5)^2=0,(qk_4k_5)^2=0,(k_1+k_6+q)^2=0,\hfill \end{array}$$
and $`|๐ฎ|`$ is the number of solutions. This expression is analogous to the formula for one-loop coefficients of box integrals . It is important to remember that this discussion is valid if
$$\rho =1,\lambda _1=\lambda _2=0.$$
Otherwise, the singularity $`1/(k_1+k_6+q)^2`$ will be replaced by a more complicated one which is not propagator-like, as it can easily be seen from eq. (4.1). The conditions given in (4.1) are satisfied if a given box has two adjacent three-particle vertices. It easy to check that this is always the case if the number of gluons is less than seven. This means that every double-box coefficient of any gluon amplitude with less than seven external lines is given by eq. (4.1). The first double-box configuration where eq. (4.1) is not satisfied appears when the number of external gluons is seven and is shown in fig. 7.
Fig. 7: The simplest double-box configuration for which the conditions in (4.1) are not satisfied.
However, even if the number of external gluons is greater than six, there are double-box configurations for which eq. (4.1) is satisfied. In such cases the eighth cut still exists and eq. (4.1) is still correct.
In fact, eq. (4.1) requires some additional explanations. Note that existence of the effective box in fig. 6 implies that either the momentum $`l`$ or the momentum $`p_1`$ in fig. 5 vanishes. In Minkowski space, this would mean that some tree level amplitudes in eq. (4.1) vanish. Moreover, in Minkowski space, the system of equations (4.1) does not have solutions, which means that we cannot see the singularities under consideration. Therefore, it is not surprising that eq. (4.1), at least naively, is meaningless in Minkowski space. In order to see the new kind of singularities, we have to analytically continue all momenta to signature $`(++)`$. But in signature $`(++)`$, the statement that a tree amplitude vanishes when one of the incoming or outgoing momentum vanishes is not correct. Each tree amplitude is constructed by using spinors. When one of the incoming or outgoing $`(++)`$ momenta vanishes, it is impossible to determine its spinors components even up to rescaling. This leaves the amplitude undetermined. For example, assume that we have a helicity configuration containing a three-gluon amplitude $`A(p^{},p_1^{},k_1^+)`$. It is given by
$$A(p^{},p_1^{},k_1^+)=\frac{p_1p^3}{pk_1k_1p_1}.$$
If $`p_1`$ vanishes, the spinor $`\lambda _{p_1}`$ cannot be uniquely determined. In fact, $`\lambda _{p_1}`$ is not uniquely defined even for non-zero $`p_1`$ as it is defined up to rescaling. However, when $`p_1=0`$ the freedom in not being able to determine $`\lambda _{p_1}`$ becomes much larger. One can always say that $`p_1=0`$ implies that $`\stackrel{~}{\lambda }_{p_1}=0`$ and $`\lambda _{p_1}`$ is arbitrary. This means that $`A(p^{},p_1^{},k_1^+)`$ becomes arbitrary. Therefore, the numerator in eq. (4.1) is a discontinuous function of momenta and we have to give a prescription on how to define it as $`l`$ or $`p_1`$ goes to zero. The natural way to define it is as follows. Consider first the loop with momentum $`p`$. Let $`A_{(1)}^{tree},A_{(2)}^{tree},A_{(3)}^{tree}`$ and $`A_{(4)}^{tree}`$ be the four tree amplitudes which depend on $`p`$. Assuming that they are all non-zero, we can solve the first four $`p`$-dependent equations in (4.1) to determine $`p`$ as a function of the external momenta and $`q`$ and then evaluate the product $`A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}`$ on these solutions. We claim that this product can be simplified in such a way that it is a well-defined function when the constraint $`(k_1+k_6+q)^2=0`$ is imposed. Below, we will present a few examples that show that this is indeed the case. Having found the product $`A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}`$ as a function of the external momenta and $`q`$, we then multiply it by the remaining two tree amplitudes $`A_{(5)}^{tree}`$ and $`A_{(6)}^{tree}`$ and evaluate the product on the solution of the remaining four equations in (4.1). We propose this as a method for calculating double-box coefficients provided conditions (4.1) are fulfilled.
A Subtlety
There is one important subtlety we have to discuss before presenting examples. Consider a helicity configuration with two adjacent three-particle vertices, one of which depends only on internal momenta and the other one depends on both internal and external momenta, with both vertices having the same helicity configuration. For example, consider the configuration in fig. 8.
Fig. 8: This helicity configuration is non-zero only if $`\lambda _q\lambda _1`$.
This configuration is non-zero only if $`\lambda _q\lambda _1`$. Therefore, the integral over $`p`$
$$d^4p\delta (p^2)\delta ((pk_1)^2)\delta ((pk_1k_2)^2)\delta ((p+q)^2)A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}$$
must be proportional to $`\delta ((k_1+q)^2)`$. In other words, the integral lacks the extra propagator-like singularity and therefore it does not contribute to the octa-cut.
4.2. Split Double-Box Scalar Integrals
When the number of gluons is greater than five, a new kind of double box integrals can appear. These were introduced in the introduction in fig. 1b. For the readerโs convenience we depict them again in fig. 9. This double box scalar integrals are such that the two boxes only share a vertex and not a propagator. This is why we will call them split double boxes.
Fig. 9: Generic split double box configuration.
The coefficients of the split double boxes are easy to compute. Since they have eight propagators and the two loop integrations are completely independent, it is straightforward to consider two quadruple cuts or equivalently an honest octa-cut. This produces eight delta-functions that localize both loop integrals. Let us see this in more detail. The quadruple cut in the $`q`$-loop fixes $`q`$ to be a solution to the following equations,
$$q^2=0,(qk_6)^2=0,(qk_5k_6)^2=0,(q+k_7)^2=0,$$
while a quadruple cut in the $`p`$-loop fixes $`p`$ to be a solution of
$$p^2=0,(pk_2)^2=0,(pk_1k_2)^2=0,(p+k_3)^2=0.$$
Each set of equations has two solutions. The coefficient of a split double-box scalar integral is then given by
$$c=\frac{1}{4}\underset{h,J_1,J_2,๐ฎ}{}(n_{J_1}n_{J_2}A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree}A_{(7)}^{tree})_h$$
where we have used that the number of solutions is $`4`$.
5. Examples
In this section we consider several examples of coefficients calculated by using octa-cuts. All of them are coefficients of scalar double boxes with seven propagators. We consider four-, five- and six-gluon amplitudes. For six-gluons we study a non-MHV amplitude with adjacent negative helicity gluons.
5.1. Four-Gluon Amplitude $`A^{2loop}(,,+,+)`$ Revisited
As a first example, let us reconsider the octa-cut of the four-gluon amplitude from section 3 in the $`s`$-channel. The octa cut in the $`t`$-channel is analogous. The additional propagator that we cut is $`\frac{1}{(q+K_1)^2}`$. See fig. 3 for notation. Taking into account the subtlety in the previous section, there are only four helicity configurations that contribute. All of them give the same answer $`A_4^{tree}s^2t`$. On the other hand, the number of solutions to eqs. (4.1) can be shown to be four. This gives $`c_s=A_4^{tree}s^2t`$ as in eq. (3.1). Note that the product of the four tree level amplitudes depending on the internal momentum $`p`$ is given by $`12^2[qq_2]^2`$ (see eq. (3.1)). This expression does not have any ambiguity in the presence of the eighth delta-function $`\delta ((q+K_1)^2)`$.
5.2. Five-Gluon Amplitude $`A^{2loop}(,,+,+,+)`$
As a next example, let us calculate the coefficient of the following five-gluon double-box configuration.
Fig. 10: A double box integral of the five-gluon amplitude $`A(1^{},2^{},3^+,4^+,5^+)`$.
In this case, there are two helicity configurations to consider. Both of them can be shown to give the same contribution. We will consider the helicity configuration shown in fig. 11.
Fig. 11: One of the two possible helicity configurations contributing to the coefficient of the integral of fig. 10.
Note that only gluons can propagate in both loops. The product of the six tree level amplitudes is as follows
$$12^2[qq_2]^2\frac{[3q_1]^3}{[q_1q_2][q_23]}\frac{qq_1^3}{q_14455q},$$
where the first two factors come from the four vertices on the right. The computation leading to the first two factors was done in the previous section in eq. (3.1). By using momentum conservation along the lines of eq. (3.1), eq. (5.1) can be reduced to
$$\frac{12^2s^2}{3445}\left([53]+[43]\frac{q4}{q5}\right),$$
Now we impose the delta-function $`\delta ((K_1+q)^2)`$. It has two possible solutions, $`\lambda _q\lambda _1`$ or $`\stackrel{~}{\lambda }_q\stackrel{~}{\lambda }_1`$. It is not hard to show that if we choose $`\stackrel{~}{\lambda }_q\stackrel{~}{\lambda }_1`$ then the expression in eq. (5.1) vanishes. Therefore, the only relevant solution is $`\lambda _q\lambda _1`$. Then the octa-cut becomes
$$\frac{12^2s^2}{3445}\left([53]+[43]\frac{14}{15}\right)๐\mu .$$
Taking into account that the system of equations (4.1) has four solutions, we find that the corresponding coefficient is
$$c_1^{(5)}=\frac{2i}{4}\frac{12^2s^2}{3445}\left([53]+[43]\frac{14}{15}\right)=\frac{1}{2}A_5^{tree}s^2t,$$
where $`s=(K_1+K_2)^2`$ and $`t=(K_2+K_3)^2`$ and $`A_5^{tree}`$ is the tree-level five-gluon amplitude
$$A_5^{tree}=i\frac{12^3}{23344551}.$$
Let us consider one more example. Let us calculate the coefficient of the five-gluon double-box configuration shown in fig. 12.
Fig. 12: Second example of a double-box configuration in the five-gluon amplitude $`A(1^{},2^{},3^+,4^+,5^+)`$.
There are two helicity configurations to consider. However, one of them gives the zero answer. The non-zero contribution comes from the helicity configuration shown in fig. 13.
Fig. 13: The only non-vanishing helicity configuration contributing to the coefficient of the double box integral of fig. 12.
Note that only gluons can propagate in both loops. The number of solutions to eqs. (4.1) can be shown to be two. Then the corresponding coefficient is given by
$$\begin{array}{cc}\hfill c_2^{(5)}=& \frac{i}{2}\frac{q_1q_2^3}{q_233q_1}\frac{[q_14]^3}{[4q][qq_1]}\frac{pl^3}{lqq55p}\frac{[q_2l]^3}{[lp_2][p_2q_2]}\hfill \\ & \frac{p_12^3}{2p_2p_2p_1}\frac{[p_1p]^3}{[p1][2p_1]}.\hfill \end{array}$$
Using momentum conservation and the fact that $`\lambda _q\lambda _4`$, eq (5.1) can be simplified to give
$$c_2^{(5)}=\frac{1}{2}A_5^{tree}stu,$$
where $`u=(K_3+K_4)^2`$. All other double-box coefficients of the five-gluon amplitude can be found by analogous calculations.
5.3. Six-Gluon Amplitude $`A^{2loop}(,,,+,+,+)`$
As our next example, let us calculate the coefficient of the six-gluon next-to-MHV double-box configuration shown in fig. 14. The additional singularity that we cut is again $`\frac{1}{(q+K_1)^2}`$. There are two helicity configurations to consider, both giving the same answer. In both of them only gluons can propagate in both loops. Let us describe the calculation of the one depicted in fig. 15.
Fig. 14: A double-box integral of the six-gluon non-MHV amplitude $`A(1^{},2^{},3^{},4^+,5^+,6^+)`$.
Fig. 15: One of the two helicity configurations contributing to the coefficient of the scalar double box of fig.14.
The product of the four $`p`$-dependent tree amplitudes gives (see eq. (3.1))
$$12^2[qq_2]^2.$$
Then the numerator of eq. (4.1) becomes
$$\begin{array}{cc}& 2i(A_{(1)}^{tree}A_{(2)}^{tree}A_{(3)}^{tree}A_{(4)}^{tree}A_{(5)}^{tree}A_{(6)}^{tree})=\hfill \\ & 12^2[qq_2]^2\frac{q_23^3}{344q_1q_1q_2}\frac{[56]^3}{[6q][qq_1][q_15]},\hfill \end{array}$$
where the factor of two comes from two helicity configurations. Using momentum conservation similar to eq. (3.1), eq. (5.1) can be simplified as follows
$$i\frac{12^2[56]^33|5+6|q]^2q_23}{34[6q]4|5+6|q]q_2|3+4|5]}.$$
Now we impose the last condition $`(q+K_1)^2=0`$. This equation has two solutions. We can either have $`\stackrel{~}{\lambda }_q\stackrel{~}{\lambda }_1`$ or $`\lambda _q\lambda _1`$. Both solutions give non-zero contributions. The first solution yields
$$i\frac{12^2[56]^33|5+6|1]^223}{34[61]4|5+6|1]2|3+4|5]}$$
while the second solution yields
$$i\frac{12^2[56]^356[42]^23|(1+2)(5+6)(3+4)|2]}{[23]5|6+1|2][5|(3+4)(1+2)(5+6)(3+4)|2]}.$$
Taking into account that the system (4.1) in this case has four solutions, the double-box coefficient becomes
$$\begin{array}{cc}\hfill c_1^{(6)}=& \frac{i}{2}(\frac{12^2[56]^33|5+6|1]^223}{34[61]4|5+6|1]2|3+4|5]}+\hfill \\ & \frac{12^2[56]^356[42]^23|(1+2)(5+6)(3+4)|2]}{[23]5|6+1|2][5|(3+4)(1+2)(5+6)(3+4)|2]}).\hfill \end{array}$$
Let us consider one more example. Let us calculate the coefficient of the six-gluon double-box configuration shown in fig. 16.
Fig. 16: Second example of a six-gluon double-box integral of the six-gluon non-MHV amplitude $`A(1^{},2^{},3^{},4^+,5^+,6^+)`$.
Fig. 17: The only non-vanishing helicity configuration contributing to the coefficient of the scalar double-box integral of fig.16.
In this case, there is only one helicity configuration contributing to the octa-cut. It is shown in fig. 17. Only gluons can propagate in both loops. The system of equations (4.1) has two solutions. Therefore, the corresponding double-box coefficient is given by
$$\begin{array}{cc}\hfill c_2^{(6)}=& \frac{i}{2}\frac{[p_1p]^3}{[p_11][1p_1]}\frac{p_12^3}{2p_2p_2p_1}\frac{pl^3}{lqqp}\frac{[q_2l]^3}{[lp_2][p_23][3q_2]}\times \hfill \\ & \frac{[q_16]^3}{[6q][qq_1]}\frac{q_1q_2^3}{q_24455q_1}.\hfill \end{array}$$
By using the first seven equations in (4.1), we can simplify (5.1) as follows
$$c_2^{(6)}=\frac{i}{2}\frac{u^3s1|q|6]}{[12][23]45564|qK_4^{[3]}|3]},$$
where
$$u=(K_4+K_5+K_6)^2,s=(K_1+K_2)^2.$$
Now we consider the last equation $`(q+K_1)^2=0`$. From fig. 17, it follows that $`\lambda _q`$ has to be proportional to $`\lambda _6`$. Therefore, $`\stackrel{~}{\lambda }_q`$ has to be proportional to $`\stackrel{~}{\lambda }_1`$. Using momentum conservation, one can find that
$$q=\frac{u}{6|4+5|1]}\lambda _6\stackrel{~}{\lambda }_1.$$
Substituting this into eq. (5.1), we obtain
$$c_2^{(6)}=\frac{i}{2}\frac{u^4st}{[12][23]45564|5+6|1]6|4+5|3]},$$
where
$$t=(K_1+K_6)^2.$$
All other double-box coefficients can be computed by similar calculations.
The calculation of the coefficient $`c_2^{(6)}`$ can be generalized for the configuration considered in fig. 18.
Fig. 18: An infinite family of $`n`$-gluon double box scalar integrals.
There is only one helicity configuration that contributes to the octa-cut. A similar calculation gives
$$\begin{array}{cc}\hfill c^{(n)}=& \frac{i}{2}\frac{u^4st}{[12][23]\mathrm{}[m1m]m+1m+2\mathrm{}n1n}\times \hfill \\ & \frac{1}{m+1|K_{m+1,m+2,\mathrm{},n}|1]n|K_{m+1,m+2,\mathrm{},n}|m]},\hfill \end{array}$$
where $`s`$ and $`t`$ are given by eq. (5.1) and u is given by
$$u=(K_{m+1,m+2,\mathrm{},n})^2=(K_{m+1}+K_{m+2}+\mathrm{}+K_n)^2.$$
6. Application to Three and Higher Loops
Ideas presented in the previous sections can be applied to higher loops. Let us consider triple-box configurations appearing at three-loops. The configurations we consider are obtained from the double-box configurations at two loops by adding three new propagators to form the third loop. This way, one can produce a ladder diagram as in fig. 19a or a double box with a pentagon as in fig. 19b. We make a slight abuse of terminology and call both kind of configurations triple-box diagrams.
Fig. 19: Three-loop triple box configurations. $`(a)`$ A triple box ladder integral. $`(b)`$ A double box with a pentagon.
Every triple box contains ten propagators<sup>7</sup> Of course, for a large enough number of gluons one can also find split triple boxes which can have 11 or 12 propagators.. Therefore it is natural to start with a ten-particle cut. This produces ten delta-functions whereas the number of integration variables is twelve. However, it follows from our previous analysis that box configurations develop additional propagator-like singularities which can also be cut (replaced by their discontinuities). A triple-box configuration naturally admits two extra propagator-type singularities which allows us to consider twelve-particle cuts. Therefore, it should be possible to completely localize all momentum integrals, at least if the number of gluons is not big enough. Then it is straightforward to write down an expression for the coefficients analogous to eq. (4.1). Obviously, a similar analysis can be performed at any number of loops. It is interesting to mention that at three loops, some of the triple boxes that enter in the calculation of the four-gluon amplitude are not scalar triple boxes . This means that the numerator in the integrand is not one but an inverse propagator. See fig. 20. This was also found to be the case for higher loops .
We start our discussion with the analysis of ladder diagrams which allow a straightforward generalization of our discussion in section 4. Then we turn to the triple-box integrals where one of the โboxesโ has five propagators and see how our analysis of singularities realizes the phenomenon mentioned above. For concreteness, we will concentrate on the four-gluon amplitude though an identical analysis can be performed regardless of the number of external lines. Let us start with the triple-box configuration in fig. 21.
Fig. 20: Schematic representation of a modified โtriple boxโ integral given in .
Fig. 21: A ladder triple-box configuration with four external gluons.
Fig. 22: One of the helicity configurations contributing to the coefficient of the ladder triple-box in fig.21.
In this case, there are twelve helicity configurations, all giving the same answer. It is enough to consider one of them, for example the one in fig. 22. Similarly to the two-loop case, it is enough to consider the ten-particle cut because the measure integral factors out and cancels. Then the corresponding coefficient is given by
$$d_1=i^3\stackrel{~}{A}_4^{tree}s^2\stackrel{~}{t}^2\frac{r_1r_2^3}{r_233r_1}\frac{[r_14]^3}{[4r][rr_1]},$$
where $`\stackrel{~}{A}_4^{tree}`$ is the tree level four-gluon amplitude with the external lines $`(K_1^{},K_2^{},r_2^+,r_1^+)`$, $`s=(K_1+K_2)^2`$ and $`\stackrel{~}{t}=(K_1r)^2`$. In eq. (6.1), we used that the product of the six tree amplitudes was computed in section 3. Using momentum conservation, eq. (6.1) can be simplified to give
$$d_1=iA_4^{tree}s^3t.$$
where $`t=(K_2+K_3)^2`$. This coincides with the answer from .
Fig. 23: A triple-box configuration.
Now let us consider the configuration in fig. 23. Note that one of the loops has five propagator and this is why we said that the configuration was a double box with a pentagon. This is the basic reason why the integral is not a scalar box integral as we will see below. Let us start our analysis with the measure integral
$$\begin{array}{cc}\hfill _0=& d^4ld^4qd^4p\delta (l^2)\delta ((lK_3)^2)\delta ((lK_3K_4)^2)\delta ((lK_3K_4K_1)^2)\hfill \\ & \delta ((pK_2l)^2)\delta ((pK_2)^2)\delta ((qK_1)^2)\delta ((p+q)^2)\delta (p^2)\delta (q^2)\hfill \end{array}$$
and perform the integration over $`p`$ and $`q`$. After integrating over $`p`$ we obtain (up to the momenta-independent factor)
$$\frac{1}{(K_2+q)^2(K_2+l)^2}.$$
Then we cut the โpropagatorโ $`1/(K_2+q)^2`$. This gives the fourth delta-function which allows us to perform the integration over $`q`$. This produces (again, we ignore the momenta-independent factors) one more factor of $`(K_2+l)^2`$ in the denominator. Thus, we obtain that this triple box configuration has a singularity
$$\left(\frac{1}{(K_2+l)^2}\right)^2.$$
Let us calculate the coefficient of this triple box integral by calculating the product of the eight-gluon amplitudes. For every non-vanishing helicity configuration, the product of the six gluon amplitudes on the right is the coefficient $`c_t`$ of the two-loop four-gluon amplitude studied in section 3 and it is given by
$$A_4^{tree}st^2,$$
where $`A_4^{tree}`$ is the tree four-gluon amplitude with external lines $`(K_1^{},K_2^{},l^+,(lK_3K_4)^+)`$ and $`t^{}`$ is given by
$$t^{}=(K_2+l)^2.$$
Therefore, in attempting to calculate the integral of the product of the tree amplitudes, the singularity (6.1) cancels out. This means that the coefficient of the triple box in fig. 23 is zero. This is in agreement with results of . On the other hand, the amplitude $`A_4^{tree}`$ has a factor
$$\frac{1}{2l}=\frac{[2l]}{(K_2+l)^2}.$$
This indicates that the actual diagram has a singularity $`\frac{1}{(K_2+l)^2}`$. In order to account for this singularity we have to introduce a slightly modified triple box integral schematically shown in fig. 20. The basic idea is to multiply in the numerator by $`(K_2+l)^2`$ in order to cancel one of the power in (6.1) and get the correct $`1/(K_2+l)^2`$ singular behavior. This shows that this triple-box integral should be in the list of scalar integrals of the amplitude under study. This is completely consistent with results of . Now we can cut this โpropagatorโ and completely localize the integral. Note that the combination
$$[2l]\delta ((K_2+l)^2)$$
is not necessarily zero. This just means that we have to choose the solution
$$\lambda _l\lambda _2.$$
The coefficient of this modified box is now straightforward to compute. The answer is
$$d_2=iA_4^{tree}s^2t.$$
This coincides with the corresponding coefficient from . All other coefficients of this amplitude can be found in a similar manner. Even though we concentrated on the four-gluon amplitude, we could do the same analysis for any amplitude admitting additional cuts.
As mentioned in the introduction, the basis of integral for $`๐ฉ=4`$ $`L`$-loop amplitudes of gluons is not known except for $`L=1`$. One can imagine that a more systematic analysis along the lines of the discussion presented in this section might give a way of obtaining such a basis. It would be interesting to explore this direction in the future.
7. Conclusion
In this paper, we observed that certain scalar double-box integrals which appear at two loops in $`๐ฉ=4`$ Yang-Mills theory possess hidden singularities. Such singularities are manifest after a quadruple cut is performed on one of the boxes. The end result is that one can straightforwardly calculate the coefficient of such integrals by an octa-cut which localizes the cut integral. The form of the coefficient is universal and it is given by the product of a certain number of tree-level amplitudes. This technique is applicable to all scalar double box integrals in amplitudes with less than seven external gluons and to a large subset of double box integrals for seven or more external gluons. The basis of integrals at two loops is not known in general. For four gluons the amplitude is given in terms of only scalar double-box integrals. If it turns out that the basis of integrals for five- and six-gluon amplitudes is also given by scalar double box integrals, then our technique gives a simple way of computing all those amplitudes for any helicity configuration. We also argued that this technique can be applied to higher loop amplitudes. At three loops we found that our technique can be easily extended to compute the coefficient of ladder diagrams. For the class of diagrams with a pentagon, our method shows that the coefficient of scalar integrals is zero and naturally gives the modified integral for which the coefficient does not vanish.
Acknowledgments
The authors would like to thank Z. Bern, R. Britto, L. Dixon, B. Feng, H. Osborn, M. Spradlin and A. Volovich for comments on the first version of the manuscript. Research of E.I.B is supported by NSF grant PHY-0070928. Research of F.C. is supported in part by the Martin A. and Helen Chooljian Membership at the Institute for Advanced Study and by DOE grant DE-FG02-90ER40542. Research at the Perimeter Institute is supported by funds from NSERC of Canada.
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warning/0506/cond-mat0506786.html | ar5iv | text | # On the Limits of Analogy Between Self-Avoidance and Topology-Driven Swelling of Polymer Loops
## I Introduction: Formulation of the Problem
The last few years have seen significant work addressing the effects of knotting on looped polymer chains. Of interest to mathematicians and physicists for good part of nineteenth and most of the twentieth centuries, knots were first seen by W. Thomson as a way to understand the nature of atoms knotted\_vorticies , and more recently as the basis for string theory. On the biological front, knots have been observed in, JBiolChem\_1985 ; probability\_DNA\_knotting , and tied into, tie\_knot\_into\_DNA-Japan ; tie\_knot\_into\_DNA-Quake , strands of DNA. Additionally, topoisomerases - proteins which act to alter the topological state of DNA - are quite common and play a significant role in cellular processes.
The requirements a knot imposes on a strand are hard to formulate in a simple way, as โinteractionsโ between neighboring strands can require highly non-local changes in the coilโs conformation to maintain topological state.
That said, the most obvious effect knotting has on a loop is in the size, commonly measured in terms of radius of gyration, $`R_g^2`$. For instance, the loop topologically equivalent to a circle, called a trivial or $`0_1`$ knot in professional parlance, is on average *larger* than the loop of the same length with any other topology. In other words, a trivial loop is larger than the phantom loop, the latter representing topology-blind average over all loops of a certain length: $`R_g^2_{triv}>R_g^2_{phantom}`$. This topology-driven swelling is operational even for very thin polymers, in the limit when volume exclusion has no effect on polymer coil size. In this case, the phantom loopโs size (which is, once again, average over all topologies) scales as $`N^{1/2}`$, while the trivial loop is larger not merely because of a larger prefactor, but because of a larger scaling exponent, its size scales as $`N^\nu `$, where $`\nu >1/2`$. The conjecture, formulated a long time ago desCloizeaux\_conj , supported by further scaling arguments Quake1 ; AG\_pred , and consistent with recent simulation data Deguchi\_2003 ; swiss\_PNAS ; Nathan\_PNAS , specifies that the scaling exponent $`\nu `$ describing topology-driven swelling of a trivial loop is exactly the same as the Flory exponent Flory\_excluded\_volume , which describes swelling driven by the self-avoidance (or excluded volume): $`\nu 0.5893/5`$.
Equality of scaling exponents for the two cases reflects the similarity of fractal properties for these systems at very large $`N1`$, because topological constraints result in self-avoidance of blobs on all length scales above a certain threshold AG\_pred . As we understand much about self-avoidance Madras , and next to nothing about knots, we would like to exploit the analogy to see if it yields any insights into knots. Specifically, it is tempting to look at the dependence of the unknotted loop size, $`R_g^2_{triv}`$, on the number of segments, $`N`$, not only in the asymptotic scaling regime of very large $`N`$, but also the corrections to scaling at not-so-large $`N`$. This is particularly important from a practical standpoint, because the asymptotic scaling limit is barely accessible computationally, and what one really computes is the value of $`R_g^2_{triv}`$ at rather moderate $`N`$. Systematic comparison of $`N`$-dependencies of $`R_g^2`$ for (trivial) knots and self-avoiding polymers over the wide range of $`N`$ is the goal of this paper.
We show that although large $`N`$ scaling appears to be identical for trivial knots and excluded volume polymers, their respective approach to the asymptotic regime is different. This points obviously to the limited character of the analogy between the two mechanisms of swelling, due to volume exclusion and due to topological constraints.
The plan of the paper is as follows. We start from a brief summary of the main results for self-avoiding polymers. Although these results are widely known, we restate them in the form most suitable for our purposes. Next, we present some heuristic analytical arguments to shed light on why trivial knots may behave differently then their excluded volume counterparts. With this insight in mind, we present our detailed computational data on the $`N`$-dependence of $`R_g^2_{triv}`$ over the wide range of $`N`$. To obtain data with the necessary degree of accuracy, it is necessary to make sure that our method of generating loops is ergodic and unbiased. Although this aspect is of decisive importance, it is purely technical, and thus it is relegated to the Appendix. Up to about section II.3 we mostly review the known results, starting from section II.4, we present our new findings.
## II Preliminary Considerations
### II.1 Swelling driven by self-avoidance: an overview
To make our work self-contained we now offer a brief review of the results for the scaling of excluded volume polymers (see further details in Madras ; AG\_Red ; Yamakawa ). We should emphasize from the beginning that the main properties of the excluded volume polymer are valid also for loops Casassa\_rings\_1965 . The simplest model for excluded volume is a system in which $`N`$ beads, each of volume $`b`$, are placed along a loop with mean separation $`\mathrm{}`$. All other forms of excluded volume, e.g. freely jointed stiff rods, worm-like filaments, etc., can be mapped to this simple rod-bead model (see.e.g., AG\_Red ). There are two scaling regimes, with crossover at the length
$$N^{}\left(\mathrm{}^3/b\right)^2.$$
(1)
In terms of $`N^{}`$, the mean squared gyration radius $`R_g^2`$ can be written as $`R_g^2=\mathrm{}^2N\rho \left(z\right)`$, where the swelling factor $`\rho `$ depends on the single variable $`z=\sqrt{N/N^{}}`$. For classical polymer applications, the large $`z`$ regime is most interesting. $`\rho (z)`$ has a branch point singularity in infinity, its large $`z`$ asymptotics are dominated by the factor $`z^{2\nu 1}`$; however, if we write $`\rho (z)=z^{2\nu 1}\varphi (z)`$, then $`\varphi (z)`$ is analytical in infinity and can be expanded in integer powers of $`1/z`$. Accordingly, the large $`N`$ asymptotics of $`R_g^2`$ follow:
$$R_g^2|_{NN^{}}\mathrm{}^2N^{2\nu }A\left[1+k_1\left(\frac{N^{}}{N}\right)^{1/2}+k_2\left(\frac{N^{}}{N}\right)^1+\mathrm{}\right].$$
(2)
Conversely, in the region $`NN^{}`$, the approximation for $`R_g^2`$ is afforded by the fact that $`\rho (z)`$ is analytical at small $`z`$ and can be expanded in integer powers of $`z`$:
$$R_g^2|_{1NN^{}}\mathrm{}^2N\frac{A^{}}{12}\left[1+k_1^{}\left(\frac{N}{N^{}}\right)^{1/2}+k_2^{}\left(\frac{N}{N^{}}\right)^1+\mathrm{}\right],$$
(3)
where prefactor $`A^{}`$ should be equal to unity (which explains why we did not absorb the factor of $`1/12`$ into $`A^{}`$). Note that the latter result is an intermediate asymptotics, which means the corresponding region exists only so long as $`N^{}1`$ is large, which means excluded volume is sufficiently small.
### II.2 Swelling driven by topology: cross-over length
With this brief summary of results in mind we now set forward, intending to systematically compare the computational results for the behavior of trivial knots to the well-understood polymer with excluded volume.
To look at the analogy between self-avoiding polymers and trivial knots, it is useful to start, AG\_pred , by identifying the cross-over length for knots, an analog of $`N^{}`$ (1), which we call $`N_0`$. For knots, it is natural to identify the cross-over value of $`N`$ with the so-called characteristic length of random knotting, $`N_0`$; the latter quantity is known as the characteristic length of the exponential decay of probability, $`w_{triv}(N)`$, of formation of a trivial knot upon random closure of the polymer ends koniaris\_muthu\_N0 : $`w_{triv}\mathrm{exp}(N/N_0)`$. Depending on the specifics of the model used, koniaris\_muthu\_N0 ; Nathan\_PNAS ; Deguchi\_universality\_1997 , the critical length varies subtly around $`N_0300`$. It is also clear qualitatively AG\_pred and seen computationally Nathan\_PNAS that this $`N_0`$ is about the length at which topological effect on loop swelling crosses over from marginality at $`N<N_0`$ to significance at $`N>N_0`$. In particular, it is at $`N>N_0`$ that the trivial knot begins to swell noticeably beyond the size of the phantom polymer Nathan\_PNAS .
### II.3 Swelling driven by topology: above the cross-over
A number of groups reported observation of the power $`\nu 3/5`$ in the scaling of trivial Deutsch ; Deguchi\_2003 ; swiss\_PNAS ; Nathan\_PNAS and other topologically simple Deguchi\_2003 ; swiss\_PNAS ; Nathan\_PNAS knots in the region $`N>N_0`$.
In the works Deguchi\_2003 ; swiss\_PNAS ; swiss\_macromolecules , following the idea suggested in RG\_style\_fitting , the $`N`$ dependence of $`R_g^2_{triv}`$ was fitted to the formula similar to equation (2) for self-avoiding polymers. No attempt was made at physical interpretation of the best fit values of the three coefficients ($`A`$, $`k_1`$, $`k_2`$) or the region of $`N`$ where the fit was examined. In this sense, fitting with equation (2) was only used as an instrument to find the scaling exponent $`\nu `$, which in these works was found to be strikingly consistent with the expected value of the self-avoidance exponent. A puzzling aspect of the situation is that, particularly in the work swiss\_PNAS , the data was fit to equation (2) not only in the region $`N>N_0`$, but across the crossover, starting from about $`N_0/3`$ to about $`3N_0`$ (see also swiss\_macromolecules ).
At present we are aware of no studies which provide a detailed comparison of excluded volume and trivial knotting at modest $`N<N_0`$. Seeking to further appraise the analogy between trivial knotting and excluded volume, in the present work we address the two systems in the region below their respective crossovers.
### II.4 Swelling driven by topology: below the cross-over
Formula (3) is the result of perturbation theory Yamakawa , in which conformations with overlapping segments represent a small part of conformational space and their exclusion is considered a small correction to Gaussian statistics. It is tempting to try a similar approach for knots. The idea would be to note that at small $`N<N_0`$, the probability of a non-trivial knot is small, which implies that restricting the loop such that it remains a trivial knot excludes only a small sector of the conformation space which therefore, comprises a small correction to Gaussian statistics.
Let us try to imagine the realization of this idea. We want to find the swelling ratio of the trivial loop:
$$\rho _{0_1}=R_g^2_{triv}/R_g^2_{phantom}.$$
(4)
We know that the (topology blind) ensemble average over all knots must, by definition, yield unity for the swelling ratio:
$$1=P_{0_1}\rho _{0_1}+P_{3_1}\rho _{3_1}+P_{4_1}\rho _{4_1}+\mathrm{},$$
(5)
where $`P_i`$ and $`\rho _i`$ are, respectively, the probability and swelling ratio of the knot $`i`$. Our plan is to consider formula (5) as the equation from which we can determine the quantity of interest, $`\rho _{0_1}`$:
$$\rho _{0_1}=\frac{1P_{3_1}\rho _{3_1}P_{4_1}\rho _{4_1}\mathrm{}}{P_{0_1}}.$$
(6)
To this point our consideration is exact, but now we switch to hand waving arguments and guesses justified by the simulation data. In the range of small $`N`$, the ensemble of loops consists mostly of $`0_1`$ knots, perturbed slightly by the presence of $`3_1`$ and higher-order or more complex knots. We consider then $`N/N_0`$ as a small parameter: $`N/N_01`$. Of course, in the case of excluded volume, the similar limit is better justified, because $`N^{}`$, equation (1), can at least in principle, be arbitrarily large, leaving room for the intermediate asymptotics $`1NN^{}`$. In the case of knots, $`N_0`$ is as large as about $`300`$, but so far we do not know why it is large, and it seems beyond our control to make it larger. Accordingly, we cannot speak of an intermediate asymptotics in a mathematically rigorous way Barenblatt . Nevertheless, we assume here that the numerically large value of $`N_0`$ allows us hope that the asymptotic argument is possible, and so we assume that $`N/N_0`$ is a small parameter. We guess then that higher order knots provide only higher order perturbation corrections with respect to this parameter, and we neglect their contributions, simplifying the ensemble by accounting for only $`0_1`$ and $`3_1`$ knots. In this case, $`P_{0_1}+P_{3_1}1`$. This is justified by the data presented in Figure 1, which shows that higher knots are very rare indeed. Since we know that $`P_{0_1}\mathrm{exp}(N/N_0)`$, we can also find $`P_{3_1}`$. Given that we consider the $`N/N_01`$ regime, we must also linearize the exponent, which yields:
$`\rho _{0_1}`$ $``$ $`{\displaystyle \frac{1(1P_{0_1})\rho _{3_1}}{P_{0_1}}}{\displaystyle \frac{1\left(1e^{N/N_0}\right)\rho _{3_1}}{e^{N/N_0}}}`$ (7)
$``$ $`\left(1\left(N/N_0\right)\rho _{3_1}\right)\left(1+N/N_0\right).`$
The next step requires thinking about $`\rho _{3_1}`$. In principle, we can come up with a chain of equations, not unlike the BBGKI chain in the theory of fluids, expressing $`\rho _{3_1}`$ in terms of higher knots, etc. A more practical course is to note that for the lowest order in perturbation, with respect to the supposedly small parameter $`N/N_0`$, since $`\rho _{3_1}`$ has already the small ($`N/N_0`$) coefficient in front of it, it is enough to replace $`\rho _{3_1}`$ with a constant at $`N/N_00`$. Thus, to the lowest order in $`N/N_01`$ we get $`\left(N/N_0\right)\rho _{3_1}\left(N/N_0\right)c`$, where $`c`$ is a constant. We therefore finally obtain
$$\rho _{0_1}1+\left(N/N_0\right)(1c),$$
(8)
or
$$R_g^2_{triv}\mathrm{}^2N\frac{1}{12}\left[1+\left(\frac{N}{N_0}\right)(1c)\right].$$
(9)
The difference between equations (3) and (9) is immediately obvious: the former is an expansion in powers of $`\sqrt{N}`$, the latter starts from the first power of $`N`$. The $`\sqrt{N}`$ term does not occur in our expansion for knots. Note that the values of the $`k_i^{}`$ coefficients in equation (3) are known Yamakawa , and this prevents the easy (and incorrect) explanation that $`k_1^{}=0`$. As regards the value of coefficient $`c`$, we do not have at present an analytical means to calculate it, we will later estimate it based on the simulation data. Thus, despite identical scaling index at large $`N`$, trivially knotted and excluded volume polymers exhibit a very different mathematical structure of $`N`$-dependence in their respective gyration radii in the region of small $`N`$.
It is possible that another manifestation of the same difference is the fact that data in the work swiss\_PNAS were successfully fitted to the equation (2) across the crossover region, where this formula for the self-avoiding polymers is not supposed to work.
Thus, our considerations suggest that there is some fundamental difference between topology and self-avoidance in terms of their respective effects on the swelling at moderate $`N`$. In what follows, we present computational tests supporting and further developing this conclusion.
## III Model and Simulation Methods
We model polymer loops as a set of $`N+1`$ vertices, $`\stackrel{}{x}_i`$, embedded in $`3D`$, where $`\stackrel{}{x}_0=\stackrel{}{x}_N`$ implies loop closure. The step between successive vertices, $`\stackrel{}{y}_i=\stackrel{}{x}_{i+1}\stackrel{}{x}_i`$ is constructed either from steps of fixed length, with probability density
$$P(\stackrel{}{y}_i)=\frac{1}{4\pi \mathrm{}^2}\delta \left(\left|\stackrel{}{y}_i\right|\mathrm{}\right),$$
(10)
or Gaussian distributed, with probability density
$$P(\stackrel{}{y}_i)=\left(\frac{3}{2\pi \mathrm{}^2}\right)^{3/2}\mathrm{exp}\left(\frac{3\left|\stackrel{}{y}_i\right|^2}{2\mathrm{}^2}\right).$$
(11)
Note that $`\mathrm{}`$, the โaverageโ steplength, is defined, $`\mathrm{}^2=P(y)y^2d^3y`$. Many methods have been used to generate loops in computer simulation over the past decade. A brief review of the methods is available in Appendix A, the details of the method implemented in this work are presented in Appendix B.
Once generated, we asses the loopโs size by calculating its radius of gyration
$$R_g^2=\frac{1}{2N^2}\underset{ij}{}\left|\stackrel{}{x}_i\stackrel{}{x}_j\right|^2.$$
(12)
The mean square average radius of gyration seen over all loops is, $`R_g^2=\frac{1}{12}(N+\beta )l^2`$, where $`\beta =1`$ for fixed steplength loops and $`\beta =1/N`$ for loops of gaussian distributed steplength. Noting that the excluded volume constraint is maintained by the condition that pair distances be larger than excluded volume bead diameter, $`r_{ij}=\left|\stackrel{}{x}_i\stackrel{}{x}_j\right|`$, $`r_{ij}d`$, we record the minimum $`r_{ij}`$ for each coil, which enables us to ascertain what maximum diameter of excluded volume, $`d`$, the loop corresponds to, footnote\_1 . Finally, we calculate the topological state of the loop by computing the Alexander determinant, $`\mathrm{\Delta }(1)`$, and Vassiliev knot invariants of degree 2 and 3, $`v2`$ and $`v3`$, the implementation of which is described in Lua\_Invariants . As the simulation progresses, averages are accumulated in a matrix, indexed over different knot types and minimum pair distances. In the end, we can collect the data to find the gyration radius for either a particular knot type irrespective of pair distances (i.e., without volume exclusion), or for a particular excluded volume value irrespective of topology.
## IV Results
### IV.1 On the functional form of $`N`$-dependence of the gyration radius in the moderate $`N`$ regime
Figure 2 provides direct comparison of the computationally determined mean square gyration radius for trivial knots and phantom loops with excluded volume (averaged over all topologies), in the latter case - for various values of the bead diameter. Note that in the figure, the gyration radius is expressed with the swelling ratio $`\rho `$, as defined in equation (6). The most striking feature of this figure is the differently shaped curves of swelling. The region of intermediate $`N`$ visible in the figure, $`1<N<N_0`$, shows the plot of trivial knot swelling passing through all excluded volume curves. As seen, the very shape of the $`\rho _{0_1}`$ curve is different. Specifically, all curves for the excluded volume loops are bent downwards, consistent with the presence of the $`\sqrt{N}`$ terms in equation (3). In contrast, the curve for the topologically restricted trivial loop is very nicely linear. A fit of the form
$$\rho _{0_1}=0.998+N/14371+0.18N/N_0,$$
(13)
consistent our estimate, equation (9), where $`N_0=255`$, is shown in Figure 3. Note that deviation from the linear form occurs as $`N`$ increases. This is entirely expected as the crossover to asymptotic swelling of the gyration radius, $`N^{2\nu }/NN^{0.19}`$, must occur as $`N`$ grows beyond $`N_0`$.
### IV.2 Which excluded volume diameter matches most closely the topological swelling of trivial knots?
The cross-over points between curves of trivially knotted loops and loops with excluded volume in Figure 2 inspired the idea of plotting the excluded volume diameter at each $`N`$ whose swelling matches the swelling of a trivial knot at the same $`N`$. As seen in Figure 4 this mapping parameter seems to approach an asymptote at the specific diameter of $`d=0.1625`$. While at present it is not computationally feasible to extend the scale of $`N`$ to significantly larger values, this asymptotic approach of trivial knot swelling to loops with excluded volume is consistent with the similar asymptotic swelling of $`N^{2\nu }`$ seen in other work Deguchi\_2003 ; swiss\_PNAS ; Nathan\_PNAS .
At the same time, it is interesting to note that although the swelling parameter due to the excluded volume at $`d0.16`$ seems to fit the topologically driven swelling, the corresponding characteristic length $`N^{}`$ (see (1)) is significantly larger than $`N_0`$. To see this, we note that the excluded volume data in figure 2 fit reasonably well to the expression $`\rho 1+1.71\sqrt{N}(d/\mathrm{})^3=1+\sqrt{N/N^{}}`$, where, therefore, $`N^{}=0.34(d/\mathrm{})^6`$. Here, we determined, based on the fit, the numerical coefficient intentionally left undetermined in formula (1). At $`d=0.16\mathrm{}`$, we get, therefore, $`N^{}20000`$, which is almost two orders of magnitude greater than $`N_0255`$. Alternatively this situation can be seen by finding the excluded volume diameter for which crossover length $`N^{}`$ matches $`N_0`$: $`N^{}=N_0`$; the corresponding $`d`$ equals $`d0.33\mathrm{}`$. It is fairly obvious that this value of excluded volume does not agree well with the data presented in figure 4. This discrepancy possibly points at yet another difference between swelling driven by topology and excluded volume.
## V Conclusions
It seems quite clear from our simulation data that the analogy between excluded volume and trivial knotting does not hold at loop sizes smaller than the crossover for knots, $`N_0`$. The nature of the swelling function, $`\rho (N)`$, in this region is yet unknown. Although our cursory explanation accounts for the trivial knot dataโs linear trend in this regime, the similar parameter for the size of more complex knots behaves non-linearly, and we currently have no explanation for this. A more systematic treatment of the problem is badly needed to understand the size behavior of knots.
That said, our data showing the mapping of excluded volume diameter to trivial knot size seems to reinforce the notion that asymptotically, the two classes of objects scale with the same power.
We express thanks to R. Lua of the University of Minnesota for the use of his Knot Analysis routines. We also wish to thank the Minnesota Supercomputing Institute for the use of their facilities. This work was supported in part by the MRSEC Program of the National Science Foundation under Award Number DMR-0212302.
## Appendix A A Brief Review of Loop Generation Methods
A number of methods exist and have been used in the literature for the computational generation of looped polymers. The goal of generation methods is to produce statistically representative and unbiased sets of mutually uncorrelated loops. The generation of a random walk is a simple matter. Steps are chosen with isotropic probability until the desired length is reached. Creating random walks with biased probability, specifically, walks which return to the origin after a specified number of steps, is a more difficult task. As many studies of the topological properties of polymer chains have been completed, we do not intend to make an exhaustive summary of all work, but rather in broad strokes summarize the generation methods used in the field.
All methods used to generate loops can be grouped into two large categories. Methods of one group start from some loop configuration which does not pretend to be random, and then transform it in some way to randomize the set of steps making the loop. Methods of the other group build more or less random loops from the very beginning.
One of the initial techniques used for the generation of loops is the dimerization method of Chen, Chen\_1 ; Chen\_3 , in which smaller sets of walks are joined end to end to form larger walks or loops. This โRing Dimerizationโ accepts the joining of smaller walks with some probability, as self-intersections between the chains are prohibited. In addition, if the generated walk is closed to form a loop, a statistical weight is calculated to account for loop closure. Several groups have used this method, koniaris\_muthu\_N0 ; Deguchi\_2002 , usually in the context of including excluded volume in the topological study.
Other workers, Deutsch ; Deguchi\_2003 , start with an initial loop conformation and then modify it by applying a number of โelbowโ pivot moves on randomly selected sections of the loop. Specifically, if the loop is defined by $`N`$ vertices, $`\{\stackrel{}{x}_i\}`$, a pivot move is performed by selecting two vertices, $`\stackrel{}{x}_j`$ and $`\stackrel{}{x}_k`$, and then rotating by a random angle the intermediate vertices $`\stackrel{}{x}_{j+1}`$ through $`\stackrel{}{x}_{k1}`$ about the axis made by $`\stackrel{}{x}_k\stackrel{}{x}_j`$.
A third method in common use, the so-called โhedgehogโ method Vologodskii\_hedgehog ; swiss\_PNAS , starts by generating $`N/2`$ pairs of mutually opposite bond vectors. The resulting set of $`N`$ vectors has zero sum, and it is tempting reshuffle them and then use as bond vectors, thus surely obtaining a closed loop. Unfortunately, such a loop has obviously correlated segments, the most striking manifestation of which is that the loop has self-intersections with a large probability of order unity (in fact, $`1/e0.37`$, combinatorics ; see also a related scaling argument in Nathan\_PNAS ). To overcome this, Dykhne Vologodskii\_hedgehog suggested imagining all $`N`$ vectors plotted from the origin and thus forming something like a hedgehog, and then randomly choosing pairs of vectors (hedgehog needles), and rotating the pair by a random angle about their vector sum. This operation does not change the sum of all $`N`$ vectors, which remains zero, and therefore, upon sufficiently many such operations and upon reshuffling all vectors, one can hope to obtain a well randomized loop.
The hedgehog method and elbow moves method are in fact quite similar. Indeed, in both cases the idea is to rotate some bond vectors around their vector sum; in the hedgehog method it is done with pairs of vectors before reshuffling, in the elbow moves method it is done after reshuffling with a set of subsequent bonds, but the idea is the same. In both cases, the evolution of loop shape can be described by Rouse dynamics, known in polymer physics (see, e.g., AG\_Red ). This allows us to make a simple estimate as to how many moves are necessary in order to wash away correlations imposed by the initial loop configuration. Rouse dynamics can be understood as diffusive motion of Fourier modes. Since the longest wave Fourier mode has wavelength which scales as $`N`$, the longest relaxation time in Rouse dynamics scales as $`N^2`$. This estimate is valid for physical dynamics in which all segments move at the same time. Translated into computational language, this implies that every monomer has to make about $`N^2`$ moves, which means that we have to make about $`N^3`$ random moves for proper removal of correlations. Unfortunately this point is rarely mentioned in the use of these algorithms, (see however, Deutsch ), and the number of moves between sampling is generally quite small, which puts into question the ergodicity of implementations of this algorithm.
To overcome this problem, we proposed in Nathan\_PNAS another method which we call the method of triangles, which does not involve any relaxation. In this method, we generate $`N/3`$ randomly oriented triplets of vectors with zero sum, reshuffle them, and connect them head-to-tail, thus obtaining a loop. As we shall explain in another publication, this method produces loops with insignificant correlations when $`N`$ is larger than a hundred or so.
Since our major attention in this article is the range of relatively small $`N`$, we have to resort to a computationally more intensive, but reliably unbiased method based on conditional probabilities. The idea is to generate step number $`i`$ in the loop of $`N`$ steps using the conditional probability that the given step arrives to a certain point provided that after $`Ni`$ more steps the walk will arrive at the origin. This method was suggested and implemented for Gaussian chains in volog . Here, we apply it for the loops with fixed step length.
## Appendix B Generation of loops with fixed steplength using the conditional probability method
### B.1 Derivation of the Conditional Probability Method
A walk is composed of $`N`$ steps between $`N+1`$ nodes, a step from nodes $`\stackrel{}{x}_i`$ to $`\stackrel{}{x}_{i+1}`$ having normalized probability, $`g(\stackrel{}{x}_i,\stackrel{}{x}_{i+1},1)`$. The probability for a random walk composed of $`N`$ such steps is described by the Green function which ties the steps together,
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)=g(\stackrel{}{x}_1\stackrel{}{x}_0)g(\stackrel{}{x}_2\stackrel{}{x}_1)\mathrm{}g(\stackrel{}{x}_N\stackrel{}{x}_{N1})๐\stackrel{}{x}_1๐\stackrel{}{x}_2\mathrm{}๐\stackrel{}{x}_{N1}$$
(14)
Note that in this notation the walk stretches from $`\stackrel{}{x}_0`$ to $`\stackrel{}{x}_N`$. The specifics of integration depend on the sort of steps which are being taken. At times, these integrations can be difficult to evaluate. In such cases the convolution theorem can be of some utility. Suppose that the Fourier transform and inverse is defined in the usual way,
$$\begin{array}{cc}g_\stackrel{}{k}=\beta g(\stackrel{}{x})\mathrm{exp}\left[i\stackrel{}{k}\stackrel{}{x}\right]๐\stackrel{}{x}& \\ g(\stackrel{}{x})=\beta g_\stackrel{}{k}\mathrm{exp}\left[i\stackrel{}{k}\stackrel{}{x}\right]๐\stackrel{}{k}.& \end{array}$$
(15)
Note that in this formulation $`\beta =(2\pi )^{3/2}`$. The convolution theorem allows for the following expression for $`N2`$,
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)=(1/\beta )^{N2}(g_\stackrel{}{k})^N\mathrm{exp}\left[i\stackrel{}{k}(\stackrel{}{x}_N\stackrel{}{x}_0)\right]๐\stackrel{}{k}.$$
(16)
If steplength is fixed to a certain distance, $`\mathrm{}`$, the probability distribution and its fourier transform are expressed,
$$\begin{array}{cc}g(\stackrel{}{x}_0,\stackrel{}{x}_1,1)_{fixed}=\frac{\delta (\left|\stackrel{}{x}_1\stackrel{}{x}_0\right|\mathrm{})}{4\pi l^2}& \\ & \\ g_\stackrel{}{k}=\beta \frac{\mathrm{sin}\left(k\mathrm{}\right)}{k\mathrm{}},& \end{array}$$
(17)
Using equations (16) and (17), along with differential volume $`d\stackrel{}{k}=2\pi k^2dkd(\mathrm{cos}\theta )`$, the probability distribution for a walk of $`N`$ fixed-length steps spanning the displacement $`\stackrel{}{x}_N\stackrel{}{x}_0`$ is,
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)_{fixed}=\beta ^24\pi _0^{\mathrm{}}\left(\frac{\mathrm{sin}\left[k\mathrm{}\right]}{k\mathrm{}}\right)^N\frac{\mathrm{sin}\left[k\left|\stackrel{}{x}_N\stackrel{}{x}_0\right|\right]}{k\left|\stackrel{}{x}_N\stackrel{}{x}_0\right|}k^2๐k.$$
(18)
If we use the definition of $`\beta `$ and express Sine terms as exponentials, also using $`d=\left|\stackrel{}{x}_N\stackrel{}{x}_0\right|/\mathrm{}`$ then,
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)_{fixed}=\frac{1}{2\pi ^2}_0^{\mathrm{}}\frac{\left(\mathrm{exp}\left[ik\mathrm{}\right]\mathrm{exp}\left[ik\mathrm{}\right]\right)^N\left(\mathrm{exp}\left[ik\mathrm{}d\right]\mathrm{exp}\left[ik\mathrm{}d\right]\right)}{(2ik\mathrm{})^{N+1}d}k^2๐k.$$
(19)
Then using the Newton binomial $`(x+y)^N=_{m=0}^N\left(\genfrac{}{}{0pt}{}{N}{m}\right)x^{Nm}y^m`$, where, $`\left(\genfrac{}{}{0pt}{}{N}{m}\right)=\frac{n!}{(nm)!m!},`$ yields a shiny prize, an analytically tractable expression:
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)_{fixed}=\frac{1}{\pi ^2}\frac{1}{2^{N+2}i^{N+1}\mathrm{}^{N+1}d}_0^{\mathrm{}}\underset{m=0}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{m}\right)\frac{(\mathrm{exp}\left[ik\mathrm{}\right])^{Nm}(\mathrm{exp}\left[ik\mathrm{}\right])^m(\mathrm{exp}\left[ik\mathrm{}d\right]\mathrm{exp}\left[ik\mathrm{}d\right])}{k^{N1}}dk.$$
(20)
At this point two further simplifications are made. The first is to extend the integration from $`\mathrm{}`$ to $`\mathrm{}`$, as the integrand is even on the real axis (with proper incorporation of the factor of $`1/2`$). The second simplification is to integrate over the dimensionless number, $`\kappa =k\mathrm{}`$. Note that the dimension of the integral remains $`1/volume`$.
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)_{fixed}=\frac{1}{\pi ^2}\frac{N!}{2^{N+3}i^{N+1}\mathrm{}^3d}_{\mathrm{}}^{\mathrm{}}\underset{m=0}{\overset{N}{}}\frac{(1)^m}{(Nm)!m!}\frac{\mathrm{exp}\left[i\kappa (N2m+d)\right]\mathrm{exp}\left[i\kappa (N2md)\right]}{\kappa ^{N1}}d\kappa .$$
(21)
The integral which remains can be evaluated as a contour integral in the complex plane. The contour along the real axis is chosen with a small bump in the $`+i`$ direction at $`\kappa =0`$. The upper or lower arch is chosen according to Jordanโs Lemma. The residue at $`\kappa =0`$ is obtained by Taylor expanding the exponent to resolve the coefficient corresponding to the $`\kappa ^1`$ term, which is the definition of a residue. The result follows,
$$_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{exp}\left[i\alpha \kappa \right]}{\kappa ^{N1}}๐\kappa =\{\begin{array}{ccc}0& \mathrm{if}& \alpha 0\\ 2\pi i\left(\frac{1}{(N2)!}(i\alpha )^{N2}\right)& \mathrm{if}& \alpha <0\end{array}.$$
(22)
Integration winnows the sum considerably, the final result is,
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)_{fixed}=\frac{N(N1)}{2^{N+2}\pi l^3d}\left(J_1(N,d)J_2(N,d)\right),$$
(23)
where
$$J_1(N,d)=\underset{m>(N+d)/2}{\overset{N}{}}\frac{(1)^m}{(Nm)!m!}(N2m+d)^{N2},$$
(24)
and
$$J_2(N,d)=\underset{m>(Nd)/2}{\overset{N}{}}\frac{(1)^m}{(Nm)!m!}(N2md)^{N2}.$$
(25)
A table of probabilities can then be composed. Note however that the probability is defined on intervals over $`d`$, listed in the right column below.
$$\begin{array}{cc}G(\stackrel{}{x}_0,0,3)_{fixed}=\{\begin{array}{cc}\frac{1}{8\pi \mathrm{}^3d}& d[0,2]\end{array}& \\ & \\ G(\stackrel{}{x}_0,0,3)_{fixed}=\{\begin{array}{cc}(1)/(8\pi \mathrm{}^3)& d[0,1]\\ (3d)/(16\pi d\mathrm{}^3)& d[1,3]\end{array}& \\ & \\ G(\stackrel{}{x}_0,0,4)_{fixed}=\{\begin{array}{cc}(83d)/(64\pi \mathrm{}^3)& d[0,2]\\ (d4)^2/(64\pi \mathrm{}^3d)& d[2,4]\end{array}& \\ & \\ G(\stackrel{}{x}_0,0,5)_{fixed}=\{\begin{array}{cc}(5d^2)/(64\pi \mathrm{}^3)& d[0,1]\\ (2d^315d^2+30d5)/(192\pi \mathrm{}^3d)& d[1,3]\\ (d5)^3/(384\pi \mathrm{}^3d)& d[3,5]\end{array}& \\ & \\ G(\stackrel{}{x}_0,0,6)_{fixed}=\{\begin{array}{cc}(5d^324d^2+96)/(1536\pi \mathrm{}^3)& d[0,2]\\ (5d^4+72d^3360d^2+672d240)/(3072\pi \mathrm{}^3d)& d[2,4]\\ (d6)^4/(3072\pi \mathrm{}^3d))& d[4,6]\end{array}.& \end{array}$$
(26)
These piecewise-defined probability distributions approach the shape of the corresponding quantity for gaussian distributed steplength,
$$G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)_{gaussian}=\left(\frac{3}{2\pi N\mathrm{}^2}\right)^{3/2}\mathrm{exp}\left[\frac{3}{2N\mathrm{}^2}(\stackrel{}{x}_N\stackrel{}{x}_0)^2\right].$$
(27)
Due to the complexity and computational expense of the conditional method, and noting the apparent similarity of the two curves, one might be tempted to substitute the Gaussian formulation, equation (27), when $`N`$ is above some threshold, $`N>N_c`$. Our own experience with this approximation leads us to discourage the intermingling of the two distributions. When included, at even the large $`N_c=30`$, a sharp discontinuity in the curve of curve for $`\rho _{0_1}`$ vs $`N`$ (Figure (3)) was visible at $`N_c`$. We hypothesize that substitution of the Gaussian formulation, equation (27), for the fixed-step formulation, equation (23), allows for slightly more inflated loop conformations and thus leads to a discontinuity when the approximation is used in the simulation code at $`N>N_c`$.
### B.2 Implementation of Conditional Probability Method
Generation of a random walk which is looped, i.e. $`\stackrel{}{x}_N\stackrel{}{x}_0=0,`$ can be achieved with the use of the already derived equations. Imagine that a walk of $`N+M`$ steps is underway and $`M`$ steps have already been taken. This means that a walk of $`N`$ steps remains, which starts at the present location, $`\stackrel{}{x}_0`$, and finishes at the starting point, $`\stackrel{}{x}_N`$. The probability distribution for the next step, from $`\stackrel{}{x}_0`$ to $`\stackrel{}{x}_1`$, can then be written,
$$P(\stackrel{}{x}_0|\stackrel{}{x}_1)=\frac{G(\stackrel{}{x}_0,\stackrel{}{x}_1,1)G(\stackrel{}{x}_1,\stackrel{}{x}_N,N1)}{G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)}$$
(28)
In principle one could generate new steps with probability isotropic in direction, accepting them with conditional probability defined by equations (28) and (23) or (27). In the interest of efficiency, a better method is to generate random steps within these probability distributions. Now discussed is the way to transform a flat random distribution (that produced by the UNIX math function drand48() for example) into the distribution above. If the flatly distributed variable is $`q`$, ie $`P(q)=1`$ on $`[0,1)`$, $`0`$ elsewhere, the following equation, with $`\stackrel{}{d}=\stackrel{}{x}_N\stackrel{}{x}_1`$, defines the transform to the conditional distribution above, $`G(\stackrel{}{x})`$,
$$_0^qP(q^{})d(q^{})=_0^{f(q)}\frac{G(\stackrel{}{x}_0,\stackrel{}{x}_1,1)G(\stackrel{}{d},0,N1)}{G(\stackrel{}{x}_0,\stackrel{}{x}_N,N)}d(\stackrel{}{d}),$$
(29)
In this statement of normalization, the function of importance is $`f(q)`$, which defines the way the two probability distributions are made equal.
In principle the problem is now solved. A complete set of probability distributions for walks of fixed or gaussian steplength has been defined, and the the formula which maps that distribution to a flat, machine-generated distribution has also been expressed. If the form of equation (29) is simple enough, meaning relatively small $`N`$, the integral equation can be solved directly for $`f(q)`$. In practice however, $`N>5`$ is an interesting regime and a different technique must be used to obtain $`f(q)`$.
For the case of finishing a random walk of fixed length steps, $`\mathrm{}`$, which is $`\stackrel{}{x}`$ away from the ending point, and has $`N`$ steps alloted to get to that point, we use the geometry shown in figure (5). In this diagram $`\stackrel{}{x}_p`$ is the new distance away from the endpoint after the present step is taken. Thus the expression above becomes,
$$_0^qP(q^{})๐q^{}=_0^{f(q)}\frac{G(\stackrel{}{l},1)G(\stackrel{}{x_p},N1)}{G(\stackrel{}{x},N)}d(\stackrel{}{x_p}),$$
(30)
where, for convenience, the following syntax is used, $`G(\stackrel{}{b},0,N)=G(\stackrel{}{b},N)`$.
Of course the single step $`G(\stackrel{}{d},1)`$ is a delta-function, $`\delta (|\stackrel{}{d}|\mathrm{})/4\pi \mathrm{}^2`$, so the integration over $`d(\stackrel{}{x_p})`$ occurs over most or all of the spherical shell created by the possible orientations of $`\mathrm{}`$. Integration over the shell (about the axis made by $`\stackrel{}{x}`$) is performed in โrings,โ each ring having circumference $`2\pi \mathrm{}\mathrm{sin}[\theta ]`$, and width, $`\mathrm{}d(\theta )`$, with resulting differential area, $`dA=2\pi \mathrm{}^2\mathrm{sin}[\theta ]d(\theta )`$. $`\theta `$ is integrated over the range, $`[0,\pi ]`$.
It should be apparent that, $`x_{p}^{}{}_{}{}^{2}=x^2+\mathrm{}^2+2x\mathrm{}\mathrm{cos}[\theta ]`$. This yields the differential transform, $`\mathrm{sin}[\theta ]d(\theta )=(x_p/x\mathrm{})d(x_p)`$. This simplification allows the integration of equation (30) in the following way,
$$_0^qP(q^{})๐q^{}=\frac{1}{2\mathrm{}xG(x,N)}_{min}^{f(q)}G(\stackrel{}{x_p},N1)x_pd(x_p),$$
(31)
This expression is normalized to $`1`$ if integrated over appropriate $`x_p`$ bounds. In most cases, those bounds are $`[x\mathrm{},x+\mathrm{}]`$, although the physical limit on the upper bound, $`x_p(N1)\mathrm{}`$ is necessary to keep the walk from straying too far from the origin. Additionally, if the walk is very close to the origin, $`x<\mathrm{}`$, the integration bounds, $`[\mathrm{}+x,\mathrm{}x]`$, are used.
As Equation (25) for fixed steplength probability is defined as a polynomial, integration of that polynomial, described by Equation (31), can be performed exactly within simulation computer code, and the resulting equation for $`f(q)`$ solved numerically. In practice we use the Gnu Multiple Precision library to represent the polynomial coefficients and values as rational numbers. From a computational standpoint this is significantly more expensive than representing coefficients as double floating point, but using rationals allows us to represent all outputs of the polynomial with great accuracy, the goal of this simulation method. At a relatively small number of steps the coefficients become quite small, for example at $`N=15`$, in the region $`x[13,15]`$, equation (25) reads, $`\frac{(d15)^{13}}{40809403514880(\mathrm{}^3d\pi )}`$. We feel the need in this routine to retain accuracy when performing operations such as $`PQ`$, where $`P1`$ and $`Q1`$ but $`(PQ)P,Q`$. In order to retain the accuracy of the conditional formulation it was imperative to perform this rational number algebra. For the interested reader we provide a table of these polynomial coefficients as supplementary materials. |
warning/0506/cond-mat0506570.html | ar5iv | text | # Emergence of atomic-density waves in a trapped Luther-Emery fermion gas
## Abstract
We present a novel and comprehensive microscopic study of Luther-Emery-paired phases in a strongly interacting atomic Fermi gas inside a parabolic trap and a one-dimensional ($`1D`$) optical lattice. Our work is based on a lattice version of density-functional theory, which uses as reference system the $`1D`$ homogeneous Hubbard model. We test our approach for repulsive interactions against Quantum Monte Carlo data and show that, for sufficiently strong attractions, an atomic-density wave (ADW) in the central portion of the trap breaks the discrete translational symmetry of the underlying lattice. We demonstrate that the emergence of an ADW has a dramatic impact on experimental observables such as the Fraunhofer diffraction pattern and the momentum distribution.
Strongly-interacting $`1D`$ quantum liquids are nowadays available in a variety of laboratory systems ranging from carbon nanotubes carbon\_nanotubes to semiconductor nanowires semiconducting\_nanowires , conducting molecules Nitzan\_03 , and gases in optical lattices cold\_atoms\_low\_D . Chiral Luttinger liquids at fractional quantum Hall edges chiral\_ll also provide a beautiful example of a $`1D`$ quantum liquid. The effective low-energy description of all these $`1D`$ systems is based on a harmonic theory of long-wavelength fluctuations haldane due to the interplay between topology and interactions. In the most interesting experimental situations translational invariance is broken by inhomogeneities of various types, such as Hall bar constrictions in the case of quantum Hall edges or trapping for ultracold atomic gases bec\_recent . These strong external perturbations induce new length scales causing novel physical behaviors relative to the corresponding unperturbed model system.
A powerful theoretical tool to study the interplay between interactions and inhomogeneous external fields of arbitrary shape is provided by density-functional theory (DFT) dft ; Giuliani\_and\_Vignale , which is based on the Hohenberg-Kohn theorem and on the Kohn-Sham mapping onto an auxiliary noninteracting system. Many-body effects enter DFT via the exchange-correlation (xc) functional, which is often treated by the local-density approximation (LDA) dft ; Giuliani\_and\_Vignale . The essence of LDA is to locally approximate the xc energy of the inhomogeneous system under study with that of an interacting homogeneous reference fluid, whose correlations are transferred by the LDA to the inhomogenous system. Most of the applications of DFT to inhomogeneous electronic systems use the interacting homogeneous electron gas (HEG) Giuliani\_and\_Vignale as the underlying reference fluid. The exact HEG xc energy is not known but can be calculated to a high degree of numerical precision with the help of Quantum Monte Carlo (QMC) data. However, there are a few interesting examples in the literature general in which either the reference system is not the HEG or the auxiliary system of the Kohn-Sham mapping is not an assembly of noninteracting particles. In particular, inhomogeneous $`1D`$ fermionic systems appear as an example in which it is appropriate to change the reference system to one that possesses ground-state (GS) Luttinger-liquid rather than Fermi-liquid-type correlations (see Ref. lima\_2003, and references therein to earlier work).
In this Letter we present the first such DFT study of trapped ultracold Fermi gases. We focus our attention on a two-component Fermi gas in a $`1D`$ optical lattice with either repulsive or attractive inter-component interactions and in the presence of a static external potential. Creating these types of systems experimentally seems to be within the reach of present-day technology experiments\_cold\_fermi\_gases . Studies of this model have only been carried out in the case of repulsive interactions rigol ; machida\_2004 ; xia\_ji\_2005 whereas the main focus of this Letter is on attractive interactions. We use a lattice version of DFT within an LDA based on the $`1D`$ homogeneous Hubbard model (HHM) lieb\_wu , which transfers Luttinger-Mott and Luther-Emery-type luther\_emery correlations to the inhomogeneous gas. A novel type of ground state is found at strong attractive coupling.
We consider a Fermi gas with $`N_\mathrm{f}`$ particles hopping on a $`1D`$ lattice with $`N_\mathrm{s}`$ lattice sites. Each site is labeled by the discrete coordinate $`\{z_i=i,i[1,N_\mathrm{s}]\}`$. We assume that the system can be described by a single-band Hubbard Hamiltonian jaksch\_98 ,
$`\widehat{}`$ $`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle \underset{\sigma }{}}t_{i,j}[\widehat{c}_\sigma ^{}(z_i)\widehat{c}_\sigma (z_j)+\mathrm{H}.\mathrm{c}.]`$ (1)
$`+`$ $`U{\displaystyle \underset{i}{}}\widehat{n}_{}(z_i)\widehat{n}_{}(z_i)+{\displaystyle \underset{i}{}}V_{\mathrm{ext}}(z_i)\widehat{n}(z_i).`$
Here $`t_{i,j}=t>0`$ if $`i,j`$ are nearest-neighbor sites and zero otherwise, $`\sigma =,`$ is a pseudospin-$`1/2`$ degree of freedom (hyperfine-state label), $`\widehat{n}_\sigma (z_i)=\widehat{c}_\sigma ^{}(z_i)\widehat{c}_\sigma (z_i)`$ is the pseudospin-resolved site occupation operator normalized to the number of particles with pseudospin $`\sigma `$, $`_i\widehat{n}_\sigma (z_i)=N_\sigma `$, and $`\widehat{n}(z_i)=_\sigma \widehat{n}_\sigma (z_i)`$ is the total site occupation with $`_i\widehat{n}(z_i)=N_\mathrm{f}`$. Finally $`V_{\mathrm{ext}}(z)`$ is an external potential that we take later below as representing a harmonic trap.
If $`V_{\mathrm{ext}}(z)=0`$ the Hamiltonian in Eq. (1) reduces to a $`1D`$ HHM that can be exactly solved using the Bethe-Ansatz (BA) technique for both repulsive ($`U>0`$) and attractive ($`U<0`$) interactions lieb\_wu . The properties of the system are determined by three parameters: (i) the filling factor $`n=N_\mathrm{f}/N_\mathrm{s}`$, (ii) the pseudospin polarization $`s=(N_{}N_{})/(2N_\mathrm{s})`$, and (iii) the dimensionless coupling constant $`u=U/t`$. The energy $`\epsilon (n,s;u)`$ per site can be calculated by solving a set of nonlinear coupled integral equations. In agreement with the Lieb-Mattis theorem lieb\_mattis the GS has $`s=0`$. In what follows the GS energy at $`s=0`$ will be denoted by $`\epsilon _{\mathrm{GS}}(n,u)=\epsilon (n,0;u)`$. For $`u>0`$ the $`1D`$ HHM describes a Luttinger liquid schulz\_1990 if $`n1,2`$: the two metallic GS branches for $`n<1`$ and $`n>1`$ are connected by particle-hole symmetry, $`\epsilon _{\mathrm{GS}}(n,u)=(n1)U+\epsilon _{\mathrm{GS}}(2n,u)`$. At half filling, i.e. for $`n=1`$, the GS is a Mott insulator for every $`u0`$, as signaled by a cusp in $`\epsilon _{\mathrm{GS}}`$ and by a gap in the charge excitation spectrum, while for $`n=2`$ the GS is a band insulator. Turning to the case $`u<0`$, the $`1D`$ HHM describes a Luther-Emery liquid luther\_emery , which exhibits a gap $`\mathrm{\Delta }_\mathrm{S}(n)`$ in the spin excitation spectrum and no charge gap. Due to an exact symmetry of the $`1D`$ HHM, the GS energy for $`u<0`$ is given by the energy of a repulsive $`1D`$ HHM at half filling and variable spin polarization $`s(n)=(1n)/2`$, $`\epsilon _{\mathrm{GS}}(n,|u|)=|U|n/2+\epsilon (1,s(n);|u|)`$. The Bardeen-Cooper-Schrieffer (BCS) approximation has been shown marsiglio\_1997 to provide a good description of the GS energy but a rather inaccurate and even qualitatively wrong description of the gap $`\mathrm{\Delta }_\mathrm{S}(n)`$.
In the presence of a confining external potential $`\widehat{}`$ cannot be diagonalized exactly. We calculate the GS properties of $`\widehat{}`$ in Eq. (1) by a lattice version of DFT, the so-called site-occupation-functional theory (SOFT). SOFT was introduced in a pioneering paper by Gunnarsson and Schรถnhammer soft in order to study the so-called band-gap problem Giuliani\_and\_Vignale in the context of ab initio theories of fundamental energy gaps in semiconductors and insulators. Within SOFT the exact GS site occupation $`n_{\mathrm{GS}}(z_i)=\mathrm{GS}|\widehat{n}(z_i)|\mathrm{GS}`$ can be obtained by solving self-consistently the Kohn-Sham Schrรถdinger equations
$$\underset{j}{}[t_{i,j}+v_{\mathrm{KS}}(z_i;[n_{\mathrm{GS}}(z_i)])\delta _{ij}]\phi _\alpha (z_j)=\epsilon _\alpha \phi _\alpha (z_i),$$
(2)
with $`v_{\mathrm{KS}}(z_i;[n_{\mathrm{GS}}(z_i)])=Un_{\mathrm{GS}}(z_i)/2+v_{\mathrm{xc}}(z_i)+V_{\mathrm{ext}}(z_i)`$, together with the closure
$$n_{\mathrm{GS}}(z_i)=\underset{\alpha ,\mathrm{occ}.}{}\mathrm{\Gamma }_\alpha \left|\phi _\alpha (z_i)\right|^2.$$
(3)
Here the sum runs over the occupied orbitals and the degeneracy factors $`\mathrm{\Gamma }_\alpha `$ satisfy the sum rule $`_\alpha \mathrm{\Gamma }_\alpha =N_\mathrm{f}`$. The first term in the effective Kohn-Sham potential is the Hartree potential, while $`v_{\mathrm{xc}}(z_i;[n_{\mathrm{GS}}(z_i)])=\delta E_{\mathrm{xc}}[n(z_i)]/\delta n(z_i)|_{\mathrm{GS}}`$ is the xc potential defined by the functional derivative of the xc energy $`E_{\mathrm{xc}}[n(z_i)]`$ evaluated at the GS site occupation. The total GS energy of the inhomogeneous system is given by $`E_{\mathrm{GS}}[n_{\mathrm{GS}}(z_i)]=_\alpha \mathrm{\Gamma }_\alpha \epsilon _\alpha _iv_{\mathrm{xc}}(z_i)n_{\mathrm{GS}}(z_i)_iUn_{\mathrm{GS}}^2(z_i)/4+E_{\mathrm{xc}}[n_{\mathrm{GS}}(z_i)]`$.
Equations (2) and (3) provide a formally exact scheme to calculate $`n_{\mathrm{GS}}(z_i)`$ and $`E_{\mathrm{GS}}[n_{\mathrm{GS}}(z_i)]`$, but $`E_{\mathrm{xc}}`$ and $`v_{\mathrm{xc}}`$ need to be approximated. The LDA has been shown to provide an excellent description of the GS properties of a variety of inhomogeneous systems dft . In the following we employ a BA-based LDA (BALDA) functional lima\_2003 ,
$$v_{\mathrm{xc}}^{\mathrm{BALDA}}(z_i;[n_{\mathrm{GS}}(z_i)])=v_{\mathrm{xc}}^{\mathrm{hom}}(n,u)|_{nn_{\mathrm{GS}}(z_i)}.$$
(4)
Here the xc potential of the $`1D`$ HHM is defined by $`v_{\mathrm{xc}}^{\mathrm{hom}}(n,u)=_n\left[\epsilon _{\mathrm{GS}}(n,u)k(n)\right]Un/2`$, with $`k(n)=4t\mathrm{sin}(n\pi /2)/\pi `$ being the kinetic energy of the noninteracting gas. We should point out two appealing features of Eq. (4) from the formal DFT viewpoint. First of all, the xc potential of the reference $`1D`$ HHM can be exactly calculated from BA integral equations footnote1 . Secondly, the Mott cusp in $`\epsilon _{\mathrm{GS}}(n,u)`$ for repulsive interactions is responsible for a discontinuity $`\mathrm{\Delta }_{\mathrm{xc}}(U)`$ in $`v_{\mathrm{xc}}^{\mathrm{hom}}`$ at $`n=1`$, $`\mathrm{\Delta }_{\mathrm{xc}}(U)=U2lim_{n1^{}}_n\epsilon _{\mathrm{GS}}(n<1,u)`$. As a consequence and contrary to the HEG-based LDA, the xc potential in Eq. (4) naturally possesses a derivative discontinuity Giuliani\_and\_Vignale ; lima\_2002 .
We demonstrate the accuracy of the BALDA scheme for trapped $`1D`$ Fermi gases in Eqs. (2)-(4) by showing in Fig. 1 some illustrative results for the GS site occupation of harmonically trapped repulsive fermions against the state-of-the-art QMC data of Rigol et al. rigol .
The agreement between our theory and the QMC results is clearly excellent footnote2 . We now turn to describe how the density profile changes when $`U`$ becomes negative. In Figs. 2 and 3 we report our theoretical predictions for the density profiles of attractive fermions. We clearly see from these figures that the consequences of Luther-Emery pairing in a confined Fermi gas are dramatic. The GS site occupation is a finite-size density wave, i.e. an oscillating function of $`z_i`$ that reflects the tendency of atoms with antiparallel pseudospins to form spin-singlet dimers.
In Fig. 2 we illustrate the role of the external harmonic potential.
For a very weak confinement the central part of the trap is occupied by a small-amplitude ADW, but in a sufficiently strong trap a region of doubly-occupied sites is formed at the center of the trap. In Fig. 3 we show the crossover from the weak-coupling BCS regime to a strong-coupling regime where the spin-singlet dimers are tightly bound.
On increasing the magnitude of the atom-atom attraction the amplitude of the Luther-Emery oscillations increases dramatically, giving rise to a large-amplitude ADW with period $`\lambda _{\mathrm{ADW}}2`$. The tendency of the system to form stable dimers is also confirmed numerically by a calculation of the pair-binding energy $`_\mathrm{P}=E_{\mathrm{GS}}(N_\mathrm{f}+2)+E_{\mathrm{GS}}(N_\mathrm{f})2E_{\mathrm{GS}}(N_\mathrm{f}+1)`$, which turns out to be negative (see Table 1).
We proceed to illustrate the impact of the above ADWs on some measurable quantities. The key point is that the real-space Luther-Emery ADWs with their intrinsic periodicity $`\lambda _{\mathrm{ADW}}1`$ break the lattice symmetry and transfer structure to momentum space at wavenumbers $`kk_{\mathrm{ADW}}=\pm 2\pi /\lambda _{\mathrm{ADW}}`$, away from the reciprocal lattice vectors $`G_\nu =2\pi \nu `$, $`\nu =0,\pm 1,\pm 2,\mathrm{}`$. As an illustrative example we show in Fig. 4 the elastic contribution to the light-scattering diffraction pattern vignolo\_2001 , i.e. the Fraunhofer structure factor $`S_{\mathrm{el}}(k)=N_\mathrm{f}^2\left|_i\mathrm{exp}(ikz_i)n_{\mathrm{GS}}(z_i)\right|^2`$, in the case of a relatively strong harmonic confinement \[note that $`S_{\mathrm{el}}(G_\nu )=1`$ and $`S_{\mathrm{el}}(k+G_\nu )=S_{\mathrm{el}}(k)`$\]. In this case, bulk ADWs with wavenumber $`k_{\mathrm{ADW}}=\pm \pi `$ characterize the GS density profile for a sufficiently strong attraction. The emergence of a satellite peak at $`k=\pm \pi `$ for $`u2`$ is evident in Fig. 4.
For the same parameters employed in Fig. 4 the momentum distribution function of an attractive Fermi gas possesses a satellite peak at $`k=\pm \pi `$, which could be observed in a time-of-flight experiment.
In addition to light scattering and time-of-flight-type experiments, trapping in a ring-shaped optical lattice amico\_2005 can also unambiguously reveal the Luther-Emery phases. The response of the system to an Aharonov-Bohm flux twisted\_boundary\_conditions can be used to characterize microscopically the nature of the ground state. In fact, if the GS has a spin gap as in the Luther-Emery phase, it has been shown seidel\_2005 that for a large system the period of the GS energy as a function of the Aharonov-Bohm flux is $`\pi `$, exactly half the expected one. Thus ultracold Fermi gases with attractive interactions could be a highly tunable system in which it is possible to observe direct evidence of Luther-Emery pairing.
In conclusion we remark that our DFT scheme is also applicable to a number of other experimentally interesting systems such as non-paramagnetic gases, gases under the effect of time-dependent external potentials, and $`1D`$ quantum electron liquids.
###### Acknowledgements.
This work was partially supported by MIUR, FAPESP, and CNPq. We thank Marcos Rigol for providing us with his QMC results and the โHLR-Stuttgart (Project DynMet)โ for the allocation of computer time. We also thank Dr. P. Vignolo for very useful discussions. |
warning/0506/quant-ph0506267.html | ar5iv | text | # Optimal estimation of group transformations using entanglement
## I Introduction
Many of the most surprising advantages offered by the new technology of Quantum Information Nielsen arise from the concept of quantum entanglement. Computational speed-up shor ; grover , quantum teleportation teleport and dense coding dense , secure protocols in cryptography Ekert , precision enhancement in quantum measurements DLP ; LorReview are just a short summary of some of the main lines of research inspired by entanglement.
After a so promising list, it is natural to expect remarkable improvements coming from entanglement also in the context of Quantum Estimation Theory Helstrom ; Holevo , in particular in the typical problem of estimating an unknown physical transformation drawn from a given set. With a heuristic argument inspired by dense coding, we could expect that the accuracy in the discrimination of a set of quantum channels can be increased by letting them act locally on a fixed side of a maximally entangled state. Even more, one is tempted to guess that a maximally entangled state is the optimal input for the estimation of an unknown black box. Even though these are both reasonable conjectures, in general they turn out to be false: for example, a maximally entangled input state is always uselessโand often suboptimalโfor the discrimination of two unitary transformations AcinUnitaries ; DLP . The question then arises: is it really possible to make some general statement about the role of entanglement in the optimal estimation of an unknown transformation?
In this paper we will answer this question in the *covariant case*, which corresponds to the estimation of unitary transformations randomly picked out from a given representation of some group. To face the problem, we will choose the Bayesian approach, assuming a uniform *a priori* distribution for the unknown group parameters, and defining optimality as the minimization of the average value of a given cost function. Within the Bayesian framework, some results about the optimality of maximally entangled states have been presented in AJV ; DLP . Other results in the same direction have been derived in Fujiwara ; Ballester within a different approach based on quantum Cramรฉr-Rao bound. However, all the mentioned results are limited to particular cases, and their extension to arbitrary representations of arbitrary groups is not straightforward.
Another nontrivial question is: which kind of entanglement is really useful for the estimation of group transformations? In Ref. AJV , it has been considered the estimation of unitary transformations $`U_g`$ in $`๐๐(d)`$ in the form $`U_g^N`$, corresponding to $`N`$ copies of the same unknown black box. The result is that the optimal performance can be attained by entangling the $`N`$ $`d`$-level systems that undergo the unknown transformation with another set of $`N`$ $`d`$-level systems playing the role of a *reference system*. However, as pointed out in Ref. refframe , the entanglement with an additional set of $`N`$ reference systems actually is not needed: what really matters is something more subtle, namely the entanglement between spaces where the action of the group is irreducible and spaces where the action of the group is trivial. In the language of group theory, what is needed is maximal entanglement between *representation spaces* and *multiplicity spaces*. This kind of entanglement can be obtained not only by adding an external reference system as in AJV , but also via the use of the multiple equivalent representations that appear in the Clebsch-Gordan decomposition of the representation $`\{U_g^N\}`$.
The concept of entanglement between representation spaces and multiplicity spaces will be the protagonist of this paper. In the following, we derive the optimal scheme for estimating an unknown group transformation, showing how this kind of entanglement allows to achieve the ultimate precision limits allowed by Quantum Mechanics. To do this, we introduce a class of cost functions that generalize the well known Holevo class for phase estimation Holevo , and show that all functions in such a class lead to the same optimal measurement. We give also an explicit expression for the average cost so that the optimization of the estimation scheme is reduced to a simple eigenvalue problem.
In Sec. II, before starting the analysis about optimal estimation strategies, we introduce the notation (II.1) and some group theoretical tools (II.2) that will be exploited throughout the paper. In Sec. III, we present the problem of estimating an unknown group transformation (III.1), introducing a generalization to arbitrary groups of the Holevo class of cost functions (III.2). The optimal input states are then derived (III.3), and the entanglement between representation and multiplicity spaces is recognized to be the basic resource for an optimal estimation strategy. In order to find the optimal measurement for the estimation of a group transformation, we show in Sec. III.4 how the special form of the optimal input states reflects on the covariance properties of the optimal measurement. Exploiting this analysis, we will show in III.5 that, for input states of the optimal form, all cost functions in the generalized Holevo class lead to the same optimal measurement. Finally, Sec. IV is devoted to applications of the general results. A first application (IV.1) is the optimality proof of the protocol refframe for the absolute alignment of two reference frames. As a second application, we derive (IV.2) the optimal estimation of a completely unknown two-qubit maximally entangled state with $`N`$ identical copies of the state. Section V concludes the paper, while the most technical proofs are provided in the Appendix.
## II Theoretical tools
### II.1 Notation for bipartite states
A simple notation can be introduced to deal with bipartite states. Given two Hilbert spaces $`_A`$ and $`_B`$, and fixed two orthonormal bases $`_A=\{|\varphi _n|n=1,\mathrm{},d_A\}`$ and $`_B=\{|\psi _n|n=1,\mathrm{},d_B\}`$ for $`_A`$ and $`_B`$ respectively, it is possible to associate in a one to one way any vector $`|C_A_B`$ with an operator $`C(_B,_A)`$ via the relation PLA
$$|C=\underset{m,n}{}\varphi _m|C|\psi _n|\varphi _m|\psi _n.$$
(1)
With this notation, one has the simple relations
$$C|D=\mathrm{Tr}[C^{}D]$$
(2)
and
$$AB|C=|ACB^T,$$
(3)
for any $`A(_A)`$ and $`B(_B)`$, where transposition $`T`$ is defined with respect to the fixed bases. Such relations allow to greatly simplify the calculation involving entangled states, and will be extensively used throughout the paper.
### II.2 Elements of group theory
Here we recall some simple tools of group theory groups that will be exploited throughout the paper.
Suppose we are given a Hilbert space $``$ and a unitary representation $`๐ฑ(๐)=\{U_g()|g๐\}`$ of a compact Lie group $`๐`$. The Hilbert space can be decomposed into orthogonal subspaces in the following way
$$\underset{\mu ๐ฒ}{}_\mu ^{m_\mu },$$
(4)
where the sum runs over the set of irreducible representations of $`๐`$ that appear in the Clebsch-Gordan decomposition of $`๐ฑ(๐)`$. The action of the group is irreducible in each *representation space* $`_\mu `$, while it is trivial in the *multiplicity space* $`^{m_\mu }`$, namely
$$U_g\underset{\mu ๐ฒ}{}U_g^\mu ๐_{๐_\mu },$$
(5)
$`๐_๐`$ denoting the identity in a $`d`$dimensional Hilbert space. The projection $`\mathrm{\Pi }_\mu `$ onto the subspace $`_\mu ^{m_\mu }`$ is given by the integral formula
$$\mathrm{\Pi }_\mu =d_\mu \mathrm{d}g\chi ^\mu (g)U_g,$$
(6)
where $`\mathrm{d}g`$ denotes the normalized invariant Haar measure ($`\mathrm{d}g=\mathrm{d}(kg)=\mathrm{d}(gk)`$ for any $`k,g๐`$), $`d_\mu \mathrm{dim}(_\mu )`$, and $`\chi ^\mu (g)\mathrm{Tr}[U_g^\mu ]`$ is the character of the irreducible representation $`\mu `$. Note that here we are considering $`๐`$ as a continuous group only for fixing notation, neverthelessโhere and all throughout the paperโ$`๐`$ can have a finite number of elements, say $`|๐|`$, and in this case we have simply to replace integrals with sums and $`\mathrm{d}g`$ with $`1/|๐|`$.
Moreover, any operator $`O()`$ in the commutant of $`๐ฑ(๐)`$โi.e. such that $`[O,U_g]=0g๐`$โhas the form
$$O=\underset{\mu ๐ฒ}{}๐_{๐_\mu }๐_\mu ,$$
(7)
where $`O_\mu `$ is a $`m_\mu \times m_\mu `$ complex matrix. In particular, the group average $`A_๐\mathrm{d}gU_gAU_g^{}`$ of a given operator $`A`$ with respect to the invariant Haar measure is in the commutant of $`๐ฑ(๐)`$, and has the form:
$$A_๐=\underset{\mu ๐ฒ}{}๐_{๐_\mu }\frac{\mathrm{๐}}{๐_\mu }\mathrm{Tr}__\mu [๐ธ],$$
(8)
where $`\mathrm{Tr}__\mu [A]`$ is a short notation for $`\mathrm{Tr}__\mu [\mathrm{\Pi }_\mu A\mathrm{\Pi }_\mu ]`$, $`\mathrm{\Pi }_\mu `$ being the projection onto $`_\mu ^{m_\mu }`$. Here and throughout the paper we assume the normalization of the Haar measure: $`_๐\mathrm{d}g=1`$.
Remark I: *entanglement between representation spaces and multiplicity spaces.*
The choice of an orthonormal basis $`๐ก^\mu =\{|\varphi _n^\mu ^{m_\mu }|n=1,\mathrm{},m_\mu \}`$ for a multiplicity space fixes a particular decomposition of the Hilbert space as a direct sum of irreducible subspaces:
$$_\mu ^{m_\mu }=_{n=1}^{m_\mu }_n^\mu ,$$
(9)
where $`_n^\mu _\mu |\varphi _n^\mu `$. In this picture, it is clear that $`m_\mu `$ is the number of different irreducible subspaces carrying the same representation $`\mu `$, each of them having dimension $`d_\mu `$. Moreover, with respect to the decomposition (4), any pure state $`|\mathrm{\Psi }`$ can be written as
$$|\mathrm{\Psi }=\underset{\mu ๐ฒ}{}c_\mu |\mathrm{\Psi }_\mu ,$$
(10)
where $`|\mathrm{\Psi }_\mu `$ is a bipartite state in $`_\mu ^{m_\mu }`$ and $`_{\mu ๐ฒ}|c_\mu |^2=1`$. With respect to the direct sum decomposition (9), the Schmidt number of such a state is the minimum number of subspaces carrying the same representation $`\mu ๐ฒ`$ that are needed to decompose $`|\mathrm{\Psi }`$.
Remark II: *maximum number of equivalent representations in the decomposition of a pure state.*
The Schmidt number of any bipartite state $`|\mathrm{\Psi }_\mu _\mu ^{m_\mu }`$ is always less then or equal to $`k_\mu =\mathrm{min}\{d_\mu ,m_\mu \}`$. This means that any pure state can be decomposed using *no more* than $`k_\mu `$ irreducible subspaces carrying the same representation $`\mu ๐ฒ`$.
## III Optimal estimation of group transformations
### III.1 Background problem
Suppose we are given a black box that performs on a system $`๐ฎ`$ an unknown unitary transformation $`U_g`$ randomly drawn from a group representation $`๐ฑ(๐)`$. In order to estimate the transformation $`U_g`$, we can prepare the system in an input state $`\rho ^๐ฎ`$, send it through the black box, and try to estimate the parameter $`g`$ from the output state
$$\rho _g^๐ฎU_g\rho ^๐ฎU_g^{}.$$
(11)
More generally, we can also exploit an additional reference system $``$ and prepare an entangled state $`\rho ^๐ฎ`$, so that the output state becomes
$$\rho _g^๐ฎ(U_g๐_{})\rho ^๐ฎ(๐_๐^{}๐_{}).$$
(12)
Our task is to find the best input states and the best estimation strategies allowed by Quantum Mechanics to determine the parameter $`g`$. Since we are interested in ultimate in-principle limits, we assume complete freedom in preparing any physical state and in realizing any quantum measurement. This means that we are allowed to choose the state $`\rho ^๐ฎ`$ with minimal stability group, reducing the set of unitaries that are not discriminable to those that differ just by a phase factor. Therefore, the stability group can be only a (nontrivial) center for $`๐`$, made of multiples of the identity, corresponding to (a subgroup of) $`U(1)`$. The quotient group is then a group itself, and in the following we will use the same symbol $`๐`$ for such a quotient group. Notice that the requirement of central stability group $`U(1)`$ is satisfied by choosing the state $`\rho ^๐ฎ`$ as pure, and with maximal Schmidt number.
The most general estimating strategy allowed by quantum mechanics, including both quantum measurements and classical data processing, can be described by a Positive Operator Valued Measure (POVM) $`M`$ that associates to any estimate $`\widehat{g}๐`$ a positive semidefinite operator $`M(\widehat{g})`$, satisfying the normalization condition
$$_๐\mathrm{d}gM(g)=๐.$$
(13)
The probability density of the estimate $`\widehat{g}`$ in the state $`\rho _g`$ is given by the usual Born rule:
$$p(\widehat{g}|g)=\mathrm{Tr}[\rho _gM(\widehat{g})].$$
(14)
In this paper, the estimation problem will be faced in the Bayesian setting with prior uniform probability density $`\mathrm{d}g`$, and the optimal estimation will be defined as the one that minimizes the average value of a given cost function $`c(\widehat{g},g)`$ that associates to any estimate $`\widehat{g}`$ a cost which increases versus the โdistanceโ of $`\widehat{g}`$ from the true value $`g`$. The average of the cost function over the prior and the conditional probability distributions will be given by
$$c=\mathrm{d}g\mathrm{d}\widehat{g}c(\widehat{g},g)p(\widehat{g}|g).$$
(15)
### III.2 A generalized Holevo class of cost functions
We will make two assumptions on the form of the cost function $`c(\widehat{g},g)`$.
*First assumption.* We require $`c`$ to be group invariant, namely
$$c(\widehat{g},g)=c(k\widehat{g},kg)\widehat{g},g,k๐$$
(16)
(left-invariance), and
$$c(\widehat{g},g)=c(gk,\widehat{g}k)\widehat{g},g,k๐$$
(17)
(right-invariance). By using Fourier analysis, one can prove (see Appendix) that this assumption is equivalent to the expansion
$$c(\widehat{g},g)=\underset{\sigma }{}a_\sigma \chi ^\sigma (\widehat{g}g^1),$$
(18)
where $`\chi ^\sigma (g)\mathrm{Tr}[U^\sigma (g)]`$ is the character of the irreducible representation $`\sigma `$, and the coefficients $`a_\sigma `$ satisfy the identity $`a_\sigma ^{}=a_\sigma \sigma `$, in order to have a real cost function.
*Second assumption.* We require all nonzero coefficients $`a_\sigma `$ in the expression (18) to be negative, with the only exception of the coefficient $`a_{\sigma _0}`$ corresponding to the trivial representation $`U^{\sigma _0}(g)=1g`$, which is allowed to be positive (the $`\sigma _0`$ term just adds a trivial constant to the cost function, since $`\chi ^{\sigma _0}(g)=1g`$).
The class of functions that satisfy our two assumptions is a direct generalization of the class of cost functions introduced by Holevo for the estimation of phase shifts Holevo . In fact, such functions have the form
$$c(\widehat{\varphi }\varphi )=\underset{k}{}a_ke^{ik(\widehat{\varphi }\varphi )},$$
(19)
where $`a_k0`$ for any $`k0`$, and $`e^{ik\varphi }`$ is the character of the unidimensional representation labeled by $`k`$.
### III.3 Optimal choice of the input state
Since the average cost (15) is a linear functional of the input state $`\rho `$, in the optimization problem we can restrict attention to *pure* input states $`\rho =|\mathrm{\Psi }\mathrm{\Psi }|`$. Then the problem becomes equivalent to the optimal discrimination problem of states in the orbit
$$๐ช=\left\{|\mathrm{\Psi }_gU_g|\mathrm{\Psi }|g๐\right\}$$
(20)
generated from $`|\mathrm{\Psi }`$ by the action of the representation $`๐ฑ(๐)`$.
Letโs consider the Clebsch-Gordan decomposition (5) of the unitaries $`U_g`$. From now on we will assume the algebraic condition
$$m_\mu =d_\mu \mu ๐ฒ.$$
(21)
###### Lemma 1
The assumption (21) can be done without any loss of generality.
Proof. Suppose $`d_\mu >m_\mu `$ for some representation $`\mu `$. In this case, we can introduce a reference system $``$ whose dimension is
$$d_{}\underset{\mu ๐ฒ}{\mathrm{max}}\left\{\frac{d_\mu }{m_\mu }\right\},$$
(22)
and replace $`U_g`$ with its extension $`U_g^{}=U_g๐_{}`$, acting in the tensor product Hilbert space $`_{}`$. In this way, $`U_g^{}`$ will satisfy the condition $`m_\mu ^{}m_\mu \times d_{}d_\mu \mu `$. On the other hand, as already mentioned at the end of Sec. II.2, any pure state $`|\mathrm{\Psi }`$ can be decomposed in the form (10) with no more than $`k_\mu =\mathrm{min}\{d_\mu ,m_\mu \}`$ irreducible subspaces for any $`\mu `$. Therefore, we can switch our attention from the whole Hilbert space $`_{}=_\mu _\mu ^{m_\mu ^{}}`$ to the invariant subspace $`^{}_\mu _\mu ^{d_\mu }`$, which contains the input state $`|\mathrm{\Psi }`$ along with its orbit (20). In other words, without loss of generality we can always consider an input state in the Hilbert space
$$^{}=\underset{\mu }{}_\mu ^{d_\mu },$$
(23)
which can be thought as embedded in a larger Hilbert space $`_{}`$. $`\mathrm{}`$
Remark. The need of adding an external reference system $``$ arises only in the case when $`d_\mu >m_\mu `$ for some irreducible representation $`\mu `$. In fact, the role of the reference system is simply to increase the number of equivalent representations until the extended Hilbert space $`_{}`$ reaches the threshold $`m_\mu d_\mu \mu `$. This observation allows to greatly reduce the dimension of the reference system with respect to the customary estimation schemes inspired by dense coding, with a reference system $`_{}`$ having the same dimension of $``$.
Now we show that the best input state $`|\mathrm{\Psi }`$ for estimating the group transformation of an unknown black box is a state of the form (10), with each $`|\mathrm{\Psi }_\mu `$ maximally entangled, namely
$$|\mathrm{\Psi }_\mu =\frac{1}{\sqrt{d_\mu }}\underset{n=1}{\overset{d_\mu }{}}|\psi _n^\mu |\varphi _n^\mu ,$$
(24)
$`๐ก_A^\mu =\{|\psi _n^\mu |n=1,\mathrm{},d_\mu \}`$ and $`๐ก_B^\mu =\{|\varphi _n^\mu |n=1,\mathrm{},d_\mu \}`$ being Schmidt bases for $`_\mu `$ and $`^{d_\mu }`$ respectively. Exploiting the notation (1)โwith fixed bases $`๐ก_A^\mu `$ and $`๐ก_B^\mu `$โthe optimal input state $`|\mathrm{\Psi }`$ must have the form
$$|\mathrm{\Psi }=\underset{\mu ๐ฒ}{}\frac{c_\mu }{\sqrt{d_\mu }}|W_\mu ,$$
(25)
with $`W_\mu _n|\psi _n^\mu \varphi _n^\mu |`$ unitary operators.
###### Theorem 1 (optimal input states)
With a suitable choice of the coefficients $`\{c_\mu \}`$, any input state of the form (25) achieves the minimum average cost.
Suppose that the minimum cost $`c^{Opt}`$ is achieved by the input state $`|\mathrm{\Phi }=_\mu c_\mu |\mathrm{\Phi }_\mu `$ along with the estimation strategy described by the POVM $`M(g)`$. The operator $`K_h_\mu ๐_\mu \sqrt{๐_\mu }\left(๐_\mu ^{}๐_๐^\mu \mathbb{\Phi }_\mu \right)^๐`$ converts the orbit of an input state (25) into the orbit of the optimal input state $`|\mathrm{\Phi }`$, since using identity (3), we have
$$K_h|\mathrm{\Psi }_g=|\mathrm{\Phi }_{gh},$$
(26)
where $`|\mathrm{\Psi }_g=U_g|\mathrm{\Psi }`$ and $`|\mathrm{\Phi }_g=U_g|\mathrm{\Phi }`$. Consider now the POVM $`M^{}(g)\mathrm{d}hK_h^{}M(gh)K_h`$. The POVM $`M^{}(g)`$ is normalized, since
$`{\displaystyle \mathrm{d}gM^{}(g)}`$ $`=`$ $`{\displaystyle \mathrm{d}g\mathrm{d}hK_h^{}M(gh)K_h}`$
$`=`$ $`{\displaystyle \mathrm{d}hK_h^{}K_h}`$
$`=`$ $`๐,`$
where we exchanged integrals over $`g`$ and $`h`$, used invariance of the Haar measure $`\mathrm{d}g`$, and finally used Eq. (8) and the normalization of bipartite states $`|\mathrm{\Phi }_\mu `$ in the form $`\mathrm{Tr}[\mathrm{\Phi }_\mu ^{}\mathrm{\Phi }_\mu ]=1`$. A state $`|\mathrm{\Psi }`$ of the form (25) along with the POVM $`M^{}(g)`$ achieves the minimum cost. In fact,
$`c`$ $`=`$ $`{\displaystyle \mathrm{d}g\mathrm{d}\widehat{g}c(\widehat{g},g)\mathrm{\Psi }_g|M^{}(\widehat{g})|\mathrm{\Psi }_g}`$
$`=`$ $`{\displaystyle \mathrm{d}g\mathrm{d}\widehat{g}\mathrm{d}hc(\widehat{g},g)\mathrm{\Phi }_{gh}|M(\widehat{g}h)|\mathrm{\Phi }_{gh}}`$
$`=`$ $`{\displaystyle \mathrm{d}g\mathrm{d}\widehat{g}\mathrm{d}hc(\widehat{g}h,gh)\mathrm{\Phi }_{gh}|M(\widehat{g}h)|\mathrm{\Phi }_{gh}}`$
$`=`$ $`{\displaystyle \mathrm{d}k\mathrm{d}\widehat{k}c(\widehat{k},k)\mathrm{\Phi }_k|M(\widehat{k})|\mathrm{\Phi }_k}`$
$`=`$ $`c^{Opt},`$
where we used right-invariance of both cost function and Haar measure. $`\mathrm{}`$
### III.4 Covariance properties of the estimating POVM
Since the whole orbit (20) is generated from the input state $`|\mathrm{\Psi }`$ by the action $`๐ฑ(๐)`$ of the group, there is no loss of generality in restricting attention to estimating POVM of the covariant form Holevo
$$M(g)=U_g\mathrm{\Xi }U_g^{}$$
(27)
with $`\mathrm{\Xi }`$ a suitable positive operator satisfying the normalization condition (13). A covariant POVM yields a left-invariant probability distribution, namely $`p(k\widehat{g}|kg)=p(\widehat{g}|g)k,\widehat{g},g๐`$. Using both the left-invariance of the probability distribution and of the cost function, the average cost (15) can be written as
$$c=\mathrm{d}gc(g,e)p(g|e)$$
(28)
where $`e`$ is the identity element of the group $`๐`$.
For superpositions of maximally entangled states as in Eq. (25), the orbit $`๐ช`$ enjoys an additional symmetry that reflects on an additional covariance property of the POVM. In fact, using the decomposition (5) and the identity (3), we can note that
$`|\mathrm{\Psi }_g`$ $`=`$ $`U_g|\mathrm{\Psi }`$
$`=`$ $`{\displaystyle \underset{\mu ๐ฒ}{}}{\displaystyle \frac{c_\mu }{\sqrt{d_\mu }}}(U_g^\mu ๐_\mu )|๐_\mu `$
$`=`$ $`{\displaystyle \underset{\mu ๐ฒ}{}}{\displaystyle \frac{c_\mu }{\sqrt{d_\mu }}}[๐_\mu (๐_\mu ^{}๐_๐^\mu ๐_\mu )^๐]|๐_\mu `$
$`=`$ $`V_g^{}|\mathrm{\Psi }g๐,`$
where
$$V_g_{\mu ๐ฒ}(๐_\mu (๐_\mu ^{}๐_๐^\mu ๐_\mu )^{}),$$
(29)
is an element of a new unitary representation $`๐ฑ^{}(๐)`$ of the group $`๐`$. Notice that the two representations $`๐ฑ(๐)`$ and $`๐ฑ^{}(๐)`$ commute among themselves. Then, the following Lemma holds:
###### Lemma 2
There is no loss of generality in assuming a covariant POVM $`M(g)=U_g\mathrm{\Xi }U_g^{}`$ with
$$[\mathrm{\Xi },U_gV_g]=0g๐,$$
(30)
where $`U_g`$ and $`V_g`$ are given in Eqs. (5) and (29), respectively.
Proof. For any possible POVM $`N(g)`$ there is a covariant POVM with the above property and with the same average cost. In fact, the group average
$$M(g)=\mathrm{d}k\mathrm{d}hU_k^{}V_h^{}N(kgh^1)V_hU_k$$
(31)
is covariantโnamely $`M(g)=U_g\mathrm{\Xi }U_g^{}`$ with $`\mathrm{\Xi }=M(e)`$โand satisfies the required commutation relation (30). Both properties follow simply from the invariance of the Haar measure. To prove that the cost of the covariant POVM $`M(g)`$ is the same as the cost of $`N(g)`$ we use the property
$$U_kV_h|\mathrm{\Psi }_g=|\mathrm{\Psi }_{kgh^1}k,h,g๐$$
(32)
of the states generated from the input (25). In this way,
$`c_M`$ $``$ $`{\displaystyle \mathrm{d}g\mathrm{d}\widehat{g}c(\widehat{g},g)\mathrm{\Psi }_g|M(\widehat{g})|\mathrm{\Psi }_g}`$
$`=`$ $`{\displaystyle }\mathrm{d}g{\displaystyle }\mathrm{d}\widehat{g}{\displaystyle }\mathrm{d}k{\displaystyle }\mathrm{d}hc(\widehat{g},g)\times `$
$`\times \mathrm{\Psi }_{kgh^1}|N(k\widehat{g}h^1)|\mathrm{\Psi }_{kgh^1}`$
$`=`$ $`{\displaystyle }\mathrm{d}g{\displaystyle }\mathrm{d}\widehat{g}{\displaystyle }\mathrm{d}k{\displaystyle }\mathrm{d}hc(k\widehat{g}h^1,kgh^1)\times `$
$`\times \mathrm{\Psi }_{kgh^1}|N(k\widehat{g}h^1)|\mathrm{\Psi }_{kgh^1}`$
$`=`$ $`{\displaystyle \mathrm{d}r\mathrm{d}\widehat{r}c(\widehat{r},r)\mathrm{\Psi }_r|N(\widehat{r})|\mathrm{\Psi }_r}`$
$``$ $`c_N,`$
where we used the left- and right-invariance of the cost function $`c(\widehat{g},g)`$. $`\mathrm{}`$
Letโs diagonalize the operator $`\mathrm{\Xi }`$ and express its (non-normalized) eigenvectors in the decomposition (4):
$`\mathrm{\Xi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{r}{}}}|\eta ^i\eta ^i|`$ (33)
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{\mu ,\nu }{}}\sqrt{d_\mu d_\nu }|\eta _\mu ^i\eta _\nu ^i|,`$
where the factor $`\sqrt{d_\mu }`$ has been inserted just for later convenience.
###### Lemma 3
Any covariant POVM $`M(g)=U_g\mathrm{\Xi }U_g^{}`$ with the commutation property (30) must satisfy the two relations:
$$\underset{i}{}\eta _\mu ^i\eta _\mu ^i=๐_\mu \mu ๐ฒ,$$
(34)
and
$$\underset{i}{}\eta _\mu ^i\eta _\mu ^i=๐_\mu \mu ๐ฒ.$$
(35)
Proof. The normalization (13) becomes
$$\mathrm{\Xi }_๐\mathrm{d}gU_g\mathrm{\Xi }U_g^{}=๐.$$
(36)
The group average $`\mathrm{\Xi }_๐`$ can be expressed using Eq. (8). In this way, Eq. (36) becomes
$$\frac{1}{d_\mu }\mathrm{Tr}__\mu [\mathrm{\Xi }]=๐_\mu \mu ๐ฒ.$$
(37)
By explicit computation,
$`{\displaystyle \frac{1}{d_\mu }}\mathrm{Tr}__\mu \left[\mathrm{\Xi }\right]`$ $`=`$ $`{\displaystyle \underset{i}{}}\mathrm{Tr}__\mu [|\eta _\mu ^i\eta _\mu ^i|]`$
$`=`$ $`{\displaystyle \underset{i}{}}\eta _\mu ^{iT}\mathrm{Tr}__\mu [|๐_\mu ๐_\mu |]\eta _\mu ^i`$
$`=`$ $`{\displaystyle \underset{i}{}}\eta _\mu ^{iT}\eta _\mu ^i.`$
Substituting this expression in (37) and taking the complex conjugate we get (34). Moreover, using the commutation relation (30), we can transform the group average with respect to $`๐ฑ(๐)`$ in a group average with respect to $`๐ฑ^{}(๐)`$, namely
$`\mathrm{\Xi }_๐`$ $`=`$ $`{\displaystyle \mathrm{d}gU_g(U_g^{}V_g^{}\mathrm{\Xi }U_gV_g)U_g^{}}`$
$`=`$ $`{\displaystyle \mathrm{d}gV_g^{}\mathrm{\Xi }V_g}.`$
In this way, using Eq. (29), Eq. (35) can be proved along the same lines used to prove Eq. (34). $`\mathrm{}`$
### III.5 The optimal POVM
We are now able to find the optimal covariant POVM for the estimation of group transformation with superpositions of maximally entangled states.
###### Theorem 2 (optimal POVM)
In the estimation of the states in the orbit $`๐ช`$ generated from the input state
$$|\mathrm{\Psi }=\underset{\mu ๐ฒ}{}\frac{c_\mu }{\sqrt{d_\mu }}|W_\mu ,$$
(38)
where $`W_\mu `$ are unitary operators, the covariant POVM given by $`\mathrm{\Xi }=|\eta \eta |`$ with
$$|\eta =\underset{\mu ๐ฒ}{}\sqrt{d_\mu }e^{i\mathrm{arg}(c_\mu )}|W_\mu $$
(39)
is optimal for any cost function $`c(\widehat{g},g)`$ of the form
$$c(\widehat{g},g)=\underset{\sigma }{}a_\sigma \chi ^\sigma (\widehat{g}g^1),$$
(40)
with $`a_\sigma 0\sigma \sigma _0`$.
The average cost corresponding to the optimal estimation strategy is
$$c^{Opt}=a_{\sigma _0}+\underset{\mu ,\nu }{}|c_\mu |C_{\mu \nu }|c_\nu |,$$
(41)
where
$$C_{\mu \nu }\underset{\sigma \sigma _0}{}a_\sigma m_\sigma ^{(\mu \nu )},$$
(42)
$`m_\sigma ^{(\mu \nu )}`$ being the multiplicity of the irreducible representation $`\sigma `$ in the Clebsch-Gordan series of the tensor product $`U_g^\mu U_g^\nu `$.
Proof. We will show that Eq. (41) gives a lower bound for the average cost, and that the POVM $`\mathrm{\Xi }=|\eta \eta |`$ with $`|\eta `$ given by Eq. (39) achieves this bound. By using identities (2) and (3), and the form (33) for the operator $`\mathrm{\Xi }`$, Eq. (28) becomes
$`c`$ $`=`$ $`{\displaystyle }\mathrm{d}gc(g,e)\times `$
$`\times {\displaystyle \underset{i}{}}{\displaystyle \underset{\mu ,\nu }{}}c_\mu ^{}c_\nu \mathrm{Tr}[W_\mu ^{}U_g^\mu \eta _\mu ^iW_\nu ^TU_g^\nu \eta _\nu ^i].`$
Letโs expand $`c(g,e)`$ as in (40). Subtracting from the average cost $`c`$ the constant term $`a_{\sigma _0}`$, which is not relevant for the optimization, we get
$`ca_{\sigma _0}`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{\mu ,\nu }{}}c_\mu ^{}c_\nu \times `$
$`\times {\displaystyle \underset{\sigma \sigma _0}{}}{\displaystyle \frac{a_\sigma }{d_\sigma }}\mathrm{Tr}\left[\mathrm{\Pi }_\sigma ^{(\mu \nu )}(\eta _\mu ^iW_\mu ^{}\eta _\nu ^iW_\nu ^T)\right],`$
where we defined
$$\mathrm{\Pi }_\sigma ^{(\mu \nu )}d_\sigma \mathrm{d}g\chi ^\sigma (g)U_g^\mu U_g^\nu .$$
(43)
According to (6), $`\mathrm{\Pi }_\sigma ^{(\mu \nu )}`$ is the projection onto the direct sum of all the subspaces of $`_\mu _\nu `$ that carry the irreducible representation $`\sigma `$ in the tensor product $`U_g^\mu U_g^\nu `$. Clearly $`\mathrm{\Pi }_\sigma ^{(\mu \nu )}`$ is nonzero if and only if the Clebsch-Gordan series of $`U_g^\mu U_g^\nu ^{}`$ contains $`\sigma `$ with nonzero multiplicity $`m_\sigma ^{(\mu \nu )}`$. Notice also that $`\mathrm{Tr}[\mathrm{\Pi }_\sigma ^{(\mu \nu )}]=d_\sigma m_\sigma ^{(\mu \nu )}`$, by definition of $`\mathrm{\Pi }_\sigma ^{(\mu \nu )}`$.
Denoting by $`_{\mu ,\nu ,\sigma }^{}`$ the sum over $`\mu ,\nu `$ and all $`\sigma `$ except $`\sigma _0`$, the average cost can be bounded as follows
$$\begin{array}{cc}\hfill ca_{\sigma _0}& \underset{\mu ,\nu ,\sigma }{}^{}\frac{a_\sigma }{d_\sigma }\left|c_\mu c_\nu \underset{i}{}\mathrm{Tr}\left[\mathrm{\Pi }_\sigma ^{(\mu \nu )}(\eta _\mu ^iW_\mu ^{}\eta _\nu ^iW_\nu ^T)\right]\right|\hfill \\ & \underset{\mu ,\nu ,\sigma }{}^{}\frac{a_\sigma }{d_\sigma }|c_\mu c_\nu |\sqrt{\left(\underset{i}{}\mathrm{Tr}\left[\mathrm{\Pi }_\sigma ^{(\mu \nu )}(\eta _\mu ^i\eta _\mu ^i๐_\nu )\right]\right)}\hfill \\ & \times \sqrt{\left(\underset{j}{}\mathrm{Tr}\left[\mathrm{\Pi }_\sigma ^{(\mu \nu )}(๐_\mu ๐_\mu ^{}\eta _\nu ^{๐๐}\eta _\nu ^๐๐_\mu ^๐)\right]\right)},\hfill \end{array}$$
since all $`a_\sigma `$ are nonpositive. The second inequality follows from Cauchy-Schwartz inequality with respect to the scalar product $`๐,๐_i\mathrm{Tr}\left[A_i^{}B_i\right]`$, where we take $`A_i^{}=\mathrm{\Pi }_\sigma ^{(\mu \nu )}(\eta _\mu ^iW_\mu ^{}๐_\nu )`$ and $`B_i=(๐_\mu \eta _\nu ^๐๐_\mu ^๐)\mathbb{\Pi }_\sigma ^{(\mu \nu )}`$. Exploiting the relations (34) and (35), and using that $`\mathrm{Tr}\left[\mathrm{\Pi }_\sigma ^{(\mu \nu )}\right]=d_\sigma m_\sigma ^{(\mu \nu )}`$, we obtain the bound
$`c`$ $``$ $`a_{\sigma _0}+{\displaystyle \underset{\mu ,\nu ,\sigma }{}^{}}a_\sigma m_\sigma ^{(\mu \nu )}|c_\mu c_\nu |`$ (44)
$``$ $`c^{Opt}.`$
It is straightforward to see that the choice of a covariant POVM with $`\mathrm{\Xi }=|\eta \eta |`$ with $`|\eta `$ given by (39) achieves this lower bound. $`\mathrm{}`$
### III.6 Remarks
Remark I Up to the constant term $`a_{\sigma _0}`$, the minimum cost (41) is simply given by the expectation value of the *cost matrix* (42) over the normalized vector $`๐ฏ(|c_\mu |)`$. Therefore the optimal input state is obtained just by finding the eigenvector corresponding to the minimum eigenvalue of the cost matrix. In other words, the optimal state for the estimation of an unknown parameter is always a superposition of maximally entangled states, with the coefficients in the superposition modulated by the particular choice of the cost function. Notice the simplification of the optimization problem provided by Theorem 2: instead of optimizing a state in the Hilbert space $`=_{\mu ๐ฒ}_\mu ^{m_\mu }`$ we need only to optimize a vector in $`^{|๐ฒ|}`$, where $`|๐ฒ|`$ is the number of irreducible representations contained in the action of the black box.
Remark II The optimal POVM of Theorem 2 is the same optimal POVM arising from the maximum likelihood criterion MLPovms ; DeGiorgi . In fact, this criterion corresponds to the particular choice of the delta cost function
$`c(\widehat{g},g)`$ $`=`$ $`\delta (\widehat{g},g)`$
$`=`$ $`{\displaystyle \underset{\sigma }{}}d_\sigma \chi ^\sigma (\widehat{g}g^1),`$
which is of the form (40). In other words, in the case of superpositions of maximally entangled states, the result of Theorem 2 can be viewed as the extension of the maximum likelihood approach of Ref. MLPovms to arbitrary cost functions.
Remark III In the optimization of covariant POVMโs it is often assumed that the operator $`\mathrm{\Xi }`$ corresponding to an optimal estimation can be taken with unit rank. However, for mixed states some counterexamples are known ExtPovms ; PhaseMixedStates , and for pure states there is no general proof that the POVM minimizing the average Bayes cost can be chosen with rank one. Therefore, it is important to emphasize that here the rank-one property of the optimal POVM of Theorem 2 is a result of the derivation, not an assumption.
## IV Applications
### IV.1 Optimal transmission of reference frames
The result of Theorem 2 can be exploited to give the definitive proof of optimality of the protocol for the absolute transmission of a Cartesian reference frame of Ref. refframe , which concludes a long debate about the optimal way of communicating a reference frame frames . Such a protocol allows two distant parties, Alice and Bob, to align their Cartesian axes in an absolute way, i.e. without the need of any kind of prior information about their relative orientations. To this purpose, Alice sends to Bob $`N`$ spin-$`\frac{1}{2}`$ particles, prepared in some fixed state. The preparation procedure of the state is related to the directions of Aliceโs Cartesian axes: for example Alice can align the angular momenta of some particles with her $`x`$axis, some with her $`y`$axis, and so on. When Bob receives the particles, since his axes are mismatched with Aliceโs ones, each particle appears rotated by the same unknown rotation. Then, instead of receiving the particles in the same state prepared by Alice, Bob receives them in a rotated state. Clearly, if he knows how the state should look in absence of rotations, he can try to estimate the difference, i.e. he can estimate the unknown rotation, inferring in this way the directions of Aliceโs axes. The precision of this scheme is defined in a Bayesian way, by taking as cost function the *transmission error*, i.e. the distance between the directions of Aliceโs axes and Bobโs axes at the end of the protocol. In terms of the estimated rotation $`\widehat{g}`$ and the true one $`g`$, the transmission error can be written as refframe
$$e(\widehat{g},g)=62\chi ^1(\widehat{g}g^1),$$
(45)
where $`\chi ^1(g)\mathrm{Tr}[U_g^1]`$ is the character of the three dimensional irreducible representation of the rotation group. It is immediate to see that the transmission error is a cost function the form (40).
What is the best precision that can be achieved with the mentioned protocol? To answer this question we need to solve two problems: the first is to find what is the optimal state for encoding rotations, and the second is to find the optimal estimation strategy. It is important to stress that, since we want to achieve an *absolute* transmission, we are not allowed to use an external reference system, whose role would correspond to a partially shared reference frame refframe . For this reason we are allowed only to exploit the entanglement coming from the multiple equivalent representations that appear in the Clebsch-Gordan series of $`U_g^N`$, where $`U_g`$ is the $`๐๐(2)`$ matrix that represent the rotation $`g`$ in the two-dimensional Hilbert space $``$ of a single spin $`\frac{1}{2}`$ particle.
The tensor product Hilbert space $`^N`$ can be decomposed as in (4)
$$^N=\underset{j=0(\frac{1}{2})}{\overset{\frac{N}{2}}{}}_j^{m_j}.$$
(46)
The irreducible representations are labeled by the quantum number $`j`$ of the total angular momentum, which ranges from $`0(\frac{1}{2})`$ to $`\frac{N}{2}`$ for $`N`$ being even (odd), respectively. The dimension of the representation space $`_j`$ is
$$d_j=2j+1,$$
(47)
while the multiplicities are given by CEM ; refframe :
$$m_j=\frac{2j+1}{\frac{N}{2}+j+1}\left(\genfrac{}{}{0pt}{}{N}{\frac{N}{2}+j}\right).$$
(48)
Since $`m_jd_j`$ for any $`j<\frac{N}{2}`$, it is possible to have maximal entanglement between representation spaces and multiplicity spaces for any $`j`$, with the only exception of $`j=\frac{N}{2}`$.
However, as shown in refframe , the contribution of the subspace with $`j=\frac{N}{2}`$ is negligible in the asymptotic limit of large $`N`$. Therefore we can restrict ourself to the subspace $`^{}=_{j=0(\frac{1}{2})}^{\frac{N}{2}1}_j^{m_j}`$, and consider the state
$$|A=\underset{j=0(\frac{1}{2})}{\overset{\frac{N}{2}1}{}}\frac{c_j}{\sqrt{d_j}}|๐_๐.$$
(49)
According to Theorem 1, and to the result of AJV , this is the optimal state in the subspace $`^{}`$ for the estimation of an unknown $`๐๐(2)`$ rotation.
Now we can use Theorem 2 to state that the optimal estimation strategy is described by the covariant POVM given by $`\mathrm{\Xi }=|\eta \eta |`$ with
$$|\eta =\underset{j=0(\frac{1}{2})}{\overset{\frac{N}{2}1}{}}\sqrt{d_j}|๐_๐.$$
(50)
The optimization of the coefficients $`c_j`$ in the state (49) has been done in refframe , where the POVM (50) was assumed by exploiting for simplicity the maximum likelihood approach. In this way, the results of Theorem 1 and 2 provide the optimality proof for the protocol proposed in Ref. refframe . Therefore, we can definitely state that the asymptotic precision
$$e=\frac{8\pi ^2}{N^2}$$
(51)
is the best that can be achieved for all input states and all POVMโs, namely it is the ultimate precision limit imposed by Quantum Mechanics in the absolute alignment of two Cartesian reference frames.
### IV.2 Optimal estimation of a completely unknown maximally entangled state
Maximally entangled states are a fundamental resource for quantum teleportation teleport and for quantum cryptography Ekert . To achieve ideal teleportation, Alice and Bob must know with precision which maximally entangled state they are sharing, otherwise the fidelity of the state received by Bob with the original state from Alice can be lowered. Similar arguments apply to the cryptographic schemes where the correlations arising from entanglement are exploited to generate a secret key.
Now we will consider the problem of estimating in the best way a completely unknown maximally entangled state, provided that $`N`$ identical copies are available. This can be done as an application of Theorem 2. Letโs consider a state $`|\psi `$, with $`\mathrm{dim}()=d`$. In terms of the notation (1), this state is maximally entangled if and only if $`\psi =\frac{1}{\sqrt{d}}U`$, where $`U`$ is some unitary operator. Using property (3), any maximally entangled state can be written as
$$|\psi _g=\frac{1}{\sqrt{d}}(U_g๐)|๐,$$
(52)
where $`U_g`$ is an element of the group $`๐๐(d)`$.
If $`N`$ identical copies of the unknown state $`|\psi _g`$ are given, then the problem becomes to find the best estimate for parameter $`g`$ that labels the states of the form $`|\mathrm{\Psi }_g=|\psi _g^N`$. Optimality is defined here in terms of maximization of the Uhlmann fidelity between the true state and the estimated one:
$$f(\widehat{g},g)=|\psi _g|\psi _{\widehat{g}}|^2.$$
(53)
Using the definition (52) and the property (2), we obtain
$$f(\widehat{g},g)=\frac{1}{d^2}|\chi (\widehat{g}g^1)|^2.$$
(54)
where $`\chi (g)=\mathrm{Tr}[U_g]`$. The maximization of the fidelity corresponds to the minimization of the cost function
$$c(\widehat{g},g)=1f(\widehat{g},g),$$
(55)
which is of the form (40). In particular, for $`d=2`$, $`|\chi (g)|^2=1+\chi ^1(g)`$, where $`\chi ^1(g)=\mathrm{Tr}[U_g^1]`$ is the character of the irreducible representation of $`๐๐(2)`$ with angular momentum $`j=1`$,whence we have
$$c(\widehat{g},g)=\frac{1}{4}\left(3\chi ^1(g^1\widehat{g})\right).$$
(56)
All the states of the form $`|\mathrm{\Psi }_g=|\psi _g^N`$ are generated from the input state
$$|\mathrm{\Psi }=\frac{1}{\sqrt{d^N}}|๐^{}$$
(57)
by the action of the representation
$$๐ฑ(๐)=\{(U_g๐)^{}|๐_๐๐๐(๐)\}.$$
(58)
Now we need to know how the input state is decomposed with respect to the invariant subspaces of this representation.
###### Lemma 4
Using suitable bases for the muliplicity spaces in decomposition (4), the input state (57) can be written as
$$|\mathrm{\Psi }=\underset{\mu ๐ฒ}{}\frac{c_\mu }{\sqrt{d_\mu }}|๐_\mu ,$$
(59)
where the sum runs over the irreducible representations of $`๐๐(d)`$ occurring in the Clebsch-Gordan series of $`๐ฑ(๐)`$ (58), and
$$c_\mu =\sqrt{\frac{d_\mu m_\mu }{d^N}},$$
(60)
$`d_\mu `$ and $`m_\mu `$ being respectively the dimension and the multiplicity of the representation $`\mu `$ in the Clebsch-Gordan series of $`\{U_g^N\}`$.
Proof. See appendix.
Thank to this lemma we can exploit directly the result of Theorem 2 to calculate the average fidelity. Now we will carry on the calculation of the optimal fidelity in the simplest case $`d=2`$. As usual, the irreducible representations of $`๐๐(2)`$ are labeled by the quantum number $`j`$, ranging from $`0(\frac{1}{2})`$ to $`\frac{N}{2}`$ for $`N`$ being even (odd), respectively. The minimum cost can be evaluated using Theorem 2 as
$$c^{Opt}=\frac{3}{4}+\underset{i,j=0(\frac{1}{2})}{\overset{\frac{N}{2}}{}}|c_i|C_{ij}|c_j|$$
(61)
Using Eq. (60) with the values of dimensions and multiplicities given by (47) and (48), the coefficients of the state become
$$c_i=g(i)\sqrt{\frac{1}{2^N}\left(\genfrac{}{}{0pt}{}{N}{\frac{N}{2}+i}\right)},$$
(62)
where
$$g(i)=\frac{2i+1}{\sqrt{\frac{N}{2}+i+1}}.$$
(63)
On the other hand, the matrix $`C_{ij}`$ is calculated according to the definition (42), namely by evaluating the multiplicity of the representation with angular momentum $`k=1`$ in the Clebsch-Gordan series of the tensor product $`U_g^iU_g^j`$ . In this way we get
$`C_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(\delta _{i,j}+\delta _{i,j+1}+\delta _{i,j1}).`$ (64)
Since $`_i|c_i|^2=1`$, we have
$$c=\frac{1}{2}\left(1\underset{j=0(\frac{1}{2})}{\overset{\frac{N}{2}1}{}}c_jc_{j+1}\right).$$
(65)
To obtain the asymptotic behavior of the optimal fidelity, we can approximate the binomial distribution in (62) with a Gaussian $`G_\sigma (x)`$ with mean $`\overline{x}=0`$ and variance $`\sigma ^2=\frac{N}{4}`$. Since the sum in (65) runs over a large interval with respect to $`\sigma `$, we can also approximate it with an integral over $`[0,+\mathrm{}]`$ . All these approximations hold up to order higher than $`\frac{1}{N}`$. Thus the evaluation of the optimal fidelity is reduced to the evaluation of the integral
$$I=_0^{\mathrm{}}\mathrm{d}xg(x)g(x+1)G_\sigma (x),$$
(66)
whose leading order can be obtained from Taylor expansion. In this way, we derive the asymptotic cost
$$c^{Opt}=\frac{3}{4N},$$
(67)
corresponding to the optimal fidelity
$$f^{Opt}=1\frac{3}{4N}.$$
(68)
Remarkably, the Bayes cost with uniform a priori distribution has the same asymptotic behavior of the cost of the optimal locally unbiased estimator obtained in Ballester , for any possible value $`g`$ of the true parameter. This means that in the present unbiased case the covariant measurement of Theorem 2 is optimal not only on average but also pointwise.
## V Conclusions
In this paper we solved the general problem of optimal estimation of group transformations in the Bayesian framework with uniform prior. For this purpose, we introduced a class of cost functions generalizing the Holevo class for phase estimation, containing the maximum likelihood strategy as a special case. For this family of cost functions, we derived the general form of the optimal input states, which involves maximal entanglement between representation and multiplicity spaces of the group action. More precisely, the form of an optimal input state is a direct sum of maximally entangled states, and for a given cost function one only needs to optimize the coefficients in the sum. Moreover, *for any state of the optimal form all invariant cost functions lead to the same optimal POVM*. In this way, it is possible to derive an explicit expression for the average cost and to reduce the optimization of the state to a simple eigenvalue problem. As applications of the general result we have given the first complete derivation of the ultimate precision limit imposed by Quantum Mechanics in the absolute alignment of two Cartesian reference frames, and we have derived the optimal estimation of a completely unknown two-qubit maximally entangled state with $`N`$ copies of the state. In the present paper we focused attention to compact groups and finite dimensional Hilbert spaces, nevertheless an extension of our results to infinite dimensional Hilbert spaces and non-compact groups is possible, in the same way as in DeGiorgi . However, since in infinite dimension the optimal states may be non-normalizable, one has to approximate them with physical states by fixing additional constraints as, typically, the energy constraint.
###### Acknowledgements.
This work has been supported by the FET European Networks on Quantum Information and Communication Contract IST-2000-29681:ATESIT, by MIUR 2003-Cofinanziamento, and by INFM PRA-2002-CLON.
## VI Appendix
### VI.1 Invariant cost functions
In this section we prove the form (18) of any invariant cost function.
###### Proposition 1
The following integral formula holds:
$$\mathrm{d}gU_g^\mu U_g^\nu =\delta _{\mu \nu }\frac{|๐_\mu ๐_\mu |}{d_\mu }.$$
(69)
Proof. Using Eq. (6), we recognize in the l.h.s. the projection onto the subspace of $`_\mu _\nu `$ that carries the trivial representation in the Clebsch-Gordan decomposition of $`U_g^\mu U_g^\nu `$. Using the orthogonality of characters, one can prove that such tensor product contains the trivial representation if and only if $`\mu =\nu `$. Moreover, if $`\mu =\nu `$, then the multiplicity of the trivial representation is one. Using the property (3), we see that the vector $`|๐_\mu `$ is invariant under $`U_g^\mu U_g^\mu `$. Therefore the r.h.s. is the projection onto the one-dimensional invariant subspace that carries the trivial representation, whence it coincides with the l.h.s. $`\mathrm{}`$
###### Proposition 2
Any invariant function $`c(\widehat{g},g)`$ has the form
$$c(\widehat{g},g)=\underset{\mu }{}a_\mu \chi ^\mu (g^1\widehat{g}),$$
(70)
where $`\chi ^\mu (g)\mathrm{Tr}[U_g^\mu ]`$.
Proof. For each irreducible representation $`\mu `$, consider the matrix elements $`u_{ij}^\mu (g)\psi _i^\mu |U_g^\mu |\psi _j^\mu `$ with respect to a fixed basis $`๐ก^\mu =\{|\psi _i^\mu |i=1,\mathrm{},d_\mu \}`$ for the representation space $`_\mu `$. Since the collection of all these matrix elements is an orthogonal basis for $`L^2(๐)`$ groups , we can expand the function $`c(\widehat{g},g)`$ as
$$c(\widehat{g},g)=\underset{\mu ,\nu }{}\underset{i,j=1}{\overset{d_\mu }{}}\underset{k,l=1}{\overset{d_\nu }{}}a_{ijkl}^{\mu \nu }u_{ij}^\mu (\widehat{g})u_{kl}^\nu (g),$$
(71)
where the complex conjugate in $`u_{kl}^\nu (g)`$ is for later convenience. Now, the function $`c`$ is both left- and right-invariant, whence it coincides with its average $`\overline{c}(\widehat{g},g)\mathrm{d}k\mathrm{d}hc(k\widehat{g}h,kgh)`$. Using Proposition 1 and Eqs. (2) and (3), we obtain
$$\begin{array}{cc}\hfill c(\widehat{g},g)& =\mathrm{d}k\mathrm{d}hc(k\widehat{g}h,kgh)\hfill \\ & =\underset{\mu ,\nu }{}\underset{i,j,k,l}{}a_{ijkl}^{\mu \nu }\delta _{\mu \nu }\times \hfill \\ & \times \psi _i^\mu |\psi _k^\mu |\frac{|๐_\mu ๐_\mu |}{d_\mu }(U_{\widehat{g}}^\mu U_g^\mu )\frac{|๐_\mu ๐_\mu |}{d_\mu }|\psi _j^\mu |\psi _l^\mu \hfill \\ & =\frac{1}{d_\mu ^2}\underset{\mu }{}\underset{i,j,l,k}{}a_{ijkl}^{\mu \mu }\delta _{ik}\delta _{jl}\mathrm{Tr}[U_{\widehat{g}g^1}^\mu ]\hfill \\ & =\underset{\mu }{}a_\mu \chi ^\mu (g^1\widehat{g}),\hfill \end{array}$$
where $`a_\mu \frac{1}{d_\mu ^2}_{i,j}a_{ijij}^{\mu \mu }`$.$`\mathrm{}`$
### VI.2 Decomposition of a product of maximally entangled states
Here we give the proof for Lemma 4.
Proof. Consider the representation $`๐ฑ(๐)=\{(U_g๐)^{}|๐_๐๐๐(๐)\}`$. It is convenient to order the $`2N`$ Hilbert spaces in the tensor product $`^{2N}`$ in such a way that the unitary operators act on the first $`N`$ spaces, while the identity operators acts on the second $`N`$ spaces. With this ordering, by defining $`_A`$ ($`_B`$) the tensor product of the first (second) $`N`$ spaces, we have $`๐ฑ(๐)=๐ฑ_A(๐)๐_๐น`$, where $`๐ฑ_A(๐)\{U_g^N|U_g๐๐(d)\}`$ is the $`N`$-fold tensor representation of $`๐๐(d)`$.
Letโs decompose now the Hilbert space $`_A`$ with respect to the action of the representation $`๐ฑ_A(๐)`$:
$$_A=\underset{\mu }{}_\mu ^{m_\mu }.$$
(72)
The tensor product $`_A_B`$ can be decomposed with respect to $`๐ฑ(๐)=๐ฑ_A(๐)๐_๐น`$ as
$$_A_B=\underset{\mu }{}_\mu ^{M_\mu },$$
(73)
where the multiplicity has been increased to $`M_\mu =m_\mu \times d^N`$, since $`_B`$ has been absorbed in the multiplicity spaces.
With respect to the factorization $`^{2N}=_A_B`$, the input state $`|\mathrm{\Psi }=|๐^{}`$ can be written as
$$|\mathrm{\Psi }=\frac{1}{\sqrt{d^N}}|๐^{},$$
(74)
where $`๐^{}`$ is the identity in $`_A^N_B`$. Here we are using notation (1), with respect to the product basis $`๐ก^N`$ for $`_A`$ and $`_B`$, $`๐ก`$ being a fixed basis for $``$. Now we want to change the basis in $`_A`$, by switching from $`๐ก^N`$ to $`๐ก^{}_\mu ๐ก_R^\mu ๐ก_M^\mu `$, where $`๐ก_R^\mu \{|\psi _n^\mu |n=1,\mathrm{},d_\mu \}`$ ($`๐ก_M^\mu \{|\varphi _n^\mu |n=1,\mathrm{},m_\mu \}`$) is a basis for the representation (multiplicity) space in Eq. (72). One has
$$\begin{array}{cc}\hfill |\mathrm{\Psi }& =\frac{1}{\sqrt{d^N}}\underset{\mu }{}\underset{m=1}{\overset{d_\mu }{}}\underset{n=1}{\overset{m_\mu }{}}|\psi _m^\mu _A|\varphi _n^\mu _A|\psi _m^\mu _B|\varphi _n^\mu _B\hfill \\ & =\frac{1}{\sqrt{d^N}}\underset{\mu }{}\underset{m=1}{\overset{d_\mu }{}}\sqrt{m_\mu }|\psi _m^\mu |\tau _m^\mu ,\hfill \end{array}$$
where we defined the normalized vector
$$|\tau _m^\mu \frac{1}{\sqrt{m_\mu }}\underset{n=1}{\overset{m_\mu }{}}|\varphi _n^\mu |\psi _n^\mu |\varphi _n^\mu .$$
(75)
Therefore, exploiting notation (1) with respect to the bases $`\{|\psi _m^\mu \}`$ and $`\{|\tau _m^\mu \}`$, we can write
$$|\mathrm{\Psi }=\underset{\mu }{}\sqrt{\frac{m_\mu }{d^N}}|๐_\mu .$$
(76)
$`\mathrm{}`$ |
warning/0506/hep-th0506028.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Quantum field theory in the space-time with boundaries has additional structures as compared to the theory when no boundaries are present -. One of the most important questions in this context is relation between the theories in the bulk and at the boundary, where the latter one describes the eigen-states (wave-functions) of the evolution operators of the former. This relation is believed to be especially important in topological and gravitational bulk theories with trivial Hamiltonians. Further, it is expected to be promoted to an even more fundamental open-closed string duality .
In this paper we consider a small subject in the story of bulk-to-boundary correspondence, which โ to the best of our knowledge โ did not attract enough attention in the literature. It concerns a possibility to form a nearly topological theory in the bulk, by substituting the bulk propagators by differences of those with Neumann and Dirichlet boundary conditions at the boundary. Such subtraction eliminates the physical pole from the propagator and โ if the boundary consists of the stable points of discrete $`Z_2`$ transformation โ substitutes it by โunphysicalโ pole at $`Z_2`$-image of the physical one. This situation is modeled by the theory of complex scalar field $`\chi `$ in the upper semi-space $`R_+^{d+1}`$ and the boundary $`R^d`$ located at $`z_0=0`$, with partition function
$$๐ต\{J\}=D\chi (z)D\overline{\chi }(z)\mathrm{exp}\{\frac{1}{g^2}_{R_+^{d+1}}(_\mu \chi _\mu \overline{\chi }+\underset{n}{}\frac{t_n}{n!}z_0^{s_n1}\chi ^n)d^{d+1}z+$$
$`+{\displaystyle \frac{1}{g^2}}{\displaystyle _{R^d}}(J_N(\stackrel{}{x})\mathrm{Re}\chi (\stackrel{}{x})+J_D(\stackrel{}{x})\mathrm{Im}_0\chi (\stackrel{}{x}))d\stackrel{}{x}\}`$ (1)
where functional integral is over the fields with mixed (Neumann-Dirichlet) boundary conditions
$`\begin{array}{c}\mathrm{Re}_0\chi (\stackrel{}{x},z_0)=0\\ \mathrm{Im}\chi (\stackrel{}{x},z_0)=0\end{array}|_{z_0=0}`$ (4)
and $`\stackrel{}{x}`$ are coordinates on the $`d`$-dimensional boundary (while coordinates in the bulk are $`z_0,\stackrel{}{z}`$). The summation over $`\mu `$ is with respect to the flat Euclidean metric. Be there no boundary, the theory would be trivial, since propagator converts $`\chi `$ into $`\overline{\chi }`$, while interaction terms are pure holomorphic: contain only $`\chi (z)`$, not $`\overline{\chi }(z)`$. Because of the boundary conditions the cancellation between the contributions of propagating $`\mathrm{Re}\chi (z)=\frac{1}{2}(\chi +\overline{\chi })(z)`$ and $`\mathrm{Im}\chi (z)=\frac{1}{2i}(\chi \overline{\chi })(z)`$ is not complete and partition function appears to be non-trivial. Still, it is particularly simple, and โ as we demonstrate below โ can be straightforwardly rewritten as an effective field theory of a single real field $`\phi (\stackrel{}{x})`$ at the boundary,
$`๐ต\{J\}={\displaystyle D\phi (\stackrel{}{x})\mathrm{exp}\left\{\frac{1}{g^2}_{R^d}\left(\frac{1}{2}\phi \sqrt{\mathrm{}}\phi +\underset{n}{}\frac{t_n}{n!}_0^{\mathrm{}}\left(e^\alpha \sqrt{\mathrm{}}\phi \right)^n\alpha ^{s_n1}๐\alpha +\left(J_N+J_D\sqrt{\mathrm{}}\right)\phi \right)๐\stackrel{}{x}\right\}}`$ (5)
Here $`\mathrm{}=_\stackrel{}{x}_\stackrel{}{x}`$ is minus Laplacian at the boundary and $`\alpha `$ is not a field, just a single auxiliary integrational variable.
The field $`\overline{\chi }`$ enters linearly in the action (1) and can work as Lagrange multiplier, providing a functional delta-function $`\delta \left((_0^2\mathrm{})\chi (z)+J(\stackrel{}{x})\delta (z_0)\right)`$. However, in the presence of boundary this condition does not fix $`\chi (z)`$ unambiguously in the bulk: a functional freedom remains in the zero-modes of the Laplace operator in $`R_+^{d+1}`$ and it is not eliminated by our mixed boundary conditions. Thus the theory remains non-trivial, just its degrees of freedom are actually those of a field on the boundary โ providing a non-trivial realization of holography idea.
This example can be easily extended to more general space-times with boundaries and boundaries can consist of the stable points of other $`Z_2`$ transformations.<sup>1</sup><sup>1</sup>1 In (1) it is the symmetry $`z_0z_0`$, which in the case of $`d+1=2`$ with complex $`z=z_1+iz_0`$ becomes complex conjugation $`z\overline{z}`$. Of considerable interest in the same dimension is the case when the role of space-time is played by a hyperelliptic Riemann surface $`y^2(z)=\mathrm{Polynomial}\mathrm{of}z`$, while the discrete transformation is different: $`y(z)y(z)`$. In this example our โunphysicalโ, Neumann minus Dirichlet, propagator is exactly the one that appears in the role of the two-point function $`\rho ^{(0|2)}(z,z^{})`$ in matrix model theory, with pole not at coincident, but at the $`Z_2`$-reflected points $`z^{}=z^{}`$, lying at two different sheets of the surface, see for recent presentations. Actually, one can use in the theory (1) the ordinary, physical, propagator, but than the interaction term becomes explicitly non-local: made from powers of $`\frac{1}{2}(\chi (z)+\overline{\chi }(z^{}))`$ instead of $`\chi (z)`$. For Euclidean space $`R^{d+1}`$ per se the index $`s_n=1`$, while for
$$s_n=(n2)\mathrm{\Delta }_{}1$$
(1) and (5) describe a sub-sector in the theory of scalar fields $`\varphi _\pm `$ of particular dimensions $`\mathrm{\Delta }_\pm =\frac{d\pm 1}{2}`$ in $`AdS_{d+1}`$, with $`\varphi _\mathrm{\Delta }_{}=z_0^\mathrm{\Delta }_{}\mathrm{Re}\chi (z)`$ and $`\varphi _{\mathrm{\Delta }_+}=z_0^\mathrm{\Delta }_{}\mathrm{Im}\chi (z)`$. Among other things, this fact manifests itself in conformal invariance of effective theory (5). Amusingly, even for the $`AdS_{d+1}`$ case, $`s_n`$ can sometime take value $`1`$: this happens when interaction shape is adjusted to space-time dimension so that $`(n2)((d+1)2)=4`$, i.e. when $`(d+1,n)=(6,3),(4,4)`$ or $`(3,6)`$. To avoid possible confusion, we emphasize that the action (1) is essentially complex, still most amplitudes are real, but the unitarity of the theory (inessential for our considerations) can be under question.
In the remaining part of this paper we briefly comment on straightforward derivations of (1) from the scalar theory in Euclidean $`AdS_{d+1}`$ and of (5) from (1).
## 2 Propagators in AdS
### 2.1 Action
The free action of real scalar field $`\varphi _\mathrm{\Delta }`$ in $`AdS_{d+1}`$ with Euclidean metric
$$ds^2=\frac{dz_0^2+d\stackrel{}{z}^2}{z_0^2}$$
is given by
$$S=\frac{1}{2}_{AdS_{d+1}}\sqrt{g}\left(g^{\mu \nu }_\mu \varphi _\mathrm{\Delta }_\nu \varphi _\mathrm{\Delta }+m_\mathrm{\Delta }^2\varphi _\mathrm{\Delta }^2\right)d^{d+1}z+S_{bound}=$$
$`={\displaystyle \frac{1}{2}}{\displaystyle _{AdS_{d+1}}}{\displaystyle \frac{d^{d+1}z}{z_0^{d1}}}\left((_0\varphi _\mathrm{\Delta })^2+(\stackrel{}{}\varphi _\mathrm{\Delta })^2+{\displaystyle \frac{m_\mathrm{\Delta }^2\varphi _\mathrm{\Delta }^2}{z_0^2}}\right)+S_{bound},`$ (6)
where the AdS mass is related to AdS dimension $`\mathrm{\Delta }`$ of $`\varphi _\mathrm{\Delta }`$ through
$$m_\mathrm{\Delta }^2=\mathrm{\Delta }(\mathrm{\Delta }d).$$
Dimension $`\mathrm{\Delta }`$ is restricted to $`\mathrm{\Delta }>\frac{d2}{2}`$ by the unitarity bound . This action can be transformed to the one in $`R_+^{d+1}`$ by rescaling of field variables<sup>2</sup><sup>2</sup>2 We thank V.Rubakov for useful discussions of this issue. $`\varphi _\mathrm{\Delta }(z)=z_0^\mathrm{\Delta }_{}\chi _\mathrm{\Delta }(z)`$:
$`S={\displaystyle \frac{1}{2}}{\displaystyle _{R_+^{d+1}}}d^{d+1}z\left((_0\chi _\mathrm{\Delta })^2+(\stackrel{}{}\chi _\mathrm{\Delta })^2+{\displaystyle \frac{m_\mathrm{\Delta }^2+\mathrm{\Delta }_+\mathrm{\Delta }_{}}{z_0^2}}\chi _\mathrm{\Delta }^2\right),`$ (7)
provided in (6)
$$S_{bound}=\frac{\mathrm{\Delta }_{}}{2}_{\left(AdS_{d+1}\right)}\frac{d^d\stackrel{}{z}}{z_0^d}\varphi _\mathrm{\Delta }^2=\frac{\mathrm{\Delta }_{}}{2}_{R^d}d^d\stackrel{}{x}\underset{z_0+0}{lim}\frac{\chi _\mathrm{\Delta }^2(z_0,\stackrel{}{x})}{z_0}$$
For special dimensions $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$ the AdS masses $`m_{\mathrm{\Delta }_\pm }^2=\mathrm{\Delta }_+\mathrm{\Delta }_{}`$, so that the last term at the r.h.s. of (7) vanishes,<sup>3</sup><sup>3</sup>3 Potentially solvable (rational Calogero) situation arises for entire integer-labeled tower of โdressedโ masses: $`M_\mathrm{\Delta }^2=\mathrm{\Delta }(\mathrm{\Delta }d)+\mathrm{\Delta }_+\mathrm{\Delta }_{}=N(N+1)`$, i.e. for $`\mathrm{\Delta }=\frac{1}{2}\left(d\pm (2N+1)\right)`$. We do not discuss this โ as many other โ obvious generalizations here. and the action converts into a free massless action for the scalar field in $`R^{d+1}`$:
$`S={\displaystyle \frac{1}{2}}{\displaystyle _{R_+^{d+1}}}d^{d+1}z\left((_0\chi _{\mathrm{\Delta }_\pm })^2+(\stackrel{}{}\chi _{\mathrm{\Delta }_\pm })^2\right)`$ (8)
One more peculiar feature of this particular choice of $`\mathrm{\Delta }`$ is that precisely at $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$ the scalar theory in the bulk acquires extended symmetry: local conformal invariance. Indeed, since the scalar curvature $`R(z)`$ of conformal metric $`g_{\mu \nu }(z)=\rho (z)\delta _{\mu \nu }`$ in $`d+1`$ dimensions is equal to
$$R=\frac{d}{\rho }\left(^2\mathrm{log}\rho +\frac{d1}{4}_\mu \mathrm{log}\rho _\mu \mathrm{log}\rho \right),$$
it follows that
$$\frac{1}{2}\sqrt{g}\left(g^{\mu \nu }_\mu \varphi _{\mathrm{\Delta }_\pm }_\nu \varphi _{\mathrm{\Delta }_\pm }+\xi R\varphi _{\mathrm{\Delta }_\pm }^2\right)d^{d+1}z$$
is invariant under the simultaneous change
$$\{\begin{array}{c}\rho (z)\lambda ^2(z)\rho (z),\\ \varphi _{\mathrm{\Delta }_\pm }(z)\lambda ^\mathrm{\Delta }_{}(z)\varphi _{\mathrm{\Delta }_\pm }(z)\end{array}$$
with arbitrary $`z`$-dependent $`\lambda (z)`$, satisfying $`_0\lambda (z_0=0)=0`$ provided
$$\xi =\frac{d1}{4d}=\frac{\mathrm{\Delta }_{}}{2d}$$
(for $`d+1=4`$ this $`\xi =\frac{1}{6}`$). For $`AdS_{d+1}`$ our $`\rho =\frac{1}{z_0^2}`$ and $`R=d(d+1)`$, so that $`\xi R(z)=\frac{(d1)(d+1)}{4}=\mathrm{\Delta }_{}\mathrm{\Delta }_+`$, i.e. exactly equals $`m_{\mathrm{\Delta }_\pm }^2`$. Interactions preserve conformal symmetry, both local and global, whenever $`s_n=1`$. An interesting question is how the local symmetry is realized in these cases in the boundary theory (5).
### 2.2 Bulk-to-boundary propagators
Whenever the boundary is Euclidean space โ as is the case with our parametrization of AdS โ it is practical to use mixed representation for Feynman diagrams: momentum along the boundary and coordinate in the orthogonal direction (inside the bulk). Fourier transform of the boundary variable $`\stackrel{}{x}`$ converts the AdS bulk-to-boundary propagator
$`\stackrel{~}{K}_\mathrm{\Delta }(w,\stackrel{}{x})={\displaystyle \frac{w_0^\mathrm{\Delta }}{\left[w_0^2+(\stackrel{}{w}\stackrel{}{x})^2\right]^\mathrm{\Delta }}}`$ (9)
into<sup>4</sup><sup>4</sup>4Schwinger parametrization is used to deal with the denominator:
$$\frac{1}{P^\alpha }=\frac{1}{\mathrm{\Gamma }(\alpha )}\underset{0}{\overset{\mathrm{}}{}}๐\lambda \lambda ^{\alpha 1}e^{\lambda P}$$
and the emerging integral is
$$\underset{0}{\overset{\mathrm{}}{}}\lambda ^{\nu 1}e^{\alpha \lambda \frac{\beta }{\lambda }}๐\lambda =2\left(\frac{\beta }{\alpha }\right)^{\frac{\nu }{2}}๐ฆ_\nu (2\sqrt{\alpha \beta })$$
$`K_\mathrm{\Delta }(w|\stackrel{}{p})={\displaystyle d^d\stackrel{}{x}e^{i\stackrel{}{p}\stackrel{}{x}}\frac{w_0^\mathrm{\Delta }}{\left[w_0^2+(\stackrel{}{w}\stackrel{}{x})^2\right]^\mathrm{\Delta }}}=w_0^\mathrm{\Delta }e^{i\stackrel{}{p}\stackrel{}{w}}\left({\displaystyle \frac{\stackrel{}{p}^2}{w_0^2}}\right)^{\frac{1}{2}\left(\mathrm{\Delta }\frac{d}{2}\right)}๐ฆ_{\mathrm{\Delta }\frac{d}{2}}\left(w_0\sqrt{\stackrel{}{p}^2}\right)`$ (10)
Here and below we systematically omit inessential factors to avoid overloading the formulas. Equalities are defined modulo such factors.
For semi-integer index $`\mathrm{\Delta }\frac{d}{2}`$ the modified Bessel function $`๐ฆ`$ turns into elementary function, especially simple for particular values of $`\mathrm{\Delta }=\mathrm{\Delta }_\pm =\frac{1}{2}(d\pm 1)`$. This follows from the intermediate formulas:
$$K_\mathrm{\Delta }(w|\stackrel{}{p})=w_0^\mathrm{\Delta }d^d\stackrel{}{x}e^{i\stackrel{}{p}\stackrel{}{x}}_0^{\mathrm{}}e^{\alpha (w_0^2+(\stackrel{}{w}\stackrel{}{x})^2)}\alpha ^{\mathrm{\Delta }1}๐\alpha =w_0^\mathrm{\Delta }e^{i\stackrel{}{p}\stackrel{}{w}}_0^{\mathrm{}}e^{\alpha w_0^2\frac{|p|^2}{4\alpha }}\alpha ^{\mathrm{\Delta }\frac{d}{2}1}๐\alpha =$$
$`=w_0^\mathrm{\Delta }e^{i\stackrel{}{p}\stackrel{}{w}}{\displaystyle _0^{\mathrm{}}}e^{\beta ^2w_0^2\frac{|p|^2}{4\beta ^2}}\beta ^{2\mathrm{\Delta }d1}๐\beta =w_0^\mathrm{\Delta }e^{i\stackrel{}{p}\stackrel{}{w}}{\displaystyle _0^{\mathrm{}}}e^{\frac{1}{4}\lambda ^2|p|^2\frac{w_0^2}{\lambda ^2}}\lambda ^{d2\mathrm{\Delta }1}๐\lambda `$ (11)
Here $`\alpha =\beta ^2=\lambda ^2`$ and $`|p|=\sqrt{(\stackrel{}{p})^2}`$. Distinguished are values of $`\mathrm{\Delta }`$ when the powers of either $`\beta `$ or $`\lambda `$ disappear from the pre-exponents, i.e. when $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$. The integral over $`\beta `$ (or $`\lambda `$) produces a factor of $`w_0^1`$ (or $`|p|^1`$) in the pre-exponent,<sup>5</sup><sup>5</sup>5Here the celebrated integral is used:
$$_0^{\mathrm{}}e^{a\lambda ^2\frac{b}{\lambda ^2}}๐\lambda =\sqrt{\frac{\pi }{4a}}e^{2\sqrt{ab}}$$
and we get:
$`K_{\mathrm{\Delta }_+}(w|\stackrel{}{p})=w_0^\mathrm{\Delta }_{}e^{i\stackrel{}{p}\stackrel{}{w}}e^{|p|w_0}`$ (12)
(since $`\mathrm{\Delta }_+1=\mathrm{\Delta }_{}`$), and
$`K_\mathrm{\Delta }_{}(w|\stackrel{}{p})={\displaystyle \frac{1}{|p|}}w_0^\mathrm{\Delta }_{}e^{i\stackrel{}{p}\stackrel{}{w}}e^{|p|w_0}={\displaystyle \frac{1}{|p|}}K_{\mathrm{\Delta }_+}(w|\stackrel{}{p}).`$ (13)
Eq.(13) implies that the amplitudes with external fields of dimension $`\mathrm{\Delta }=\mathrm{\Delta }_{}`$ at the boundary can be obtained from those of $`\mathrm{\Delta }=\mathrm{\Delta }_+`$ by insertion of appropriate $`|p|^1`$ factors.
### 2.3 Bulk-to-bulk propagators
Propagators in the space-time with boundary depend on boundary conditions. The natural physical requirement is that the flow $`\chi _\mathrm{\Delta }_0\chi _\mathrm{\Delta }`$ of field $`\chi _\mathrm{\Delta }`$ in (7) through the boundary vanishes . This means that only Dirichlet or Neumann boundary conditions can be imposed on the field $`\chi _\mathrm{\Delta }`$ at $`z_0=0`$, and our theory can be quantized for each value of $`\mathrm{\Delta }`$, corresponding to the same value of mass. Thus, the propagator of the scalar field $`\varphi _\mathrm{\Delta }`$ in $`AdS_{d+1}`$ satisfies
$`\left(\mathrm{}_{AdS}(w)+m_\mathrm{\Delta }^2\right)G_\mathrm{\Delta }(w,z)={\displaystyle \frac{1}{\sqrt{g(z)}}}\delta ^{(d+1)}(wz)=(z_0)^{d+1}\delta ^{(d+1)}(wz)`$ (14)
with additional requirement that
$`G_\mathrm{\Delta }(w,z)=w_0^\mathrm{\Delta }\mathrm{as}w_0+0.`$ (15)
Our consideration in subsection 2.1 implies that for particular values of $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$ this propagator is โ modulo the factor of $`(w_0z_0)^\mathrm{\Delta }_{}`$ โ the massless scalar propagator in $`R^{d+1}`$, i.e. Fourier transform of $`\left(p_0^2+\stackrel{}{p}^2\right)^1`$. The difference between $`\mathrm{\Delta }_+`$ and $`\mathrm{\Delta }_{}`$ is in the boundary condition (15):
$$G_{\mathrm{\Delta }_+}(w,z)=(w_0z_0)^\mathrm{\Delta }_{}\frac{d\stackrel{}{p}dp_0}{p_0^2+\stackrel{}{p}^2}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}\mathrm{sin}(p_0z_0)\mathrm{sin}(p_0w_0)=$$
$`=(w_0z_0)^\mathrm{\Delta }_{}\left\{{\displaystyle \frac{1}{\left((w_0z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2\right)^\mathrm{\Delta }_{}}}{\displaystyle \frac{1}{\left((w_0+z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2\right)^\mathrm{\Delta }_{}}}\right\}={\displaystyle \frac{1}{u^\mathrm{\Delta }_{}}}{\displaystyle \frac{1}{(u+2)^\mathrm{\Delta }_{}}}`$ (16)
satisfies Dirichlet, while
$$G_\mathrm{\Delta }_{}(w,z)=(w_0z_0)^\mathrm{\Delta }_{}\frac{d\stackrel{}{p}dp_0}{p_0^2+\stackrel{}{p}^2}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}\mathrm{cos}(p_0z_0)\mathrm{cos}(p_0w_0)=$$
$`=(w_0z_0)^\mathrm{\Delta }_{}\left\{{\displaystyle \frac{1}{\left((w_0z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2\right)^\mathrm{\Delta }_{}}}+{\displaystyle \frac{1}{\left((w_0+z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2\right)^\mathrm{\Delta }_{}}}\right\}={\displaystyle \frac{1}{u^\mathrm{\Delta }_{}}}+{\displaystyle \frac{1}{(u+2)^\mathrm{\Delta }_{}}}`$ (17)
โ Neumann boundary conditions at the boundary of $`AdS`$ (at $`w_0=0`$ or $`z_0=0`$). Integrals in (16) and (17) are evaluated by introduction of auxiliary integration,
$$\frac{1}{p_0^2+\stackrel{}{p}^2}=_0^{\mathrm{}}e^{\alpha (p_0^2+\stackrel{}{p}^2)}๐\alpha ,$$
followed by Gaussian integration over $`p_0`$ and $`\stackrel{}{p}`$. Switching to $`\lambda =(4\alpha )^1`$ we obtain:
$$G_{\mathrm{\Delta }_+}(w,z)=(w_0z_0)^\mathrm{\Delta }_{}_0^{\mathrm{}}\lambda ^{\mathrm{\Delta }_{}1}๐\lambda \left(\mathrm{exp}\left[\lambda \left((w_0z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2\right)\right]\mathrm{exp}\left[\lambda \left((w_0+z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2\right)\right]\right)$$
and integral over $`\lambda `$ provides (16). The last representations in formulas (16) and (17) are in terms of the usual AdS variables
$$u=\frac{(w_0z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2}{2w_0z_0}\mathrm{and}u+2=\frac{(w_0+z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2}{2w_0z_0}$$
For the sake of completeness, in Appendix A at the end of the paper, we present alternative derivation of these propagators: as solutions of the hypergeometric equation.
## 3 The difference of AdS propagators and the theory (1)
Now we are ready to introduce our simplified AdS theory. Suggestion is to substitute the propagators (17) and (16) by their peculiar linear combination: subtract one from another and define
$$G_0(w,z)\frac{1}{2}\left(G_\mathrm{\Delta }_{}(w,z)G_{\mathrm{\Delta }_+}(w,z)\right)=(w_0z_0)^\mathrm{\Delta }_{}\frac{d\stackrel{}{p}dp_0}{p_0^2+\stackrel{}{p}^2}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}\mathrm{cos}\left(p_0(w_0+z_0)\right)=$$
$$=\frac{1}{(u+2)^\mathrm{\Delta }_{}}=\left(\frac{w_0z_0}{(w_0+z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2}\right)^\mathrm{\Delta }_{}$$
as the propagator of a new conformally invariant scalar field theory. The spatial Fourier transform of such bulk-to-bulk propagator is especially simple:
$`G_0(w,z)=(w_0z_0)^\mathrm{\Delta }_{}{\displaystyle \frac{d^d\stackrel{}{p}}{|p|}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}e^{|p|(w_0+z_0)}}`$ (18)
The fact that $`z_0`$ and $`w_0`$ appear in the exponent as a simple sum<sup>6</sup><sup>6</sup>6 Unlike they show up in Fourier transform of the physical propagator $`๐ข=\frac{1}{2}(G_\mathrm{\Delta }_{}+G_{\mathrm{\Delta }_+})1/u^\mathrm{\Delta }_{}`$, which contains a far more sophisticated factor of $`e^{i|p||w_0z_0|}`$, see appendix C below., makes convolutions of propagators in expressions for Feynman diagrams very simple and leads to especial simplicity of effective theory (5).
What is the way to realize such projection โ from two different propagators to their difference? A possible answer is provided by the theory of two real scalar fields $`\varphi _\mathrm{\Delta }_{}`$ and $`\varphi _{\mathrm{\Delta }_+}`$, subjected to Neumann and Dirichlet boundary conditions respectively with the action
$$S=d^{d+1}z\sqrt{g}\left[\frac{1}{2}g^{\mu \nu }_\mu \varphi _\mathrm{\Delta }_{}_\nu \varphi _\mathrm{\Delta }_{}+\frac{1}{2}m_{\mathrm{\Delta }_\pm }^2\varphi _\mathrm{\Delta }_{}^2+\frac{1}{2}g^{\mu \nu }_\mu \varphi _{\mathrm{\Delta }_+}_\nu \varphi _{\mathrm{\Delta }_+}+\frac{1}{2}m_{\mathrm{\Delta }_\pm }^2\varphi _{\mathrm{\Delta }_+}^2+\underset{n}{}\frac{t_n}{n!}(\varphi _\mathrm{\Delta }_{}+i\varphi _{\mathrm{\Delta }_+})^n\right]$$
To prove that this theory reproduces the bulk-to-bulk propagator $`G_0(w,z)`$, consider the simplest case, when only $`t_30`$ in the interaction terms. In the figure 1, the tree contribution to the four-point function of $`\varphi _{\mathrm{\Delta }_+}`$ is shown.
There are two diagrams: one with the bulk-to-bulk propagator $`G_\mathrm{\Delta }_{}(w,z)`$, another one with $`G_{\mathrm{\Delta }_+}(w,z)`$. The vertex contribution for $`\varphi _{\mathrm{\Delta }_+}^3`$ interaction term is $`V_{+++}=i^3t_3\sqrt{g}`$, while for $`\varphi _{\mathrm{\Delta }_+}^2\varphi _\mathrm{\Delta }_{}`$ is $`V_{++}=i^2t_3\sqrt{g}`$. Thus, $`V_{+++}=iV_{++}`$. This extra $`i`$ factor, being squared (since there are two vertices in each diagram), gives relative minus sign between these two diagrams. Hence, effectively, the propagator in this theory is $`G_\mathrm{\Delta }_{}(w,z)G_{\mathrm{\Delta }_+}(w,z)=G_0(w,z)`$. The generalization to other tree and loop diagrams is straighforward.
In terms of fields
$`\varphi =\varphi _\mathrm{\Delta }_{}+i\varphi _{\mathrm{\Delta }_+}`$
$`\overline{\varphi }=\varphi _\mathrm{\Delta }_{}i\varphi _{\mathrm{\Delta }_+}`$
this action can be rewritten as
$$S=d^{d+1}z\sqrt{g}\left(\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \overline{\varphi }+\frac{m_{\mathrm{\Delta }_\pm }^2}{2}\varphi \overline{\varphi }+\underset{n}{}\frac{t_n}{n!}\varphi ^n\right).$$
(19)
Making rescaling of variables
$`\varphi _\mathrm{\Delta }_{}=z_0^\mathrm{\Delta }_{}\chi _\mathrm{\Delta }_{}`$
$`\varphi _{\mathrm{\Delta }_+}=z_0^\mathrm{\Delta }_{}\chi _{\mathrm{\Delta }_+}`$
from the section 2.1, we obtain, up to boundary terms, the action
$`S={\displaystyle d^{d+1}z\left[\frac{1}{2}_\mu \chi _\mu \overline{\chi }+\underset{n}{}\frac{t_n}{n!}z_0^{s_n1}\chi ^n\right]}`$ (20)
for the fields $`\chi =\chi _\mathrm{\Delta }_{}+i\chi _{\mathrm{\Delta }_+}`$ and $`\overline{\chi }=\chi _\mathrm{\Delta }_{}i\chi _{\mathrm{\Delta }_+}`$. Index $`s_n=(n2)\mathrm{\Delta }_{}1`$.
## 4 Tree diagrams
In this section we evaluate arbitrary tree correlators of the fields of dimension $`\mathrm{\Delta }_+`$ on the boundary of $`AdS_{d+1}`$ in the complex-scalar field theory (1), or equivalently (19) with holomorphic interaction. It turns out that all the tree diagrams are expressed as simple rational functions of $`d`$-dimensional (rather than $`d+1`$-dimensional) momenta: one extra dimension can be explicitly integrated out and provide a non-local, but rather simple effective theory (5). Moreover, the same effective theory appears to describe all loop diagrams as well.
### 4.1 Single-vertex diagram
We begin from the simplest diagram with a single vertex of valence $`n+1`$, see Fig.2.
The corresponding amplitude is given by
$$A(w|\stackrel{}{p}_1,\mathrm{},\stackrel{}{p}_n)=t_{n+1}\frac{dz_0d^d\stackrel{}{z}}{z_0^{d+1}}G_0(w,z)K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_1)\mathrm{}K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_n)\stackrel{(\text{18})\&(\text{12})}{=}$$
$$=t_{n+1}\frac{dz_0d^d\stackrel{}{z}}{z_0^{d+1}}\left((w_0z_0)^\mathrm{\Delta }_{}\frac{d\stackrel{}{p}}{|p|}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}e^{|p|(w_0+z_0)}\right)\left(z_0^\mathrm{\Delta }_{}e^{i\stackrel{}{p}_1\stackrel{}{z}}e^{|p_1|z_0}\right)\mathrm{}\left(z_0^\mathrm{\Delta }_{}e^{i\stackrel{}{p}_n\stackrel{}{z}}e^{|p_n|z_0}\right)=$$
$`=t_{n+1}{\displaystyle \frac{w_0^\mathrm{\Delta }_{}}{|p_{1n}|\left(|p_1|+\mathrm{}+|p_n|+|p_{1n}|\right)^{s_{n+1}}}}e^{i\stackrel{}{w}(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)}e^{|p_{1n}|w_0}`$ (21)
where $`|p_{1n}|\sqrt{(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)^2}`$. Note, that the use of projection (18) instead of (16) or (17) eliminates the poles at collinear momenta, when at $`|p_{1n}|=|p_1|+\mathrm{}+|p_n|`$, which normally occur in AdS Feynman diagrams, see Appendix $`C`$ below.
### 4.2 Multi-vertex diagrams
Comparison of (21) and (12) shows that their $`w`$-dependencies are exactly the same. This implies universality of the vertex insertion and the possibility of recursive evaluation of tree diagrams. Namely we can factor out all the $`w`$ dependence in a simple and universal manner:
$`A^\mathrm{\Gamma }(w|\stackrel{}{p}_1,\mathrm{},\stackrel{}{p}_n)=w_0^\mathrm{\Delta }_{}e^{i\stackrel{}{w}(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)}e^{w_0|\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n|}B^\mathrm{\Gamma }(\stackrel{}{p}_1,\mathrm{},\stackrel{}{p}_n)=K_{\mathrm{\Delta }_+}(w|\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)B^\mathrm{\Gamma }(\stackrel{}{p}_1,\mathrm{},\stackrel{}{p}_n),`$ (22)
where $`\mathrm{\Gamma }`$ labels the graph = diagram, which in our case is the rooted tree with the total momentum $`\stackrel{}{p}_\mathrm{\Gamma }\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n`$.
Eq.(22) is pictorially represented in Fig.3. It shows that in the theory (1) or equivalently (19) any number of bulk-to-boundary propagators, attached to an intermediate vertex at $`z`$ in the bulk, can after integration over $`z`$ be substituted by a single bulk-to-boundary propagator, leading to the point $`w`$ in the bulk and carrying the total along-the-boundary momentum, multiplied by the factor $`B^\mathrm{\Gamma }(\stackrel{}{p}_1,\mathrm{},\stackrel{}{p}_n)`$.
At the next step we evaluate the second-level diagram like the left one in Fig.4 with several single-vertex sub-diagrams meeting at the single vertex at point $`y`$ in the bulk. As shown in the picture, repeated application of above-described contraction provides an amplitude, again proportional to a single bulk-to-boundary propagator.
We can apply this algorithm as many times as necessary, converting arbitrary tree diagram with the root at $`w`$ inside the bulk into a single bulk-to-boundary propagator $`K_{\mathrm{\Delta }_+}(w|\mathrm{total}\mathrm{momentum})`$. The coefficient function arises from recurrent relation
$`B^\mathrm{\Gamma }=t_{n+1}{\displaystyle \frac{B^{\mathrm{\Gamma }_1}\mathrm{}B^{\mathrm{\Gamma }_n}}{|p_\mathrm{\Gamma }|\left(|p_{\mathrm{\Gamma }_1}|+\mathrm{}+|p_{\mathrm{\Gamma }_n}|+|p_\mathrm{\Gamma }|\right)^{s_{n+1}}}}`$ (23)
From this relation it is clear, that induced diagrammatic technique, is local: each vertex contribution depends on momenta, incoming into this particular vertex only. This locality is due to specific choice of propagator $`G_0(w,z)`$.
To finish evaluation of the diagram in Fig.5 it remains to attach the last bunch of bulk-to-boundary propagators to the root $`w`$ of the tree (Fig.4 left) and integrate over $`w`$. Such convolution of $`m`$ propagators is equal to:
$`{\displaystyle K_{\mathrm{\Delta }_+}(w|\stackrel{}{p}_1)\mathrm{}K_{\mathrm{\Delta }_+}(w|\stackrel{}{p}_m)\frac{dw_0d\stackrel{}{w}}{w_0^{d+1}}}={\displaystyle \frac{\delta (\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_m)}{\left(|p_1|+\mathrm{}+|p_m|\right)^{s_m}}}`$ (24)
Following this procedure one can explicitly evaluate any particular tree diagram.
### 4.3 Feynman rules and effective action
The result of our consideration can be formulated as simple Feynman rules for the $`B^\mathrm{\Gamma }`$ functions:
| propagator (line) | $`\frac{1}{|p|}`$ |
| --- | --- |
| $`n`$-vertex | $`t_n\frac{\delta ^{(d)}(\stackrel{}{p}_1+\stackrel{}{p}_2+\mathrm{}+\stackrel{}{p}_n)}{\left(|p_1|+|p_2|+\mathrm{}+|p_n|\right)^{s_n}}`$ |
External lines carry the same factors $`|p|^1`$ provided external fields have dimension $`\mathrm{\Delta }_{}`$, while for dimension $`\mathrm{\Delta }_+`$ external lines carry no factors of momentum. Moreover, as explained in section 5 below, the same rules work for loop calculations.
This diagram technique can be summarized in the form of a simple effective action: tree and loop diagrams for scalars of dimension $`\mathrm{\Delta }_\pm `$ at the boundary and peculiar propagator $`G_0(w,z)`$ in the bulk coincide with the tree and loop diagrams in the boundary theory (5):
$`๐ต\{J\}={\displaystyle D\phi (\stackrel{}{x})\mathrm{exp}\left\{\frac{1}{g^2}๐\stackrel{}{x}\left(\frac{1}{2}\phi \sqrt{\mathrm{}}\phi +\left(J_N+J_D\sqrt{\mathrm{}}\right)\phi +\underset{n}{}\frac{t_n}{n!}_0^{\mathrm{}}\left(e^\alpha \sqrt{\mathrm{}}\phi \right)^n\alpha ^{s_n1}๐\alpha \right)\right\}}`$ (25)
This theory is conformal invariant: rescalings of field $`\phi `$, which has dimension $`\mathrm{\Delta }_{}`$ are accompanied by the transformation of auxiliary integration variable $`\alpha `$, which has dimension $`1`$. In variance with the boundary models, usually considered in the context of AdS/CFT correspondence -, , -, the theory (25) is non-local. Also trees (loops) in the bulk are the same trees (loops) in the theory (25).
We emphasize that matching between theories (1) and (25) is valid for arbitrary value of index $`s_n`$ (not restricted to $`s_n=(n2)\mathrm{\Delta }_{}1`$). Only for particular choice of $`s_n=(n2)\mathrm{\Delta }_{}1`$, the bulk theory (1) describes the AdS theory (19).
## 5 Loops
It is straightforward to see that our Feynman rules โ and thus the effective theory (5) at the boundary โ reproduce expressions for loop diagrams in AdS theory. Again, in variance with the usual AdS/CFT correspondence, loops in (19) are loops in (25).
### 5.1 Sample 1-loop diagram.
We provide just an illustration. Consider the diagram in Fig.6.
Original expression is
$`t_nt_m{\displaystyle \frac{dz_0d^d\stackrel{}{z}}{z_0^{d+1}}\frac{dw_0d^d\stackrel{}{w}}{w_0^{d+1}}K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_1)\mathrm{}K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_{n2})G_0^2(z,w)K_{\mathrm{\Delta }_+}(w|\stackrel{}{q}_1)\mathrm{}K_{\mathrm{\Delta }_+}(w|\stackrel{}{q}_{m2})}`$ (26)
According to (12) and (18) the product of propagators is equal to
$$z_0^{\mathrm{\Delta }_{}(n2)}e^{i\stackrel{}{z}_i\stackrel{}{p}_i}e^{z_0_i|p_i|}\left((z_0w_0)^\mathrm{\Delta }_{}\frac{d^d\stackrel{}{r}}{|r|}e^{i\stackrel{}{r}(\stackrel{}{z}\stackrel{}{w})}e^{|r|(w_0+z_0)}\right)^2w_0^{\mathrm{\Delta }_{}(m2)}e^{i\stackrel{}{w}_j\stackrel{}{q}_j}e^{w_0_j|q_j|}$$
and after integration over $`\stackrel{}{z}`$ and $`\stackrel{}{w}`$ (26) turns into
$$t_nt_m\frac{d^d\stackrel{}{r}}{|r|}\frac{d^d\stackrel{}{r}^{}}{|r^{}|}\delta ^{(d)}(\underset{i}{}\stackrel{}{p}_i+\stackrel{}{r}+\stackrel{}{r}^{})\delta ^{(d)}(\underset{j}{}\stackrel{}{q}_j\stackrel{}{r}\stackrel{}{r}^{})\times $$
$$\times _0^{\mathrm{}}z_0^{(n2)\mathrm{\Delta }_{}2}dz_0_0^{\mathrm{}}w_0^{(m2)\mathrm{\Delta }_{}2}dw_0e^{(_i|p_i|+|r|+|r^{}|)z_0}e^{(_j|q_j|+|r|+|r^{}|)w_0}=$$
$$=t_nt_m\delta ^{(d)}\left(\underset{i}{}\stackrel{}{p}_i+\underset{j}{}\stackrel{}{q}_j\right)\frac{d^d\stackrel{}{r}}{|r||r^{}|\left(_i|p_i|+|r|+|r^{}|\right)^{s_n}\left(_j|q_j|+|r|+|r^{}|\right)^{s_m}}$$
where in the last formula $`|r^{}|=\sqrt{\left(_i\stackrel{}{p}_i+\stackrel{}{r}\right)^2}=\sqrt{\left(_j\stackrel{}{q}_j\stackrel{}{r}\right)^2}`$. This expression is exactly the one which arises in effective theory (5) for the loop diagram of the same topology, Fig.6, .
Generalizations to all other loop diagrams is straightforward.
### 5.2 Tadpoles and divergencies
A few comments are deserved by tadpole diagrams, like those shown in Figs.7,8. First of all, the projected propagator $`G_0(z,w)`$ โ in variance from the usual ones, like $`G_{\mathrm{\Delta }_\pm }(z,w)`$ โ is not singular at coincident points:
$$G_0(z,z)=z_0^{2\mathrm{\Delta }_{}}\frac{d^d\stackrel{}{p}}{|p|}e^{2|p|z_0}=const$$
(this is also clear from its expression through $`u`$-variable, since at $`w=z`$ this $`u=0`$, but $`u+2=20`$). Thus the UV divergences are absent, as one expects in the non-local theory (5) with exponential damping of interactions at large momenta.
However, some peculiar divergences still survive. For example, for the diagram in Fig.7 we have, according to our Feynman rules:
$`t_n\delta ^{(d)}\left({\displaystyle \underset{i=1}{\overset{n2}{}}}\stackrel{}{p}_i\right){\displaystyle \frac{d^d\stackrel{}{q}}{|q|\left(2|q|+_{i=1}^{n2}|p_i|\right)^{s_n}}}=t_n{\displaystyle \frac{\delta ^{(d)}\left(_{i=1}^{n2}\stackrel{}{p}_i\right)}{\left(_{i=1}^{n2}|p_i|\right)^{(n4)\mathrm{\Delta }_{}1}}}`$ (27)
More accurately, expression in the r.h.s. is valid only when $`(n4)\mathrm{\Delta }_{}>1`$, otherwise the integral in the l.h.s. diverges at large $`|q|`$. (Note, that vanishing sum of the space momenta $`\stackrel{}{p}`$ does not imply that the sum of their moduli, $`|p|`$, is zero. This quantity is always positive.) If we rewrite (27) in coordinate representation, making use of $`G_0(z,z)=const`$, we get:
$$t_n\frac{dz_0d^d\stackrel{}{z}}{z_0^{d+1}}K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_1)\mathrm{}K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_{n2})G_0(z,z)=t_n\delta ^{(d)}\left(\underset{i=1}{\overset{n2}{}}\stackrel{}{p}_i\right)_0^{\mathrm{}}e^{z_0_i|p_i|}z_0^{(n4)\mathrm{\Delta }_{}2}๐z_0$$
In this representation it is clear that divergence, though looking like UV from the point of theory (5), comes from the region of small $`z_0`$, from the vicinity of the boundary of AdS and thus actually an IR divergence.
Particular diagram in Fig.7 converges if $`n6`$, we consider $`d2`$, and divergence can be eliminated, say, by putting $`t_3=t_4=t_5=0`$. However, the two-loop diagram in Fig.8, associated with the coupling $`t_6`$, diverges for exactly the same reason as the one loop in Fig.7.
### 5.3 Coleman-Weinberg potential
The boundary theory (5) can be used to sum up the one-loop diagrams in the background of constant field $`\mathrm{\Phi }`$. The result of this summation โ the Coleman-Weinberg action for our theory โ is<sup>7</sup><sup>7</sup>7We use the standard background field technique. The field $`\phi `$ is decomposed into the sum of background and quantum field: $`\phi =\mathrm{\Phi }+\phi _q`$. We find the part of the action proportional to $`\phi _q^2`$, then diagonalize it in momentum representation. Arrow implies the use of relation $`\mathrm{log}\mathrm{Det}=\mathrm{Tr}\mathrm{log}`$ and subtraction of the $`\mathrm{\Phi }`$-independent part from $`S_{CW}(\mathrm{\Phi })`$. The one-loop contribution to $`\mathrm{log}๐ต\{J\}`$ is related to $`S_{CW}(\mathrm{\Phi })`$ by Legendre transform.
$`S_{CW}\{\mathrm{\Phi }\}\mathrm{log}\mathrm{Det}\left(\sqrt{\mathrm{}}+{\displaystyle \underset{n}{}}{\displaystyle \frac{t_n\mathrm{\Phi }^{n2}}{(n2)!}}{\displaystyle \frac{1}{(2\sqrt{\mathrm{}})^{s_n}}}\right){\displaystyle \mathrm{log}\left(1+\underset{n}{}\frac{t_n}{(n2)!}\frac{\mathrm{\Phi }^{n2}}{|q|(2|q|)^{s_n}}\right)d^d\stackrel{}{q}}`$ (28)
For $`s_n+1>d`$, i.e. for $`(n4)\mathrm{\Delta }_{}>1`$ the integral converges at large $`|q|`$ (otherwise there is our familiar IR divergence at small $`z_0`$), and it converges for small $`|q|`$, though individual terms of power expansion in $`\mathrm{\Phi }`$ are divergent.
It is instructive to compare this Coleman-Weinberg action with the naive open-sector boundary effective action, associated with the free-field theory with quadratic vertex operators:
$`Z_{op}\{I\}={\displaystyle D\stackrel{~}{\phi }(\stackrel{}{x})\mathrm{exp}\left\{\frac{1}{g^2}d^d\stackrel{}{x}\left(\stackrel{~}{\phi }\mathrm{}\stackrel{~}{\phi }+I\stackrel{~}{\phi }^2\right)\right\}}\mathrm{Det}^{1/2}\left(\mathrm{}+I\right)`$ (29)
For constant source $`I(\stackrel{}{x})=const`$
$`S_{op}(I)=\mathrm{log}Z_{op}\{I\}{\displaystyle \mathrm{log}\left(1+\frac{I}{q^2}\right)d^d\stackrel{}{q}}`$ (30)
and does not have anything to do with $`\mathrm{log}๐ต\{J\}`$. Still, for distinguished value of index $`s_n=1`$ and for $`I=\mathrm{\Phi }^{n2}`$, it coincides with $`S_{CW}(\mathrm{\Phi })`$ โ the one-loop contribution to Legendre transform of $`\mathrm{log}๐ต\{J\}`$ at constant $`\mathrm{\Phi }`$. Note also that the dimensions $`\mathrm{\Delta }=d2`$, considered in coincide with our $`\mathrm{\Delta }_\pm =\frac{d\pm 1}{2}`$ for $`d=3`$ and $`d=5`$.
## 6 Appendix A. Propagators as hypergeometric functions
Usually in the literature eq.(14) is not solved explicitly for particular dimensions $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$, as we did in this paper. Instead one uses the transcendental expression for the bulk-to-bulk propagator through hypergeometric series ,
$`G_\mathrm{\Delta }(w,z)={\displaystyle \frac{1}{(4\pi )^{\frac{d+1}{2}}}}{\displaystyle \frac{\mathrm{\Gamma }(\mathrm{\Delta })\mathrm{\Gamma }(\mathrm{\Delta }\mathrm{\Delta }_{})}{\mathrm{\Gamma }(2\mathrm{\Delta }2\mathrm{\Delta }_{})}}\left({\displaystyle \frac{2}{u}}\right)^\mathrm{\Delta }F(\mathrm{\Delta },\mathrm{\Delta }\mathrm{\Delta }_{},2\mathrm{\Delta }2\mathrm{\Delta }_{};{\displaystyle \frac{2}{u}})`$ (31)
In this section we briefly explain, how our simple calculations are related to this standard approach.
For particular dimensions $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$, up to an overall normalization factor $`1/(4\pi )^{\frac{d+1}{2}}`$, we get<sup>8</sup><sup>8</sup>8 Here we use the definition of Gauss hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$,
$$F(a,b,c;z)=\frac{\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(n+a)\mathrm{\Gamma }(n+b)}{n!\mathrm{\Gamma }(n+c)}z^n,$$
which solves the hypergeometric equation,
$$z(1z)F^{\prime \prime }+\left(c(a+b+1)z\right)F^{}abF=0,$$
and Newtonโs binomial expansion
$$(1z)^s=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(n+s)}{n!\mathrm{\Gamma }(s)}z^n=F(s,c,c;z)$$
from (31) :
$$G_\mathrm{\Delta }_{}=\frac{2^{\mathrm{\Delta }_{}+1}\mathrm{\Gamma }(\mathrm{\Delta }_{})}{u^\mathrm{\Delta }_{}}F(\mathrm{\Delta }_{},0,0;\frac{2}{u})=$$
$`={\displaystyle \frac{2^\mathrm{\Delta }_{}\mathrm{\Gamma }(\mathrm{\Delta }_{})}{u^\mathrm{\Delta }_{}}}\left(1+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n+\mathrm{\Delta }_{})}{\mathrm{\Gamma }(\mathrm{\Delta }_{})n!}}\left({\displaystyle \frac{2}{u}}\right)^n\right)=2^\mathrm{\Delta }_{}\mathrm{\Gamma }(\mathrm{\Delta }_{})\left({\displaystyle \frac{1}{u^\mathrm{\Delta }_{}}}+{\displaystyle \frac{1}{(u+2)^\mathrm{\Delta }_{}}}\right)`$ (32)
and
$$G_{\mathrm{\Delta }_+}=\frac{2^{\mathrm{\Delta }_+}\mathrm{\Gamma }(\mathrm{\Delta }_+)}{u^{\mathrm{\Delta }_+}}F(\mathrm{\Delta }_+,1,2;\frac{2}{u})=\frac{2^{\mathrm{\Delta }_+}}{u^{\mathrm{\Delta }_+}}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }(n+\mathrm{\Delta }_+)\mathrm{\Gamma }(n+1)}{n!\mathrm{\Gamma }(n+2)}\left(\frac{2}{u}\right)^n=$$
$`={\displaystyle \frac{2^{\mathrm{\Delta }_+}}{2u^\mathrm{\Delta }_{}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n+\mathrm{\Delta }_{})}{n!}}\left({\displaystyle \frac{2}{u}}\right)^n=2^\mathrm{\Delta }_{}\mathrm{\Gamma }(\mathrm{\Delta }_{})\left({\displaystyle \frac{1}{u^\mathrm{\Delta }_{}}}{\displaystyle \frac{1}{(u+2)^\mathrm{\Delta }_{}}}\right)`$ (33)
An extra coefficient $`2`$ in the first line of (32) comes from the fact that the ratio $`\frac{\mathrm{\Gamma }(\mathrm{\Delta }\mathrm{\Delta }_{})}{\mathrm{\Gamma }(2\mathrm{\Delta }2\mathrm{\Delta }_{})}2`$ as $`\mathrm{\Delta }\mathrm{\Delta }_{}`$. For the same reason the $`n=0`$ term in the sum enters with an extra coefficient $`2`$, this is taken into account by an extra item $`1`$ in the first formula in the second line of (32).
Up to a common normalization factor the formulas (32) and (33) reproduce (17) and (16) respectively.
## 7 Appendix B. Derivations a la
We illustrate this standard method by re-examining the example of (21), with $`G_0(w,z)`$ replaced by $`G_\mathrm{\Delta }(w,z)`$. We put $`n=3`$ (triple vertex) and $`d=5`$ to simplify the formulas. The results of this appendix can be generalized to arbitrary values of $`n`$ and $`d`$ in an obvious way. Following we extract this quantity from solution of the equation (14):
$`(\mathrm{}_{AdS}+m_\mathrm{\Delta }^2)A_\mathrm{\Delta }(w|\stackrel{}{p}_1,\stackrel{}{p}_2)=K_{\mathrm{\Delta }_\pm }(w|\stackrel{}{p}_1)K_{\mathrm{\Delta }_\pm }(w|\stackrel{}{p}_2)e^{(|p_1|+|p_2|)w_0}e^{i(\stackrel{}{p}_1+\stackrel{}{p}_2)\stackrel{}{w}}w_0^{2\mathrm{\Delta }_{}}`$ (34)
the subscript $`\mathrm{\Delta }`$ for $`A_\mathrm{\Delta }(w|\stackrel{}{p}_1,\stackrel{}{p}_2)`$ denotes dimension of the bulk-to-bulk propagator, eq.(35) below restricts it to be $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$. Proportionality sign in (34) shows that the possible $`|p_1|^1`$ and $`|p_2|^1`$ factors from $`K_\mathrm{\Delta }_{}`$ are ignored. Substituting $`A_\mathrm{\Delta }(w|\stackrel{}{p}_1,\stackrel{}{p}_2)=e^{i\stackrel{}{w}(\stackrel{}{p}_1+\stackrel{}{p}_2)}e^{t\mu }F(t)`$ with $`t=w_0\sqrt{(\stackrel{}{p}_1+\stackrel{}{p}_2)^2}=w_0|p_{12}|`$ and $`\mu =\frac{|p_1|+|p_2|}{|p_{12}|}`$, we obtain for $`F(t)`$:
$$t^2F^{\prime \prime }+(2\mu t^2+(d1)t)F^{}+\left((1\mu ^2)t^2\mu (d1)t+\mathrm{\Delta }(\mathrm{\Delta }d)\right)F=t^{2\mathrm{\Delta }_{}}|p_{12}|^{2\mathrm{\Delta }_{}}$$
Substituting further $`F(t)=t^\mathrm{\Delta }_{}f(t)`$, we get:
$`f^{\prime \prime }+2\mu f^{}+(1\mu ^2)f=t^{2+\mathrm{\Delta }_{}}|p_{12}|^{2\mathrm{\Delta }_{}}`$ (35)
Generic solution of this equation is:
$`A_\mathrm{\Delta }(w|\stackrel{}{p}_1,\stackrel{}{p}_2)={\displaystyle \frac{w_0^\mathrm{\Delta }_{}}{|p_{12}|^2(|p_1|+|p_2|)^2}}e^{i\stackrel{}{w}(\stackrel{}{p}_1+\stackrel{}{p}_2)}\left[e^{w_0(|p_1|+|p_2|)}+c_1e^{|p_{12}|w_0}+c_2e^{|p_{12}|w_0}\right]`$ (36)
Parameter $`c_2`$ should vanish, $`c_2=0`$, to prevent growth inside the bulk. Near the boundary, as $`w_00`$, $`A_\mathrm{\Delta }(w)w_0^\mathrm{\Delta }`$. Therefore for $`\mathrm{\Delta }=\mathrm{\Delta }_+=\mathrm{\Delta }_{}+1`$ we need $`c_1=1`$ so that the asymptotics of two terms in square brackets cancel each other at $`w_0=0`$. For $`\mathrm{\Delta }=\mathrm{\Delta }_{}`$ asymptotic itself is correct for any $`c_11`$, and is not enough to choose the right solution. However, as $`w_00`$ the $`\stackrel{}{w}`$-Fourier transform of $`A_\mathrm{\Delta }_{}`$ should be symmetric in all the three momenta, provided that $`K_\mathrm{\Delta }_{}(w|\stackrel{}{p}_1)`$ and $`K_\mathrm{\Delta }_{}(w|\stackrel{}{p}_2)`$ are chosen for the external legs:
$$S(\stackrel{}{p}_1,\stackrel{}{p}_2,\stackrel{}{p}_3)\underset{w_00}{lim}w_0^\mathrm{\Delta }_{}A_\mathrm{\Delta }_{}(w_0,\stackrel{}{w}|\stackrel{}{p}_1,\stackrel{}{p}_2)e^{i\stackrel{}{w}\stackrel{}{p}_3}๐\stackrel{}{w}=$$
$$=\frac{\delta ^{(d)}(\stackrel{}{p}_1+\stackrel{}{p}_2+\stackrel{}{p}_3)}{|p_1||p_2||p_3|(|p_1|+|p_2|+|p_3|)}\times \frac{1+c_1}{|p_3||p_1||p_2|}$$
$`\delta ^{(d)}`$-function allows to substitute $`|p_3|=\sqrt{\stackrel{}{p}_3^2}`$ instead of $`|p_{12}|=\sqrt{(\stackrel{}{p}_1+\stackrel{}{p}_2)^2}`$. The last ratio breaks the symmetry unless $`c_1=\frac{|p_1|+|p_2|}{|p_{12}|}H(\stackrel{}{p}_1,\stackrel{}{p}_2,\stackrel{}{p}_3)`$. Here $`H(\stackrel{}{p}_1,\stackrel{}{p}_2,\stackrel{}{p}_3)`$ \- is some symmetric function of its arguments. To find this function we need to fix the asymptotic of $`A_\mathrm{\Delta }(w|\stackrel{}{p}_1,\stackrel{}{p}_2)Cw_0^\mathrm{\Delta }`$ at $`w_00`$. This can be done by direct evaluation of integral in (21) and in appendix C for $`w_0=0`$. The result is $`H(\stackrel{}{p}_1,\stackrel{}{p}_2,\stackrel{}{p}_3)=1`$. Thus, $`c_1=\frac{|p_1|+|p_2|}{|p_{12}|}`$. With such choice of $`c_1`$, for $`\mathrm{\Delta }=\mathrm{\Delta }_\pm `$, the pole at $`|p_{12}|=|p_1|+|p_2|`$, i.e. at collinear momenta $`\stackrel{}{p}_1||\stackrel{}{p}_2`$ disappears from the amplitude (but only at $`w_0=0`$).
## 8 Appendix C. Single-vertex diagram in conventional AdS theory
It is instructive to repeat the calculation from s.4.1 for conventional AdS theory, with the same fields of dimensions $`\mathrm{\Delta }_\pm `$, but with the physical bulk-to-bulk propagator
$`๐ข(w,z)=\left({\displaystyle \frac{w_0z_0}{(w_0z_0)^2+(\stackrel{}{w}\stackrel{}{z})^2}}\right)^\mathrm{\Delta }_{}=(w_0z_0)^\mathrm{\Delta }_{}{\displaystyle \frac{d^d\stackrel{}{p}}{|p|}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}e^{|p||w_0z_0|}}`$ (37)
instead of our projected $`G_0(w,z)`$ from (18). The technical difference is that $`w_0z_0`$ can change sign and thus the integral over $`z_0`$ in the analogue of (21) is more sophisticated:
$$๐(w|\stackrel{}{p}_1,\mathrm{},\stackrel{}{p}_n)=\frac{dz_0d^d\stackrel{}{z}}{z_0^{d+1}}๐ข(w,z)K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_1)\mathrm{}K_{\mathrm{\Delta }_+}(z|\stackrel{}{p}_n)\stackrel{(\text{37})\&(\text{12})}{=}$$
$$=\frac{dz_0d^d\stackrel{}{z}}{z_0^{d+1}}\left((w_0z_0)^\mathrm{\Delta }_{}\frac{d\stackrel{}{p}}{|p|}e^{i\stackrel{}{p}(\stackrel{}{w}\stackrel{}{z})}e^{|p||w_0z_0|}\right)\left(z_0^\mathrm{\Delta }_{}e^{i\stackrel{}{p}_1\stackrel{}{z}}e^{|p_1|z_0}\right)\mathrm{}\left(z_0^\mathrm{\Delta }_{}e^{i\stackrel{}{p}_n\stackrel{}{z}}e^{|p_n|z_0}\right)=$$
$$=\frac{w_0^\mathrm{\Delta }_{}}{|p_{1n}|}e^{i\stackrel{}{w}(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)}(e^{|p_{1n}|w_0}_0^{w_0}e^{(|p_1|+\mathrm{}+|p_n||p_{1n}|)z_0}z_0^{s_{n+1}1}dz_0+$$
$$+e^{|p_{1n}|w_0}_{w_0}^{\mathrm{}}e^{(|p_1|+\mathrm{}+|p_n|+|p_{1n}|)z_0}z_0^{s_{n+1}1}dz_0)=$$
$$=\frac{w_0^\mathrm{\Delta }_{}}{|p_{1n}|}e^{i\stackrel{}{w}(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)}\{\frac{\mathrm{\Gamma }(s_{n+1})}{(|p_1|+\mathrm{}+|p_n||p_{1n}|)^{s_{n+1}}}e^{|p_{1n}|w_0}+$$
$$\underset{k=0}{\overset{s_{n+1}1}{}}\frac{\mathrm{\Gamma }(s_{n+1})}{\mathrm{\Gamma }(s_{n+1}k)}w_0^{s_{n+1}k1}(\frac{1}{(|p_1|+\mathrm{}+|p_n|+|p_{1n}|)^{k+1}}\frac{1}{(|p_1|+\mathrm{}+|p_n||p_{1n}|)^{k+1}})e^{(|p_1|+\mathrm{}+|p_n|)w_0}\}$$
Here $`|p_{1n}|\sqrt{(\stackrel{}{p}_1+\mathrm{}+\stackrel{}{p}_n)^2}`$. Note that this expression has poles at $`|p_{1n}|=|p_1|+\mathrm{}+|p_n|`$, i.e. when all the $`n`$ momenta $`\stackrel{}{p}_1,\mathrm{}\stackrel{}{p}_n`$ are collinear.
This complicated formula is typical for the amplitudes in conventional AdS theory . In this case it is somewhat tedious, though possible, to perform the recursion from section 4.2 and evaluate multi-vertex diagrams. Even more difficult is to find an effective boundary theory and the concept of ordinary AdS/CFT correspondence and/or open-closed string duality remains obscure, at least at the level of explicit multi-vertex calculations.
## 9 Acknowledgements
It is a pleasure to thank A.Alexandrov, V.Alexandrov, N.Amburg, S.Demidov, D.Gorbunov, A.Gorsky, M.Konushikhin, D.Levkov, M.Libanov, A.Losev, V.Lysov, A.Mironov, M.Rotaev, V.Rubakov, G.Rubtsov, S.Sibiryakov, T.Tomaras and M.Vasiliev for useful discussions.
This work is partly supported by the Federal Program of the Russian Ministry for Industry, Science and Technology No 40.052.1.1.1112 and the RFBR grant 04-02-16880. The work of D.K. is also supported by the RFBR grant 05-02-17363 and the studentship of Dynasty Foundation. |
warning/0506/math0506159.html | ar5iv | text | # Vector partition function and representation theory
## 1. Introduction
In this note, we are interested in the two following computational problems for classical Lie algebras $`A_r`$, $`B_r`$, $`C_r`$, $`D_r`$:
* The multiplicity $`c_\lambda ^\mu `$ of the weight $`\mu `$ in the representation $`V(\lambda )`$ of highest weight $`\lambda `$.
* Littlewood-Richardson coefficients, that is the multiplicity $`c_{\lambda \mu }^\nu `$ of the representation $`V(\nu )`$ in the tensor product of representations of highest weights $`\lambda `$ and $`\mu `$.
Softwares LฤฑE (from van Leeuwen et al. \[vL94\]) and GAP \[GAP\]), and Maple packages coxeter/weyl (from Stembridge \[S95\]), use Freudenthalโs and Klimykโs formulae, and work for any semi-simple Lie algebra (not only for classical Lie algebras). Unfortunately, these formulae are really sensitive to the size of coefficients of weights. Moreover, they do not lead to the computation of associated quasipolynomials $`(\lambda ,\mu )c_\lambda ^\mu `$ and $`(\lambda ,\mu ,\nu )c_{\lambda \mu }^\nu `$.
Here the approach to these two problems is through *vector partition function*, that is the function computing the number of ways one can decompose a vector as a linear combination with nonnegative integral coefficients of a fixed set of vectors. For example the number $`p(x)`$ of ways of counting $`x`$ euros with coins, that is
$$p(x)=\mathrm{}\{n_+^8;x=200n_1+100n_2+50n_3+20n_4+10n_5+5n_6+2n_7+n_8\},$$
can be seen as the partition of the $`1`$-dimensional vector $`(x)`$ with respects to the set $`\{(200),(100),(50),(20),(10),(5),(2),(1)\}`$ of 1-dimensional vectors. In the case of the decomposition with respects to the set of positive roots of a simple Lie algebra, we speak of *Kostant partition function*.
Recall that any $`d`$-dimensional rational convex polytope can be written as the set $`P(\mathrm{\Phi },a)`$ of nonnegative solutions $`x=(x_i)^N`$ of an equation $`_{i=1}^Nx_i\varphi _i=a`$, for a matrix $`\mathrm{\Phi }`$ with columns $`\varphi _i^r`$ and $`a^r`$ ($`d=Nr`$). It follows that evaluating the vector partition is equivalent to computing the number of integral points in a rational convex polytope.
The vector partition function arises in many areas of mathematics: representation theory, flows in networks, magic squares, statistics, crystal bases of quantum groups. Its complexity is polynomial in the size of input when the dimension of the polytope is fixed, and NP-hard if it can vary \[B94, B97, BP99\].
There are several approaches to the vector partition problem. For example Barvinokโs decomposition algorithm \[B94\], recently implemented by the LattE team \[DHTY03, L\], works for general sets of vectors. Beck-Pixton \[BP03\] also created an algorithm dedicated to the vector set arising from the Birkhoff polytope, counting the number of semi-magic squares.
In this note, we use recent results of Baldoni-Beck-Cochet-Vergne \[BBCV05\] to obtain a fast algorithm for Kostant partition function via inverse Laplace formula. These results involve DeConcini-Procesiโs *maximal nested sets* (or in short MNSs \[DCP04\]) and iterated residues of rational functions computed by formal power series development.
We combine resulting procedures with Kostantโs and Steinbergโs formulae giving $`c_\lambda ^\mu `$ and $`c_{\lambda \mu }^\nu `$ in terms of vector partition function. We then obtain a Maple program computing for classical Lie algebras ($`A_r`$, $`B_r`$, $`C_r`$, $`D_r`$), the multiplicity of a weight in an irreducible finite-dimensional representation, as well as decomposition coefficients of the tensor product of two irreducible finite-dimensional representations. To the best of our knowledge, they are also the only ones able to compute associated piecewise-defined quasipolynomials $`(\lambda ,\mu )c_\lambda ^\mu `$ and $`(\lambda ,\mu ,\nu )c_{\lambda \mu }^\nu `$.
These programs (available at \[C\]) are specially designed for large parameters of weights. Indeed although only written in Maple they can perform examples with weights with 5 digits coordinates, far beyond classical softwares written in C++. We also stress that our programs are absolutely clear, easy to use and require no installation of exotic package or program. Retro-compatibility has been checked downto Maple Vr5. They are fully commented, so that a curious user can figure out their internal mechanisms.
However, certain other softwares and packages are not limited by the rank of the algebra like our programs. For example computation of non-trivial examples in Lie algebras of rank $`10`$ is possible with the software LฤฑE, whereas our programs are efficient up to rank $`5`$$`7`$. These facts make our programs complementary to traditional softwares.
Remark that Kostantโs and Steinbergโs formulae have already been implemented once in the case of $`A_r`$ \[C03\]. This previous program relies on results of Baldoni-Vergne \[BV01\] implemented by Baldoni-DeLoera-Vergne \[BdLV03\], computing Kostant partition function only in the case of $`A_r`$. Tools were *special permutations* and again iterated residues of rational fraction.
A new technique for Littlewood-Richardson coefficients has been recently designed by DeLoera-McAllister \[DM05\]. For $`A_r`$, they wrote an algorithm using hive polytopes \[KT99\]. For $`B_r`$, $`C_r`$, $`D_r`$, they implemented Berenstein-Zelevinsky polytopes \[BZ01\]. They can also evaluate stretched Littlewood-Richardson coefficients $`c_{t\lambda t\mu }^{t\nu }`$. These two methods consist in computing a tensor product coefficient as the number of lattice points in just one specific convex rational polytope. However our programs based on multidimensional residues are faster, and can reach examples not available by their method.
This paper is organized as follows. Section 2 recalls representation theory problems we are interested in and links them with algebraic combinatorics. Section 3 describes more precisely rational convex polytopes and formulae counting their integral points. Section 4 introduces maximal nested sets and formulae that were used in our programs. Finally in Section 5 we perform tests of our programs.
## 2. Representation theory and convex polytopes
Let us fix the notations once and for all. Let $`๐ค`$ be a semi-simple Lie algebra of rank $`r`$. Choose a Cartan subalgebra $`๐ฑ`$ of $`๐ค`$ and denote by $`L๐ฑ^{}`$ the weight lattice.
Let $`\mathrm{\Delta }^+`$ be a positive roots system. The root lattice is defined as $`[\mathrm{\Delta }^+]`$. Let $`C(\mathrm{\Delta }^+)`$ be the cone spanned by linear combinations with nonnegative coefficients of positive roots. The Weyl group of $`๐ค`$ for $`๐ฑ`$ is denoted by $`W`$.
There exist only four simple Lie algebras $`A_r`$, $`B_r`$, $`C_r`$, $`D_r`$ of rank $`r`$, called *classical Lie algebras* of rank $`r`$ \[Bou68\], and determined by their positive roots systems:
$`A_r:`$ $`\mathrm{\Delta }^+=\{e_ie_j|\mathrm{\hspace{0.17em}1}i<jr+1\}^{r+1},`$
$`B_r:`$ $`\mathrm{\Delta }^+=\{e_ie_j|\mathrm{\hspace{0.17em}1}i<jr\}\{e_i|\mathrm{\hspace{0.17em}1}ir\}^r,`$
$`C_r:`$ $`\mathrm{\Delta }^+=\{e_ie_j|\mathrm{\hspace{0.17em}1}i<jr\}\{2e_i|\mathrm{\hspace{0.17em}1}ir\}^r,`$
$`D_r:`$ $`\mathrm{\Delta }^+=\{e_ie_j|\mathrm{\hspace{0.17em}1}i<jr\}\{e_i+e_j|\mathrm{\hspace{0.17em}1}i<jr\}^r.`$
The character of a representation $`V`$ of $`๐ค`$ is $`\mathrm{ch}(V)=_{\mu L}dim(V_\mu )e^\mu `$. Recall that the irreducible finite-dimensional representation of $`๐ค`$ of highest weight $`\lambda `$ is denoted by $`V(\lambda )`$. Hence the weight multiplicity $`c_\lambda ^\mu `$ is defined as $`dim(V(\lambda )_\mu )`$ for any weight $`\mu `$ such that $`\lambda \mu `$ is in the root lattice. Multiplicities $`c_\lambda ^\mu `$ are called Kostka numbers when $`๐ค=A_r=๐ฐ๐ฉ_{r+1}()`$.
On the other hand, multiplicities of representations $`V(\nu )`$ in the tensor product $`V(\lambda )V(\mu )`$ are called Littlewood-Richardson coefficients (or Clebsch-Gordan coefficients). Here $`\nu `$ is a dominant weight such that $`\lambda +\mu \nu `$ is in the root lattice.
Evaluating weight multiplicities and Littlewood-Richardson coefficients is a difficult task. For $`A_1`$, computing Kostka numbers is immediate and Clebsch-Gordanโs formula gives Littlewood-Richardson coefficients. For $`A_2`$, one can still compute some small examples. But for general $`X_r`$ ($`r3`$) or for weights which components are big (say, with two digits), direct computation is usually intractable.
There exist many formulae from representation theory for $`c_\lambda ^\mu `$ and $`c_{\lambda \mu }^\nu `$. The first one, valid in any complex semi-simple Lie algebra $`๐ค`$, is Weylโs character formula
$`\mathrm{ch}(V(\lambda ))`$ $`=`$ $`{\displaystyle \frac{A_{\lambda +\rho }}{A_\rho }},\text{ where }A_\mu ={\displaystyle \underset{wW}{}}(1)^{\epsilon (w)}e^{w(\mu )},`$
where $`\rho `$ is half the sum of positive roots for $`๐ค`$. Littlewood-Richardson coefficients are obtained from this formula, since the character of $`V(\lambda )V(\mu )`$ is
$$\mathrm{ch}(V(\lambda )V(\mu ))=\mathrm{ch}(V(\lambda ))\times \mathrm{ch}(V(\mu ))=\underset{\nu L;\lambda +\mu \nu [\mathrm{\Delta }^+]}{}c_{\lambda \mu }^\nu \mathrm{ch}(V(\nu )).$$
But these two formulae do not lead to efficient computations when the rank of $`๐ค`$ or the size of coefficients of weights grow. Moreover, computing the whole character is untractable: for $`๐ค=A_3=๐ฐ๐ฉ_4()`$ and $`\lambda =(2,1,0,3)`$, the character $`\mathrm{ch}(V(\lambda ))`$ has 9 monomials but the character $`\mathrm{ch}(V(10\lambda ))`$ has 2903 monomials.
Let us describe Kostantโs and Steinbergโs formulae in the case of any semi-simple Lie algebra $`๐ค`$. Denote by $`k_๐ค(a)`$ the number of ways one can write a vector $`a`$ as a nonnegative linear combination of positive roots. Remark that $`k_๐ค(a)=0`$ unless $`a`$ is in the root lattice $`[\mathrm{\Delta }^+]`$. This number satisfies the equation
$`{\displaystyle \frac{1}{_{\alpha \mathrm{\Delta }^+}(1e^\alpha )}}={\displaystyle \underset{a[\mathrm{\Delta }^+]}{}}k_๐ค(a)e^a.`$
Let $`\lambda `$ and $`\mu `$ be respectively a dominant weight and a weight such that $`\lambda \mu [\mathrm{\Delta }^+]`$. A Weyl group element $`wW`$ is valid for $`\lambda `$ and $`\mu `$ if the root lattice element $`w(\lambda +\rho )(\mu +\rho )`$ is in the cone $`C(\mathrm{\Delta }^+)`$. The set of such $`w`$โs is denoted by $`\mathrm{Val}(\lambda ,\mu )`$. Then Kostantโs formula asserts that the weight multiplicity $`c_\lambda ^\mu `$ equals
(2.1)
$$c_\lambda ^\mu =\underset{w\mathrm{Val}(\lambda ,\mu )}{}(1)^{\epsilon (w)}k_๐ค(w(\lambda +\rho )(\mu +\rho )).$$
Similarly let $`\lambda `$, $`\mu `$, $`\nu `$, be three dominant weights such that $`\lambda +\mu \nu [\mathrm{\Delta }^+]`$. The couple $`(w,w^{})W\times W`$ is valid for $`\lambda `$, $`\mu `$, $`\nu `$, if the root lattice element $`w(\lambda +\rho )+w^{}(\mu +\rho )(\nu +2\rho )`$ is in $`C(\mathrm{\Delta }^+)`$. The set of such couples is denoted by $`\mathrm{Val}(\lambda ,\mu ,\nu )`$. Then Steinbergโs formula asserts that the Littlewood-Richardson coefficient equals
(2.2)
$$c_{\lambda \mu }^\nu =\underset{(w,w^{})\mathrm{Val}(\lambda ,\mu ,\nu )}{}(1)^{\epsilon (w)+\epsilon (w^{})}k_๐ค(w(\lambda +\rho )+w^{}(\mu +\rho )(\nu +2\rho )).$$
Sets of valid Weyl group elements and valid couples of Weyl group elements turn out to be relatively small, when compared to $`W`$ and $`W\times W`$ (which size is exponential in the rank). Remark that Kostantโs (resp. Steinbergโs) formula also work when $`\lambda \mu `$ (resp. $`\lambda +\mu \nu `$) is not in the root lattice, since Kostant partition function vanishes on vectors that are not in the root lattice.
From now on, let $`X_r`$ be a classical Lie algebra of rank $`r`$. Here $`X`$ stands for $`A`$, $`B`$, $`C`$, $`D`$. Its positive roots system will be denoted by $`X_r^+`$.
Multiplicities $`c_\lambda ^\mu `$ and $`c_{\lambda \mu }^\nu `$ behave nicely, in function of the parameters. More precisely, there exists a decomposition of the space $`๐ฑ^{}๐ฑ^{}๐ฑ^{}`$ in union of closed cones $`C`$, such that the restriction of $`c_{\lambda \mu }^\nu `$ to each cone $`C`$ is given by a quasi-polynomial function. This follows from theorems of Knutson-Tao \[KT99\] (for $`A_r`$), Berenstein-Zelevinsky \[BZ01\] (for any semi-simple Lie algebra) giving $`c_{\lambda \mu }^\nu `$ as the number of points in a rational convex polytope. In the case of $`A_r`$, the fact that $`c_{\lambda \mu }^\nu `$ is given on each cone $`C`$ by a polynomial function is proven in Rassart \[Ras04\], and the case of $`A_3`$ is treated as an illustration. The description of the decomposition of $`๐ฑ^{}๐ฑ^{}`$ in cones $`C`$, where the function $`c_\lambda ^\mu `$ is polynomial for $`A_r`$, was given for low ranks by Billey-Guillemin-Rassart \[BGR03\]. See also Rassartโs website \[R\] for wonderful slides.
The common point to Kostantโs and Steinbergโs formulae is the function counting the number of decompositions of a root lattice element as a linear combination with nonnegative integral coefficients of positive roots of the Lie algebra. The next section deals with an efficient method to compute it.
## 3. Counting integral points in rational convex polytopes
### 3.1. Vector partition function
Let $`E^r`$ and $`\mathrm{\Phi }`$ be an integral matrix with set of columns $`\mathrm{\Delta }^+=\{\varphi _1,\mathrm{},\varphi _N\}E^{}`$. Choose $`a^r`$. The rational convex polyhedron associated to $`\mathrm{\Phi }`$ and $`a`$ is
$`P(\mathrm{\Phi },a)`$ $`=`$ $`\left\{x^N;{\displaystyle \underset{i=1}{\overset{N}{}}}x_i\varphi _i=a,x_i0\right\}.`$
###### Remark 3.1.
Every convex polyhedron can be realized under the form $`P(\mathrm{\Phi },a)`$, that is as a set satisfying equality constraints on nonnegative variables. Indeed any inequality can be replaced by an equality by adding a new variable. For example polytopes $`\{(x,y)^2;x0,y0,x+y1\}`$ and $`\{(x,y,z)^3;x0,y0,z0,x+y+z=1\}`$ are isomorphic and have the same number of integral points.
We assume that $`a`$ is in the cone $`C(\mathrm{\Phi })`$ spanned by nonnegative linear combinations of the vectors $`\varphi _i`$, so that $`P(\mathrm{\Phi },a)`$ in non-empty. We also assume that the kernel of $`\mathrm{\Phi }`$ intersects trivially with the positive orthant $`_+^N`$, so that the cone $`C(\mathrm{\Phi })`$ is acute and $`P(\mathrm{\Phi },a)`$ is a polytope (i.e. bounded). Finally, we assume that $`\mathrm{\Phi }`$ has rank $`r`$. The vector partition function is by definition
$`k(\mathrm{\Phi },a)`$ $`=`$ $`\left|P(\mathrm{\Phi },a)_+^N\right|,`$
that is the number of nonnegative integral solutions $`(x_1,\mathrm{},x_N)`$ of the equation $`_{i=1}^Nx_i\varphi _i=a`$. If $`\mathrm{\Phi }=\mathrm{\Phi }(X_r)`$ is the matrix which columns are positive roots for a classical Lie algebra $`X_r`$, then $`ak(\mathrm{\Phi }(X_r),a)`$ is the Kostant partition function. For example
$$\mathrm{\Phi }(A_2)=\left(\begin{array}{ccc}1& 1& 0\\ 1& 0& 1\\ 0& 1& 1\end{array}\right)\text{and}\mathrm{\Phi }(B_2)=\left(\begin{array}{cccc}1& 1& 1& 0\\ 1& 1& 0& 1\end{array}\right).$$
Note that the matrix $`\mathrm{\Phi }(A_r)`$ has rank $`r`$ (and not $`r+1`$), since sums on lines are zero.
A *basic subset* of $`\mathrm{\Delta }^+`$ is a basis $`\sigma =\{\alpha _1,\mathrm{},\alpha _r\}`$ of $`E^{}`$ constituted with elements of $`\mathrm{\Delta }^+`$. Let $`B(\mathrm{\Delta }^+)`$ be the collection of all basic subsets of $`\mathrm{\Delta }^+`$. For such a $`\sigma `$, let $`C(\sigma )`$ be the cone of linear combinations with nonnegative coefficients of $`\alpha _i`$โs. Denote by $`\mathrm{Sing}(\mathrm{\Delta }^+)`$ the reunion of the facets of cones $`C(\sigma )`$, $`\sigma B(\mathrm{\Delta }^+)`$; this is the set of *singular* vectors. Let $`C_{\mathrm{reg}}(\mathrm{\Delta }^+):=C(\mathrm{\Delta }^+)\mathrm{Sing}(\mathrm{\Delta }^+)`$ be the set of *regular* vectors. A *combinatorial chamber* $`๐ `$ is by definition a connected component of $`C_{\mathrm{reg}}(\mathrm{\Delta }^+)`$. Combinatorial chambers are regions of quasi-polynomiality of the vector partition function $`ak(\mathrm{\Phi },a)`$. Figure 1 represents cones $`C(A_3^+)`$ and $`C(B_3^+)`$, and their chamber decompositions.
### 3.2. Brion-Szenes-Vergne formula for classical Lie algebras
Let us describe the formula, computing the number of integral points in rational convex polytopes $`P(\mathrm{\Phi }(X_r),a)`$ associated to a classical algebra $`X_r`$, that was implemented in our program.
Let $`E=๐ฑ`$ and consider the set $`\mathrm{\Delta }^+`$ of positive roots for $`X_r`$. Denote by $`\mathrm{\Delta }`$ the set $`\mathrm{\Delta }^+(\mathrm{\Delta }^+)`$ of all roots. Let $`R_\mathrm{\Delta }`$ be the vector space of fractions with poles on the hyperplanes defined as kernels of forms $`\alpha \mathrm{\Delta }`$. Let $`S_\mathrm{\Delta }`$ be the vector space generated by fractions $`f_\sigma :=\frac{1}{_{\alpha \sigma }\alpha }`$, $`\sigma B(\mathrm{\Delta }^+)`$. Brion-Vergne \[BV97\] proved that $`R_\mathrm{\Delta }`$ decomposes as the direct sum $`S_\mathrm{\Delta }(R_\mathrm{\Delta })`$. We define the *Jeffrey-Kirwan residue* of the chamber $`๐ `$ as the linear form $`\mathrm{JK}_๐ `$ on $`S_\mathrm{\Delta }`$:
$$\mathrm{JK}_๐ (f_\sigma ):=\{\begin{array}{cc}\mathrm{vol}(\sigma )^1,\hfill & \text{if }๐ C(\sigma ),\hfill \\ 0,\hfill & \text{if }๐ C(\sigma )=\mathrm{},\hfill \end{array}$$
where $`\mathrm{vol}(\sigma )`$ is the volume of the parallelopiped $`_{\alpha \sigma }[0,1]\alpha `$. We extend the JK residue to a linear form on $`R_\mathrm{\Delta }`$ by setting it to $`0`$ on $`(R_\mathrm{\Delta })`$, and to a linear form on the space of formal series $`\widehat{R_\mathrm{\Delta }}`$ by setting it to $`0`$ on homogeneous elements of degree different from $`r`$. For example, for the system $`\mathrm{\Delta }^+=\{e_1,e_2,e_1+e_2,e_1e_2\}^2`$ of positive roots for $`B_2`$ and the chamber $`๐ =_+e_1_+(e_1+e_2)`$ we have
$$\mathrm{JK}_๐ \left(\frac{e^{xy}}{xy^2}\right)=\mathrm{JK}_๐ \left(\frac{xy}{xy^2}\right)=\mathrm{JK}_๐ \left(\frac{1}{y^2}\frac{1}{xy}\right)=1,$$
since $`๐ C(\{e_1,e_2\})`$.
Let $`T`$ be the torus $`E/E_{}`$, where $`E_{}E`$ is the dual of the root lattice. Given a basic subset $`\sigma `$, we define $`T(\sigma )`$ as the set of elements $`gT`$ such that $`e^{\alpha ,2i\pi G}=1`$ for all $`\alpha \sigma `$; here $`G`$ is a representative of $`gE/E_{}`$. Now let
$$(g,a)(u):=\frac{e^{a,2i\pi G+u}}{_{\alpha \mathrm{\Delta }}(1e^{\alpha ,2i\pi G+u})}.$$
###### Theorem 3.2 (Brion-Szenes-Vergne \[BV99, SV04\]).
Let $`FT`$ be a finite set such that $`T(\sigma )F`$ for all $`\sigma B(\mathrm{\Delta }^+)`$. Fix a combinatorial chamber $`๐ `$. Then for all $`a[\mathrm{\Delta }^+]\overline{๐ }`$, we have:
$$k(\mathrm{\Phi },a)=\underset{gF}{}\mathrm{JK}_๐ ((g,a)).$$
Now that we linked vector partition function and Jeffrey-Kirwan residue, we describe in Section 4 an efficient way to compute the latter.
## 4. DeConcini-Procesiโs maximal nested sets (MNS) \[DCP04\]
We keep the same notations as in Section 3. A subset $`S\mathrm{\Delta }^+`$ is *complete* if $`S=S\mathrm{\Delta }^+`$. A complete subset is *reducible* if one can find a decomposition $`E=E_1E_2`$ such that $`S=S_1S_2`$ with $`S_1E_1`$ and $`S_2E_2`$; else $`S`$ is said *irreducible*. Let $``$ be the collection of irreducible subsets.
A collection $`M=\{I_1,I_2,\mathrm{},I_s\}`$ of irreducible subsets $`I_j`$ of $`\mathrm{\Delta }^+`$ is *nested*, if: for every subset $`\{S_1,\mathrm{},S_m\}`$ of $`M`$ such that there exist no $`i`$, $`j`$ with $`S_iS_j`$, the union $`S_1S_2\mathrm{}S_m`$ is complete and the $`S_i`$โs are its irreducible components. Note that a maximal nested set (MNS in short) has exactly $`r`$ elements.
Assume $`\mathrm{\Delta }^+`$ irreductible and fix a total order on it. For $`M=\{I_1,\mathrm{},I_s\}`$, $`I_j\mathrm{\Delta }^+`$, take for every $`j`$ the maximal element $`\beta _jI_j`$. This defines an application $`\varphi (M):=\{\beta _1,\mathrm{},\beta _s\}\mathrm{\Delta }^+`$. A maximal nested set $`M`$ is *proper* if $`\varphi (M)`$ is a basis of $`E^{}`$. Denote by $`๐ซ`$ the collection of maximal proper nested sets (MPNS in short). We sort $`\varphi (M)`$ and get an ordered list $`\theta (M)=[\alpha _1,\mathrm{},\alpha _r]`$. Thus $`\theta `$ is an application from the collection of MPNSs to the collection of ordered basis of $`E^{}`$. For a given $`M`$, let then
$`C(M)`$ $`:=`$ $`C(\alpha _1,\mathrm{},\alpha _r),`$
$`\mathrm{vol}(M)`$ $`:=`$ $`\mathrm{vol}\left(_{i=1}^r[0,1]\alpha _i\right),`$
$`\mathrm{IRes}_M`$ $`:=`$ $`\mathrm{Res}_{\alpha _r=0}\mathrm{}\mathrm{Res}_{\alpha _1=0}.`$
###### Example 4.1.
Let $`e_i`$ be the canonical basis of $`^r`$, with dual basis $`e^i`$ ($`i=1`$, โฆ, $`r`$), and define $`E`$ as the subspace of vectors which sum of coordinates vanish. Consider the set $`\mathrm{\Delta }^+=\{e^ie^j|\mathrm{\hspace{0.17em}1}i<jr\}`$ of positive roots for $`A_{r1}`$. Irreducible subsets of $`\mathrm{\Delta }^+`$ are indexed by subsets $`S`$ of $`\{1,2,\mathrm{},r\}`$, the corresponding irreducible subset being $`\{e^ie^j|i,jS,i<j\}`$. For instance $`S=\{1,2,4\}`$ parametrizes the set of roots given by $`\{e^1e^2,e^2e^4,e^1e^4\}`$.
A nested set is represented by a collection $`M=\{S_1,S_2,\mathrm{},S_k\}`$ of subsets of $`\{1,2,\mathrm{},r\}`$ such that if $`S_i`$, $`S_jM`$ then either $`S_iS_j`$ is empty, or one of them is contained in another.
For example one can easily compute that for the set of positive roots for $`A_3`$ (see Figure 1) there are only 7 MPNS, namely
$$\begin{array}{cc}M_1=\{[1,2],[1,2,3],[1,2,3,4]\},\hfill & M_2=\{[2,3],[1,2,3],[1,2,3,4]\},\hfill \\ M_3=\{[2,3],[2,3,4],[1,2,3,4]\},\hfill & M_4=\{[3,4],[2,3,4],[1,2,3,4]\},\hfill \\ M_5=\{[1,3],[2,4],[1,2,3,4]\},\hfill & M_6=\{[1,2],[3,4],[1,2,3,4]\}.\hfill \end{array}$$
Now we can quote the Theorem for the Jeffrey-Kirwan residue computation:
###### Theorem 4.2 (DeConcini-Procesi).
Let $`๐ `$ be a combinatorial chamber and fix $`fR_\mathrm{\Delta }`$. Take any regular vector $`v๐ `$. Then:
$$\mathrm{JK}_๐ (f)=\underset{M๐ซ:vC(M)}{}\frac{1}{\mathrm{vol}(M)}\mathrm{IRes}_M(f).$$
See \[BBCV05\] for a detailed description of how formulae from Theorems 3.2 and 4.2 were implemented.
## 5. Our programs
### 5.1. Description and implementation
Initial data for weight multiplicity and Littlewood-Richardson coefficients are only vectors (respectively two and three). Our programs work with weights represented in the canonical basis of $`E^{}`$, and not in the fundamental weights basis for $`X_r`$. Translation between these two bases is performed via straightforward procedures FromFundaToCanoX(r,vโ) and FromCanoToFundaX(r,v) (where one replaces X by A, B, C, D, according to the algebra).
Computation of the weight multiplicity $`c_\lambda ^\mu `$ and of the Littlewood-Richardson coefficient $`c_{\lambda \mu }^\nu `$ is done by typing in
MultiplicityX(lambda,mu);
TensorProductX(lambda,mu,nu);
where $`\lambda `$, $`\mu `$, $`\nu `$ are suitable weights. The syntax for computing quasipolynomials is slightly different. Assume that we want to evaluate $`(\lambda ^{},\mu ^{})c_\lambda ^{}^\mu ^{}`$ in a neighborhood of a couple $`(\lambda ,\mu )`$, and $`(\lambda ^{},\mu ^{},\nu ^{})c_{\lambda ^{}\mu ^{}}^\nu ^{}`$ in a neighborhood of a triple $`(\lambda ,\mu ,\nu )`$. Let $`\lambda _F=[x_1,\mathrm{},x_r]`$, $`\mu _F=[y_1,\mathrm{},y_r]`$, $`\nu _F=[z_1,\mathrm{},z_r]`$, be three formal vectors where $`x_i`$โs, $`y_i`$โs and $`z_i`$โs are variables. Then we use the command lines
$`\mathrm{๐ฟ๐๐๐ข๐๐๐๐๐๐๐ผ๐๐๐๐๐๐๐๐๐๐๐ข๐}(\mathrm{๐๐๐๐๐๐},\mathrm{๐๐๐๐๐๐๐ต},\mathrm{๐๐},\mathrm{๐๐๐ต});`$
$`\mathrm{๐ฟ๐๐๐ข๐๐๐๐๐๐๐๐๐๐๐๐๐ฟ๐๐๐๐๐๐๐}(\mathrm{๐๐๐๐๐๐},\mathrm{๐๐๐๐๐๐๐ต},\mathrm{๐๐},\mathrm{๐๐๐ต},\mathrm{๐๐},\mathrm{๐๐๐ต});`$
So for the polynomial $`(\lambda ^{},\mu ^{})c_\lambda ^{}^\mu ^{}`$ with $`\lambda =(3,2,1,6)`$ and $`\mu =(2,2,2,2)`$ for $`A_3`$ we enter
$`\mathrm{๐ฟ๐๐๐ข๐๐๐๐๐๐๐ผ๐๐๐๐๐๐๐๐๐๐๐ข๐ฐ}(`$
$`[\mathrm{๐น},\mathrm{๐ธ},\mathrm{๐ท},\mathrm{๐ผ}],[๐ก[\mathrm{๐ท}],๐ก[\mathrm{๐ธ}],๐ก[\mathrm{๐น}],๐ก[\mathrm{๐บ}]],[\mathrm{๐ธ},\mathrm{๐ธ},\mathrm{๐ธ},\mathrm{๐ธ}],[๐ข[\mathrm{๐ท}],๐ข[\mathrm{๐ธ}],๐ข[\mathrm{๐น}],๐ข[\mathrm{๐บ}]]);`$
and get instantly
$$\frac{1}{6}(3x_12y_1+1)(3x_12y_1+2)(3x_1+6x_22y_16y_2+3).$$
Remark that quasipolynomials $`c_{t\lambda }^{t\mu }`$ and $`c_{t\lambda t\mu }^{t\nu }`$ are obtained by setting $`x_i=t\lambda _i`$, $`y_i=t\mu _i`$, $`z_i=t\nu _i`$, so that
$`\mathrm{๐ฟ๐๐๐ข๐๐๐๐๐๐๐ผ๐๐๐๐๐๐๐๐๐๐๐ข๐ฐ}(`$
$`[\mathrm{๐น},\mathrm{๐ธ},\mathrm{๐ท},\mathrm{๐ผ}],[\mathrm{๐น}๐,\mathrm{๐ธ}๐,๐,\mathrm{๐ผ}๐],[\mathrm{๐ธ},\mathrm{๐ธ},\mathrm{๐ธ},\mathrm{๐ธ}],[\mathrm{๐ธ}๐,\mathrm{๐ธ}๐,\mathrm{๐ธ}๐,\mathrm{๐ธ}๐]);`$
returns $`(t+1)(t+2)(t+3)/6`$.
Now some words about implementation. There are two main parts in our programs. The first one is the implementation of Theorems 3.2 and 4.2; it is described in \[BBCV05\]. The second one is the implementation of Kostantโs (2.1) and Steinbergโs (2.1) formulae using valid Weyl group elements and valid couples of Weyl group elements; it is a generalization for classical Lie algebras of what has been done for $`A_r`$ in \[C03\].
### 5.2. Comparative tests
Figure 2 describes efficiency area of the software LฤฑE and of our programs using MNS; any area located to the left of a colored line represents the range where a program can compute examples in a reasonable time. Figures 45 present precise comparative tests of the software LฤฑE, of DeLoera-McAllisterโs script \[DM05\] using LattE \[L\] and of our programs using MNS.
All examples were runned on the same computer, a Pentium IV 1,13GHz with 2Go of RAM memory. Remark that computation times for LattE and LฤฑE are slower than those shown in \[DM05\], due to different computers. However, we performed exactly same examples for comparison purposes.
As in \[DM05\], in Tables 44 weights are for $`๐ค๐ฉ_{r+1}()`$ and not $`๐ฐ๐ฉ_{r+1}()`$ (coordinates do not add to zero). However the sum of coordinates of $`\lambda +\mu \nu `$ vanish.
Now some words about quasipolynomials computation. Let us examine the first example for $`B_3`$ in \[DM05\], that is the evaluation of the quasipolynomial $`c_{t\lambda t\mu }^{t\nu }`$ for weights $`\lambda =[0,15,5]`$, $`\mu =[12,15,3]`$ and $`\nu =[6,15,6]`$ expressed in the basis of fundamental weights. In canonical basis, these data become $`\lambda =(35/2,35/2,5/2)`$, $`\mu =(57/2,33/2,3/2)`$, $`\nu =(24,18,3)`$. The program using the MNS algorithm returns the quasipolynomial
$`c_{t\lambda t\mu }^{t\nu }`$ $`=`$ $`\left({\displaystyle \frac{203}{256}}+{\displaystyle \frac{53}{256}}(1)^t\right)+\left({\displaystyle \frac{1515}{128}}+{\displaystyle \frac{197}{128}}(1)^t\right)t`$
$`+\left({\displaystyle \frac{35353}{384}}+{\displaystyle \frac{881}{128}}(1)^t\right)t^2+\left({\displaystyle \frac{13405}{32}}\right)t^3`$
$`+\left({\displaystyle \frac{407513}{384}}\right)t^4+\left({\displaystyle \frac{68339}{64}}\right)t^5`$
in 1099,4s. On the other hand, the computation of the full quasipolynomial $`c_{\lambda \mu }^\nu `$ with formal vectors $`[x_1,x_2,x_3]`$, $`[y_1,y_2,y_3]`$, $`[z_1,z_2,z_3]`$ leads to a 87 pages result, obtained in only 1158,6s. With LattE, on our computer, one obtains the quasipolynomial $`c_{t\lambda t\mu }^{t\nu }`$ in only 825,8s.
As announced in the introduction, our program is really efficient for weights with huge coefficients. Note that in the particular case of $`A_r`$ the MNS algorithm allows us to compute examples one rank further than the $`\mathrm{Sp}(๐)`$ algorithm.
The translation of the program using MNS in the language of the symbolic calculation software MuPAD is in progress. A version using distributed calculation on a grid of computers is in the air; it will considerably increase the speed of computations. |
warning/0506/cond-mat0506153.html | ar5iv | text | # A single measurement of a quantum many-body system of bosons
## I Introduction
In order to simulate a single measurement of a position of a quantum particle described by a wave function $`\varphi (x)`$ it is enough to randomly draw a position $`x=\xi `$ with the probability density $`|\varphi (x)|^2`$. The outcome of the measurement, the particle at a point $`x=\xi `$, is different from the statistical average over many such measurements, given by $`|\varphi (x)|^2`$. This kind of difference, trivial in the case of measurements performed on a single particle, becomes especially interesting in the case of quantum systems that consist of many bosons. One prominent example is given by two colliding Bose-Einstein condensates Ketterle , where every single measurement of the system reveals an interference pattern while no pattern is present in the average over many measurements.
For a generic quantum state of a many-body system it is inexpensive to come up with a prediction for an average outcome of many measurements. On the other hand, finding out possible results of a single measurement is extremely difficult. If the state is given by a many-body wave function $`\varphi (x_1,\mathrm{},x_N)`$, where $`x_i`$ stands for a coordinate of an $`i`$-th particle, one needs to draw a set of $`N`$ positions (one for each particle) from the $`N`$-dimensional<sup>2</sup><sup>2</sup>2Here and in the rest of this note we consider one-dimensional particles probability density $`|\varphi (x_1,\mathrm{},x_N)|^2`$. For large number of particles, $`N`$, the direct sampling of the corresponding multidimensional probability density is very difficult, if possible at all. The direct sampling of the $`N`$-dimensional probability density can be replaced by sampling of $`N`$ one-dimensional conditional probability densities Javanainen . In practical applications, however, this clever method is suitable to handle large number of particles only if they occupy very few modes.
In the following I propose an approximate method of simulating outcomes of single measurements. It is designed for those quantum many-body states that involve many macroscopically occupied modes.
A reader, novice to the subject, may establish a necessary background by contemplating two first pages of Javanainen and three first pages of Rzazewski .
Suppose that a state $`|\psi `$ of a system of $`N`$ bosons is spanned on $`M`$ orthonormal modes $`u_i(x)`$, so that the bosonic field operator can be written as
$$\widehat{\mathrm{\Psi }}(x)=\underset{i=1}{\overset{M}{}}\widehat{a}_iu_i(x),$$
where $`\widehat{a}_i`$ annihilates a boson from the mode $`u_i(x)`$. The joint probability density for $`N`$ bosons in the state $`|\psi `$ is given by
$$p_{|\psi }^{(N)}=\frac{1}{N!}\psi |\widehat{\mathrm{\Psi }}^{}(x_1)\mathrm{}\widehat{\mathrm{\Psi }}^{}(x_N)\widehat{\mathrm{\Psi }}(x_N)\mathrm{}\widehat{\mathrm{\Psi }}(x_1)|\psi .$$
(1)
This is equivalent to the modulus squared of the $`N`$-body wave function. The task is to repeatedly generate a set of $`N`$ numbers $`\xi _1,\mathrm{},\xi _N`$ according to (1). As noticed by Javanainen and Yoo Javanainen , the probability density above can be decomposed into a product of one-dimensional conditional probabilities
$$p_{|\psi }^{(N)}=p_{|\psi }^{(1)}(x_1)p_{|\psi }^{(2)}(x_2|\xi _1)\mathrm{}p_{|\psi }^{(N)}(x_N|\xi _{N1},\mathrm{},\xi _1).$$
(2)
If one desires to generate a set of $`N`$ positions according to the probability $`p_{|\psi }^{(N)}`$, it is enough to first pick a position $`x_1=\xi _1`$ with the probability $`p_{|\varphi }^{(1)}(x_1)`$ then a position $`x_2=\xi _2`$ with the $`p_{|\varphi }^{(2)}(x_2|\xi _1)`$ and so on.
Each one-dimensional probability density of (2) takes the following functional form
$$p_{|\psi }^{(r)}(x)=\underset{i,j=1}{\overset{M}{}}c_{ij}^{(r)}u_i^{}(x)u_j(x),$$
(3)
where the coefficients $`c_{ij}^{(r)}=c_{ji}^{(r)}{}_{}{}^{}`$ are calculated from (1) with $`x_1=\xi _1,\mathrm{},x_{r1}=\xi _{r1}`$.
The important observation Javanainen ; Images is that for large $`N`$ the coefficients $`c_{ij}^{(r)}`$ assume approximately constant values after certain critical number $`N_{crit}`$ of positions have been drawn, i.e. $`c_{ij}^{(r)}=\text{const}_{ij}(r)`$ for $`r>N_{crit}`$. Assuming that $`N_{crit}N`$, these constant values vary from one position measurement of $`N`$ particles to another. In other words, first $`N_{crit}`$ particles determine the shape of the one-dimensional probability densities $`p_{|\psi }^{(r)}`$. The explanation of this โlocalizationโ of values of the coefficients is provided in the following.
## II The critical number of atoms.
A condensate is the quantum many-body state of $`N`$ bosons in the form $`(\widehat{c}^{})^N|0`$, with $`\widehat{c}^{}`$ denoting a creation operator of a particle in a single-particle wave function $`c(x)`$.
Let us consider a superposition of two condensates
$$|\psi =(\widehat{c}_1^{})^N|0+(\widehat{c}_2^{})^N|0,$$
where the operators $`\widehat{c}_1`$ and $`\widehat{c}_2`$ annihilate a boson in normalized single-particle wave functions $`c_1(x)`$ and $`c_2(x)`$, respectively. These wave functions need not be orthogonal. An unimportant normalization factor of $`|\psi `$ is omitted. Let us also define a spatial overlap of $`c_1`$ and $`c_2`$
$$o|c_1(x)||c_2(x)|\text{d}x.$$
(4)
This quantity assumes values from the interval $`[0,1]`$; it vanishes when the wave functions are spatially separated and is equal to 1 when $`|c_1(x)|=|c_2(x)|`$. In the later case, a density measurement does not distinguish the condensates and the corresponding coefficients (3) would never converge to constant values. Bellow we focus on the cases where $`|c_1(x)||c_2(x)|`$.
The $`N`$-dimensional probability density for this state is given by
$`p_{|\psi }^{(N)}`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}|c_1(x_i)|^2+{\displaystyle \underset{i=1}{\overset{N}{}}}|c_2(x_i)|^2+2{\displaystyle \underset{i=1}{\overset{N}{}}}|c_1(x_i)||c_2(x_i)|`$ (5)
$`\times `$ $`\mathrm{cos}\left({\displaystyle \underset{j=1}{\overset{N}{}}}(\mathrm{arg}(c_1(x_j))\mathrm{arg}(c_2(x_j)))\right).`$
One sees that the total probability due to the last term in (5) is limited from above by $`2o^N`$, while the first and the second term give 2. Thus, if $`|c_1(x)||c_2(x)|`$ then $`o<1`$ and $`o^N0`$ when $`N`$ increases, so the last term can be neglected. The remaining two terms define two sectors of the $`N`$-dimensional space where the probability density is substantial. Again, if $`|c_1(x)||c_2(x)|`$ then these sectors are spatially separated for large enough $`N`$. It can be easily seen, by extending the definition of the overlap (4) to multidimensional products of single-particle wave functions, $`_{i=1}^N|c_1(x_i)||c_2(x_i)|\text{d}x_i`$. Such an overlap is given by $`o^N`$ and vanishes exponentially fast with $`N`$. Finally, if one randomly chooses a point from the $`N`$-dimensional space according to the probability (5), the point will belong to the sector described by either the first term or to the sector described by the second term in (5). Thus, one samples either the condensate in $`c_1(x)`$ or the condensate in $`c_2(x)`$.
The natural requirement for this separation of the two sectors is the small value of the overlap in $`N`$-dimensions. Namely, $`o^N0.01`$, for example. The critical number of particles saturates this inequality
$$N_{crit}=\frac{\mathrm{ln}(0.01)}{\mathrm{ln}(o)},$$
(6)
where the number 0.01 is arbitrary (must be less than 1) and due to logarithm in the numerator of (6), the $`N_{crit}`$ is not sensitive to its precise value. The (6) provides an estimate of an average $`N_{crit}`$ over many repeated measurements. The critical number of particles for a single measurement can differ significantly from (6).
To better understand the meaning of $`N_{crit}`$ consider the case of the Gaussian wave functions
$$|c_1(x)|^2=\frac{1}{\sqrt{\pi }\sigma }e^{\frac{(x+L)^2}{\sigma ^2}},|c_2(x)|^2=\frac{1}{\sqrt{\pi }\sigma }e^{\frac{(xL)^2}{\sigma ^2}}.$$
The overlap is a function of the parameters $`L`$ and $`\sigma `$
$$o=e^{\frac{L^2}{\sigma ^2}}.$$
(7)
From (6) and (7)
$$N_{crit}=\mathrm{ln}(0.01)\frac{\sigma ^2}{L^2}.$$
Reversing the problem one can ask, which Gaussians can be distinguished given the total number of particles $`N`$ and the width $`\sigma `$. The answer: these displaced by more than $`2\sigma \sqrt{\mathrm{ln}(0.01)/N}`$. So, if the state $`|\psi `$ was an equal superposition of all of Gaussians, labeled by the parameter $`L`$, then a single measurement of positions of $`N`$ particles will converge to a subset of Gaussians centered within the distance $`\sigma \sqrt{\mathrm{ln}(0.01)/N}`$ from some random point $`L=L_0`$.
Now, it is clear why the coefficients in (3) assume constant values only approximately. The โlocalizationโ of their values improves with the number of already drawn particles $`r`$ like $`1/\sqrt{r}`$.
Another important example involves two condensates that show shifted interference-like patterns
$$|c_1(x)|^2=\frac{1}{2\pi }(1+\mathrm{cos}(x+\varphi )),|c_2(x)|^2=\frac{1}{2\pi }(1+\mathrm{cos}(x+\rho ))$$
for $`x[0.2\pi ]`$. Assuming $`\pi \rho \varphi \pi `$, the overlap of $`c_1`$ and $`c_2`$ is a function of the relative phase $`\varphi \rho `$
$$o=\left(1+\frac{\varphi \rho }{\pi }\right)\mathrm{cos}\left(\frac{\varphi \rho }{2}\right)+\frac{2}{\pi }\mathrm{sin}\left(\frac{\varphi \rho }{2}\right).$$
For small phase shifts $`\varphi \rho 1`$ the distinction of the two condensates requires at least the number of particles
$$N_{crit}\frac{8\mathrm{ln}(0.01)}{(\varphi \rho )^2}.$$
This example is important because $`N_{crit}`$ above is a critical number of atoms for the measurement to โdecideโ which interference pattern is realized in a single density measurement of two colliding condensates Ketterle ; Javanainen .
## III The method
The previous part of this work establishes that for $`NN_{crit}`$ almost all the positions of particles in a single simulation are randomly drawn from the same one-dimensional probability density. The single measurement on the state $`|\psi `$ can be formally replaced by the measurement on a Bose-Einstein condensate โ the state where all the particles are described by the same single-particle wave function. Indeed, the measurement of particle positions performed on a condensate consists on repeated sampling of a single probability density.
This idea leads to the following postulate:
The act of position measurement on a many-body state $`|\psi `$ is equivalent to the measurement performed on a condensate with the single-particle wave function
$$c(x;Q_1,\mathrm{},Q_M)=\underset{i=1}{\overset{M}{}}Q_iu_i(x),$$
(8)
where $`Q_iq_ie^{i\theta _i}`$ are random complex parameters satisfying the normalization condition
$$\underset{i=1}{\overset{M}{}}|Q_i|^2=1.$$
(9)
The simulation of a single measurement of positions of $`N`$ particles is now reduced either to random generation of $`N`$ positions, each with the probability density $`|c(x;Q_1,\mathrm{},Q_M)|^2`$, or to usage of the $`N|c(x;Q_1,\mathrm{}.,Q_M)|^2`$ as a smoothened density of measured particles.
The only thing left is to provide a probability distribution of random parameters $`Q_i`$. This distribution arises from the projection of the $`N`$-body state of interest on a condensate described by the wave function (8). The operator that annihilates a particle from the condensate mode $`c(x;Q_1,\mathrm{},Q_M)`$ reads
$$\widehat{c}=\underset{i=1}{\overset{M}{}}Q_i^{}\widehat{a}_i.$$
(10)
The condensate of $`N`$ particles is given then by
$$|c;Q_1,\mathrm{},Q_M\frac{1}{\sqrt{N!}}(\widehat{c}^{})^N|0.$$
(11)
Finally, the probability distribution of $`Q_i`$โs for the state $`|\psi `$
$$P(Q_1,\mathrm{},Q_M)=๐ฉ_{|\psi }|c;Q_1,\mathrm{},Q_M|\psi |^2,$$
(12)
with the constrain (9). The normalization factor $`๐ฉ_{|\psi }`$ has to be computed for each state $`|\psi `$ separately, so that $`P(Q_1,\mathrm{},Q_M)\text{d}Q_1\mathrm{}\text{d}Q_M=1`$.
The proposed method consist of two stages: (i) random generation of $`M`$ complex parameters $`Q_i`$ from the probability distribution (12) and under the constrain (9), (ii) using this set of $`Q_i`$โs the wave function (8) is constructed and its modulus squared is taken either as a probability density for generating $`N`$ particle positions or as a smoothened density of particles in a single measurement.
At this point three issues need to be clarified. First, the condensates (11) form an overcomplete and nonorthogonal basis in the Hilbert space of all $`N`$-body bosonic states spanned on the modes $`u_i(x)`$. The nonorthogonality and overcompletness of the condensates are vital if they are to describe the outcomes of measurements that are also โnonorthogonalโ and โovercompleteโ, in the sense that a weighted sum of two or more outcomes is also a valid outcome. Conversely, if the basis formed by the condensates was orthogonal then every result of measurement would be described by one of the basis states and the outcomes corresponding to combinations of condensates would have been unjustly excluded.
Second, the postulate above relies on the observation that for large number of particles $`N`$ the system in a state $`|\psi `$ is driven to a condensate state by the position measurement. In an actual experiment, the positions of all particles are measured at the same instant and the initial state is destroyed due to absorption of light by particles. It is conceivable, though, that starting from a state with macroscopically occupied modes one can create a condensate by measuring positions of a fraction $`N_{crit}`$ of particles. In the case of interacting particles, such a condensate might be short lived since it is not an eigenstate of the Hamiltonian of the system.
Third, the form of a possible outcome of a single density measurement (8) resembles the wave function of a hypothetic condensate in the spontaneous symmetry breaking approach, see Javanainen pages 3 and 4. This work provides yet another justification for this approach: It gives correct description of measurements if $`NN_{crit}`$.
## IV Examples
Fock state. Lets check how the new method works for the Fock state
$$|\psi =|n_1,\mathrm{},n_M.$$
From (12)
$$P(Q_1,\mathrm{},Q_M)\left|Q_1^{n_1}\times \mathrm{}\times Q_M^{n_M}\right|^2,$$
(13)
where $`M`$ is a number of modes of the state. The probability above does not depend on phases $`\theta _i`$ of $`Q_i`$โs, so the phases will be uniformly and independently drawn from the interval $`[0,2\pi )`$. Having said that, one can restrict the probability density to the $`M`$-dimensional space of moduli $`q_i`$. The maximum of the probability (under the constrain (9)) is located at $`q_i=\sqrt{n_i/N}`$. For large number particles per mode the probability is sharply peaked around the maximum point. In fact, one can fix the values of moduli $`q_i=\sqrt{n_i/N}`$ and simply draw the phases. Thus the expected density of particles found in a single experiment is given by
$$|c(x;\theta _1,\mathrm{},\theta _M)|^2=\left|\underset{j=1}{\overset{M}{}}\sqrt{\frac{n_j}{N}}e^{i\theta _j}u_j(x)\right|^2.$$
This example can be immediately extended to the mixtures of Fock states.
In particular, when one considers the state $`|\psi =|N/2,N/2`$ as in Javanainen , with $`u_1(x)=\frac{1}{\sqrt{2\pi }}e^{ix}`$ and $`u_2(x)=\frac{1}{\sqrt{2\pi }}e^{ix}`$ one gets the prediction for a single measurement
$$|c(x;\theta _1,\theta _2)|^2=\frac{1}{2\pi }(1+\mathrm{cos}(x+(\theta _1\theta _2))).$$
Since both phases are random, their difference is also random and the results from Javanainen are recovered: Each density measurement will reveal interference fringes shifted by a random phase.
Schrรถdinger cat state. Since the presented method has ambition to generate outcomes of a single measurement performed on any state with macroscopically occupied modes, it should also work for the Schrรถdinger cat state
$$|\psi =|N,0+|0,N.$$
If corresponding modes $`u_1(x)`$ and $`u_2(x)`$ are spatially separated, the exact calculation shows (see the section where $`N_{crit}`$ is derived), that the single measurement will reveal either all particles in the first mode (in the state $`|N,0`$) or all particles in the second mode (state $`|0,N`$).
The probability density for $`Q_i`$โs
$$P(Q_1,Q_2)q_1^{2N}+q_2^{2N}+2q_1^Nq_2^N\mathrm{cos}(N(\theta _1\theta _2)),$$
under the constrain (9), possesses two well isolated global maxima at $`(q_1=1,q_2=0)`$ and $`(q_1=0,q_2=1)`$ independent of the phases $`\theta _i`$. For large number of particles the probability is appreciable only in the narrow vicinity of these global maxima. This means that one can see either โaliveโ or โdeadโ cat in a single measurement and that superposition of โdeadโ and โaliveโ is strongly suppressed. The approximate character of the method manifests itself in the fact that even if the modes $`u_1(x)`$ and $`u_2(x)`$ are spatially separated, there is nonzero probability for the superposition.
Bogoliubov vacuum. A more complex but very important example is the Bogoliubov vacuum state. For this state, there is also another method of predicting outcomes of a single density measurement available DS . The Bogoliubov vacuum is especially useful to describe evolving Bose-Einstein condensates. It can be expressed in the particle representation as
$$|\psi =\left(\lambda _1\widehat{a}_1^{}\widehat{a}_1^{}+\mathrm{}+\lambda _M\widehat{a}_M^{}\widehat{a}_M^{}\right)^{N/2}|0,$$
where $`\lambda _i`$โs are real numbers, see Zstate ; DS . The corresponding probability function assumes the following form
$$P(Q_1,\mathrm{},Q_M)\left|\left(\underset{i=1}{\overset{M}{}}\lambda _iQ_i^2\right)^{N/2}\right|^2.$$
In the case of just one depletion mode, $`M=2`$, $`\lambda _1=1,|\lambda _2|<1`$ and the probability becomes
$$P(\{Q\})\left(q_1^4+\lambda _2^2q_2^4+2\lambda _2q_1^2q_2^2\mathrm{cos}(2(\theta _1\theta _2))\right)^{N/2}.$$
The last expression shows a single peak at $`(q_1=1,q_2=0)`$ of the variance $`(\mathrm{\Delta }q_2)^21/N`$.
## V Summary
The method presented above consists on sampling of a $`2M`$-dimensional probability density (12). What one has to draw is a set of $`M`$ complex numbers $`Q_i`$ to determine an outcome of a single measurement. This sampling is much simpler, however, than the direct sampling of particle positions from an $`N`$-dimensional probability density. First of all, for the states with macroscopic occupation of modes $`MN`$ and the dimensionality of the sampling problem is strongly reduced. Second, as shown in the examples above, the probability distribution (12) often assumes a simple shape and either analytical calculations can be performed or the importance sampling numerical techniques can be applied.
The justification of the method provides also a link between sampling of a $`N`$-body probability density and the corresponding spontaneous symmetry breaking guess for the outcome of the single density measurement.
## VI Acknowledgments
I am grateful to Jacek Dziarmaga and Krzysztof Sacha for inspirations and enlightening discussions. Work supported by LDRD X1F3 program. |
warning/0506/math0506002.html | ar5iv | text | # Closed and exact functions in the context of Ginzburg-Landau models
## 1. Introduction
Statistical physics has developed a whole variety of interacting particle systems that capture some aspects of the movement of particles on the microscale. An interacting particle system is usually a complex Markov process with a finite or infinite state space. By taking an appropriate scaling limit of an interacting particle system we expect to derive the evolution of the system on the macroscale, in general a nonlinear partial differential equation. It is fairly well understood the transition from the microscopic scale to macroscopic scale, at least for some systems, and in this notes we take for granted this step.
The most interesting microscopic models constructed so far, lack the so called gradient condition. This condition corresponds to the Fickโs law of fluid dynamics according to which the instantaneous current $`w`$ of particles over a bond is the gradient $`\tau hh`$ of some local function. Since the work of Varadhan , Quastel and Varadhan and Yau on nongradient systems, new ideas have been imposed in the field. The main idea is that a nongradient system has a generalized Fickโs law, also called the fluctuation-dissipation equation, of the form
$$w\widehat{a}(m)(\tau hh)+Lg,$$
where $`\widehat{a}`$ is the transport coefficient depending on the particle density $`m`$ in a microscopic cube, $`h`$ is some local function, and $`L`$ is the generator of the microscopic dynamics. The $`Lg`$ part of the approximate equation above is negligible on the macroscale, and is called the fluctuation part of the equation.
One of the main difficulty in finding the scaling limit of a nongradient system is to make rigorous sense of the fluctuation-dissipation equation. As it has been shown in , , the current w, the gradient $`\tau hh`$ and the fluctuations $`Lg`$ are elements of the Hilbert space of โclosed functionsโ and the fluctuation-dissipation equation is a consequence of a direct-sum decomposition of this Hilbert space. The gradient part $`\tau hh`$ of the current $`w`$ that survives after taking the scaling limit of the model is just the projection of $`w`$ onto a one dimensional subspace of the Hilbert space of โclosed functionsโ. The remaining negligible fluctuations $`Lg`$ are vectors of the Hilbert subspace of โexact functionโ.
The purpose of the present paper is not to show how the Hilbert space of โclosed functionsโ and โexact functionsโ arises in the context of interacting particle systems, but rather to motivate the direct-sum decomposition of the Hilbert space of โclosed functionsโ and to find the codimension of the space of โexact functionsโ inside the space of โclosed functionsโ. We calculate this codimension for an arbitrarily chosen vector field. The three continuum models known as the Glauber system, the second-order Ginzburg-Landau system and the continuum solid-on-solid model, also called the fourth-order Ginzburg-Landau system are covered by our general result. Our approach to establish the direct-sum decomposition of the Hilbert space of โclosed functionsโ is new and differs from the approach used before to study the first two models (see Varadhan ). We have followed a different path based on Fourier analysis that has allowed us to handle a general vector field.
## 2. The decomposition theorem
In this section we introduce some terminology and state the main result.
The Hermite polynomials provide an orthogonal basis for the Hilbert space of functions defined on the real axis, that are square integrable with respect to the Gaussian probability measure $`\frac{1}{\sqrt{2\pi }}\mathrm{exp}(\frac{x^2}{2})dx`$. The $`i`$th Hermite polynomial is defined through
$$H_i(x)=\frac{(1)^i}{i!}\mathrm{exp}\left(\frac{x^2}{2}\right)\left(\frac{d^i}{dx^i}\mathrm{exp}\left(\frac{x^2}{2}\right)\right),i.$$
We stress that $`H_i`$ is not normalized to have $`L^2`$ norm 1 with respect to the probability measure $`\frac{1}{\sqrt{2\pi }}\mathrm{exp}(\frac{x^2}{2}))dx`$, but rather $`\frac{1}{\sqrt{i!}}.`$
There is an extension of Hermite polynomials to more variables. A multi-index is a double-sided infinite sequence $`I=\{i_n\}_n`$ of positive integers, with at most finitely many non-zero entries. The degree of a multi-index is $`|I|=_ni_n`$. Call $``$ the set of multi-indices and $`_N`$ the set of multi-indices of fixed degree $`N`$. The multidimensional Hermite polynomials are
$$H_I(x)=\mathrm{\Pi }_nH_{i_n}(x_n),I.$$
We make the convention that if a multi-index $`I`$ has some strictly negative entries then $`H_I=0`$. Together the multidimensional Hermite polynomials, $`\{H_I\}_I`$ form an orthogonal basis for the Hilbert space of functions defined on $`^{}`$, that are square integrable with respect to the probability measure
$$d\nu _0^{gc}=\underset{i}{}\frac{1}{\sqrt{2\pi }}\mathrm{exp}\left(\frac{x_i^2}{2}\right)dx_i.$$
It is interesting to note that this Hilbert space is a model for the symmetric Fock space over the space of square summable, double-sided sequences $`l^2()`$, and decomposes as a direct sum of the degree $`N`$ subspaces
$$_N=\{H_I||I|=N\}^c.$$
The superscript on the line above, means that we take the closed linear span of the set.
The shift $`\tau `$ acts on configurations as $`(\tau (x))_n=x_{n+1}`$ and on functions as $`(\tau f)(x)=f(\tau x)`$. $`\tau ^n`$ stands for the $`n`$-fold composition $`\tau \mathrm{}\tau `$. If a multi-index $`I=(i_n)_n`$ has $`i_n=0`$ for all $`n<0`$ we shall say that the multi-index is supported on the set of positive integers. We shall use the notation $`\delta _n`$ for the multi-index that corresponds to the configuration with a single particle at the site $`n`$. Two multi-indices can be added and the addition is point-wise.
The action of the annihilation, creation, and shift operators on the multidimensional Hermite polynomial $`H_I`$ is very simple:
$$_nH_I(x)=H_{I\delta _n}(x),(x_n_n)H_I(x)=H_{I+\delta _n}(x),\tau H_I=H_{\tau ^1I}.$$
Above $`_n`$ stands for the partial derivative with respect to the $`n`$th coordinate.
Given a double-sided sequence of real numbers $`(a_k)_k`$, that are all but finitely many zero we introduce the vector field $`D_0=_ka_k_k`$ with constant coefficients. Translating $`a`$โs to the left or to the right produces a new sequence that defines the vector field $`D_n=_ka_k_{k+n}`$, $`n`$. Now we have the setup needed to introduce the closed and exact functions.
###### Definition 2.1.
We shall say that a function $`\xi L^2(^{},d\nu _0^{gc})`$ is closed (or more precisely, $`D_0`$-closed) if it satisfies in the weak sense
(1)
$$D_n(\tau ^m\xi )=D_m(\tau ^n\xi )$$
for all integers $`m`$ and $`n`$. Let $`๐_D`$ denote the space of all $`D_0`$-closed functions.
###### Definition 2.2.
We shall say that a function $`\xi ^gL^2(^{},d\nu _0^{gc})`$ is exact (or more precisely $`D_0`$-exact) if there is a local function $`g`$, a function that depends on finitely many co-ordinates, such that
(2)
$$\xi ^g=D_0\left(\underset{k}{}\tau ^kg\right)=\underset{k}{}D_0(\tau ^kg).$$
Let $`_D`$ denote the closed linear span of the set of $`D_0`$-exact functions.
Although the infinite sum $`_k\tau ^kg`$ does not make sense, after applying the differential operator $`D_0`$ we get a meaningful expression. Since $`g`$ is a local function, the vector field $`D_0`$ kills all but finitely many terms of the infinite formal sum.
The terminologies of exact and closed functions are not arbitrarily chosen. We can define formally the form $`w=_n\tau ^n\xi dx_n`$ and the boundary operator $`df=_nD_n(f)dx_n.`$ It is not hard to see, with these new definitions, that the form $`w`$ is closed ($`dw=0`$), in the vector calculus sense, if and only if $`D_n(\tau ^m\xi )=D_m(\tau ^n\xi )`$, i.e., if and only if $`\xi `$ is a closed function.
Knowing that any exact function is closed a natural question to ask is about the codimension of the space of exact functions inside the space of closed functions. In this paper we provide the answer for this question.
###### The Decomposition Theorem 2.1.
Let $`D_0=_ka_k_k`$ be a vector field with constant real coefficients. All but finitely many numbers in the sequence $`(a_k)_k`$ are zero. The following decomposition results hold:
* If the sum of the coefficients of the vector field $`D_0`$ is not equal to zero then
$$๐_D=_D.$$
* If the sum of the coefficients of the vector field $`D_0`$ is equal to zero then
$$๐_D=\mathrm{๐}_๐.$$
Idea of the proof for the decomposition theorem 2.1. We outline the main ideas used to prove the decomposition theorem. We shall show later that a function $`\xi `$ is $`D_0`$-closed if and only if the projections $`\mathrm{Proj}__N\xi `$, $`N0`$ are $`D_0`$-closed. Degree $`0`$ subspace is easy to analyze since it is one dimensional. Any constant function is always $`D_0`$-closed, but is exact if and only if the sum of the coefficients of $`D_0`$ is not equal to zero. If the sum of the coefficients of the $`D_0`$ is equal to zero, then any $`D_0`$-closed function is orthogonal on the degree $`0`$ subspace. Therefore the result of the theorem holds if we can prove that a given $`D_0`$-closed function $`\xi `$ in $`_N`$, $`N1`$, the function $`\xi `$ can be approximated with $`D_0`$-exact functions.
We shall investigate the properties of the Fourier coefficients of closed and exact functions, and we shall rather establish that the Fourier coefficients of a closed functions can be approximated in the appropriate sense with Fourier coefficients of exact functions. The ideas will be elaborated in the following sections.
Note. In two cases relevant for statistical physics questions, namely the second-order Ginzburg-Landau vector field $`Y_0=_1_0`$ and the fourth-order Ginzburg-Landau vector field $`X_0=_12_0+_1`$, the decomposition result of Theorem 2.1 is equivalent with the fluctuation-dissipation equation mentioned in the introduction section of the paper.
Note. To get a flavor of the result stated in Theorem 2.1 we give some examples of exact and closed functions in the case of the fourth-order Ginzburg-Landau field, $`X_0=_12_0+_1`$: $`x_n+x_n2x_0`$ are $`X_0`$-exact, $`\mathrm{๐}`$, $`x_0`$, $`x_n+x_n`$ are examples of $`X_0`$-closed but not $`X_0`$-exact functions. A strange phenomena appears for: besides the function 1 there exists another function that is $`X_0`$-closed and not $`X_0`$-exact, namely $`x_0`$. Therefore one may expect that the codimension of the space of exact functions is two. This is not the case and $`x_0`$ can be approximated with exact functions.
## 3. The set of multi-indices
A multi-index $`I=\{i_n\}_n`$ can be thought of as a configuration of particles sitting on the sites of the lattice $``$. On top of the site $`n`$ sit $`i_n`$ particles. Rather than saying how many particles are at each site, we give the positions of the particles. This way we obtain a vector
(3)
$$z_I=(\underset{i_{n_1}}{\underset{}{n_1\mathrm{}n_1}},\mathrm{},\underset{i_{n_k}}{\underset{}{n_k\mathrm{}n_k}}).$$
that lists, in increasing order, all occupied sites of $`I`$ repeated according to the number of particles that occupy the site. We assume the only non-zero entries of the multi-index $`I`$ are $`i_{n_1},\mathrm{}i_{n_k}`$. Note that the dimension of the vector $`z_I`$ is the degree of the multi-index $`I`$. If the multi-index has zero degree then $`z_I`$ is just a point. We say that $`z_I`$ is a new coding of the multi-index $`I`$. This correspondence shows that the set $`_N`$ is bijective with the set of vectors of $`^N`$ with entries in increasing order or is in bijection with the quotient space $`^N/S_N`$, where $`S_N`$ is the group of permutations of $`N`$ letters.
For the results that follow we need to say more about the set of multi-indices. We partition the set of multi-indices into orbits with the help of the group action
(4)
$$\times ^{}^{}(n,I)nI:=\tau ^n(I\delta _n+\delta _0).$$
When restricted to $`\times `$ the map (4) is not an action any more since the multi-indices that enumerate the basis of the $`L^2`$ space are constrained to have positive entries.
The orbits of the action (4) provide a partition of the set of multi-indices $`^{}`$. For each multi-index $`I`$ we define $`o(I)`$ to be the shadow of the orbit of $`I`$ on the set $``$, i.e., $`o(I)=\{J|J=nIn\}=\{J|J=nIns(I)\}.`$ Here, $`s(I)=\{n|i_n0\}`$ is the finite set of occupied positions of $`I`$. From now on we will refer to $`o(I)`$ as the orbit of $`I`$, although this is just a part of the actual orbit of the action. It has the advantage of being finite since the multi-index $`I`$ has all but finitely many entries zero and there are just finitely many $`n`$โs that after acting on $`I`$ give rise to a multi-index with positive entries. All the multi-indices in the same orbit have the same degree. The orbits partition $``$ and $`_N`$. Call $`๐ช`$ the set of orbits, and $`๐ช_N`$ the set of orbits containing multi-indices of degree $`N`$.
It is worth mentioning that inside each orbit $`o(I)`$ there exists a unique representative supported on the positive integers. Denote this multi-index by $`R(o(I))`$. To see that this is true let us assume that $`I`$ has some particles in some negative position, i.e., there exists some $`n<0`$ such that $`i_n1`$. If $`k`$ is the leftmost occupied position of $`I`$ and $`k<0`$, then $`kIo(I)`$ is supported on the positive integers. Let us call $``$ the set of all representatives, and $`_N`$ the set of degree-$`N`$ representatives.
So far the orbit space $`\{o(I)\}_I`$ is an abstract object. Fortunately we are able to give a concrete description of the orbit space. For this purpose it is very useful to know that each orbit, $`o`$ has a unique representative $`R(o)`$ supported on the positive semi-axis. The vector $`z_{R(o)}`$ is a point in the positive cone $`๐_N^+=\{z^N|z=(z_1,\mathrm{},z_N),\mathrm{\hspace{0.33em}0}z_1\mathrm{}z_N\}`$. Therefore the set of representatives, and in particular the set of orbits $`๐ช_N`$, are bijective with the cone $`๐_N^+`$. Since there is only one multi-index with zero degree, $`\mathrm{๐}=(0)_n`$, the sets $`_0`$, $`๐ช_0`$ and $`_0`$ contain just a single element. By convention, $`๐_0^+`$ is just one-point set.
We can say even more about this picture. The cone $`๐_N^+`$, itself is an orbit space, which we shall describe below.
Let us define the transformations that acts on the lattice $`^N`$, for any $`1i,jN`$,
(5)
$$\sigma _{i,j}:^N^N,\sigma _{i,j}(z_1,\mathrm{},z_i,\mathrm{},z_j\mathrm{},z_N)=(z_1,\mathrm{},z_j,\mathrm{},z_i,\mathrm{},z_N)$$
$$\mathrm{and}\gamma _1:^N^N,\gamma _1(z_1,z_2,\mathrm{},z_N)=(z_1,z_2z_1\mathrm{},z_Nz_1).$$
The smallest group generated by $`\sigma _{i,j}`$, $`1i,jN`$ and $`\gamma _1`$ will be denoted by $`\stackrel{~}{S}_N`$. To see that $`\stackrel{~}{S}_N`$ is isomorphic with the group of permutations of $`N`$ letters, we write down the basic relations among the generating transformations: $`(\gamma _1\sigma _{1,2})^3=\mathrm{๐ข๐}`$ and $`\gamma _1\sigma _{i,i+1}=\sigma _{i+1,i}\gamma _1`$, $`1iN1`$. The group $`\stackrel{~}{S}_N`$ has $`S_N`$, the group of permutations of $`N`$ letters as subgroup, and $`\stackrel{~}{S}_N`$ decomposes into left cosets with respect to $`S_N`$, as $`\stackrel{~}{S}_N=S_N\gamma _1S_N\mathrm{}\gamma _NS_N`$, where the transformations $`\gamma _i`$ are
(6)
$$\gamma _i:^N^N,\gamma _i(z_1,z_2,\mathrm{},z_N)=(z_i,z_2z_i\mathrm{},z_Nz_i).$$
It is interesting to note that $`^N/\stackrel{~}{S}_N`$ is bijective with the cone $`๐_N^+`$, as the next argument proves. Any orbit of $`^N/\stackrel{~}{S}_N`$ contains at least one vector, let say $`z`$, with components in increasing order. If this vector does not have positive co-ordinates, it means $`z_1<0`$. But $`(z_1,z_2z_1,\mathrm{},z_Nz_1)`$ is still a point in the orbit of $`z`$ under the action of $`\stackrel{~}{S}_N`$. We can rearrange the coordinates of the new vector to be in increasing order and hence the orbit of $`z`$ under the action of $`\stackrel{~}{S}_N`$ contains at least one vector of the cone $`๐_N^+`$. To see that the orbit of $`z`$ does not contain more than one vector of $`๐_N^+`$ we use the coset decomposition of $`\stackrel{~}{S}_N`$. If $`z`$ is in $`๐_N^+`$, then rearranging the co-ordinates of $`z`$ we obtain either the vector $`z`$ or some vector outside the cone $`๐_N^+`$. If we act on $`z`$ or some other vector obtained from $`z`$ by changing the places of the co-ordinates, with either of the transformations $`\gamma _1,\mathrm{}\gamma _N`$ we get a vector that has at least one negative co-ordinate, so does not belong to $`๐_N^+`$.
Now we can say that the set of orbits $`๐ช_N`$ is bijective with the cone $`๐_N^+`$, and hence with the quotient space $`^N/\stackrel{~}{S}_N`$. The bijection is $`o๐ช_Nz_{R(o)}๐_N^+`$.
In addition, if $`I`$ and $`J`$ are two multi-indices in the same orbit of the action (4) then $`z_I`$ and $`z_J`$ are in the same orbit of the action of $`\stackrel{~}{S}_N`$ on $`^N`$. Assume that $`J=n_jI`$ with $`I=_{i=1}^ka_i\delta _{n_i}`$, where $`a_i0`$ and $`n_1\mathrm{}n_k`$. Then $`J=_{i=1,\mathrm{},k,ij}a_i\delta _{n_in_j}+(a_j1)\delta _0+\delta _{n_j}`$, and so
$$z_I=(\underset{a_1}{\underset{}{n_1,\mathrm{},n_1}},\mathrm{},\underset{a_k}{\underset{}{n_k,\mathrm{},n_k}})$$
$$z_J=(n_j,\underset{a_1}{\underset{}{n_1n_j,\mathrm{},n_1n_j}},\mathrm{},\underset{a_j1}{\underset{}{0,\mathrm{},0}},\mathrm{},\underset{a_k}{\underset{}{n_kn_j,\mathrm{},n_kn_j}}).$$
It follows that $`z_J`$ is the image of $`z_I`$ under some element of $`\stackrel{~}{S}_N`$.
We shall denote by $`z\stackrel{S_N}{}z^{}`$ and $`z\stackrel{\stackrel{~}{S}_N}{}z^{}`$ two lattice points $`z`$ and $`z^{}`$ that have the same image in the quotient space $`^N/S_N`$ and $`^N/\stackrel{~}{S}_N`$, respectively.
Before we leave this section it is important to notice the following crucial facts. Let $`N1`$. Since $`_N`$ is identified with $`^N/S_N`$ we can think of any function $`\widehat{\xi }:_N`$ as being a $`S_N`$-invariant function $`\widehat{\xi }:^N`$, where $`\widehat{\xi }(z)=\widehat{\xi }(I)`$ if $`z\stackrel{S_N}{}z_I`$. Similarly, since $`๐ช_N`$ is identified with $`^N/\stackrel{~}{S}_N`$ we can think of any function $`c:๐ช_N`$ as being a $`\stackrel{~}{S}_N`$-invariant function $`\stackrel{~}{c}:^N`$, where $`\stackrel{~}{c}(z)=c(o)`$ if there exists a multi-index $`Io`$ such that $`z\stackrel{S_N}{}z_I`$.
## 4. Properties of closed functions and of exact functions
This section contains a detailed study of closed and exact functions.
Closed functions. We start with a very simple but important property of closed functions.
###### Lemma 4.1.
Assume $`D_0`$ is a vector field with constant coefficients, $`D_0=_ka_k_k`$. All but finitely many coefficients of the vector field $`D_0`$ are zero. A function $`\xi L^2(^{},d\nu _0^{gc})`$ is $`D_0`$-closed if and only if the projection $`Proj__N\xi `$ onto the degree $`N`$ subspace $`_N`$ is $`D_0`$-closed, for any $`N0`$.
Proof. We denote by $`_j`$ the differential operator with respect to the $`j^{\mathrm{th}}`$ coordinate, and by $`_j^{}=_j+x_j`$ the adjoint operator of $`_j`$. The adjoint is taken with respect to the inner product $`<,>`$ of $`L^2(^{},d\nu _0^{gc})`$. The operators $`_j`$ and $`_j^{}`$ are bounded operators when restricted to a degree subspace, although they are unbounded on the whole $`L^2(^{},d\nu _0^{gc})`$ space.
If $`\xi _N`$, with Fourier series $`\xi =_{I_N}\widehat{\xi }_IH_I`$, then the image of $`\xi `$ under the operator $`_j`$ is $`_j(\xi )=_{I_N}\widehat{\xi }_IH_{I\delta _j},`$ with the convention that if the multi-index $`I\delta _j`$ has some negative entries then $`H_{I\delta _j}=0`$. For a function $`fL^2(^{},d\nu _0^{gc})`$ denote by $`f=\sqrt{<f,f>}`$ the $`L^2`$ norm of $`f`$.
The operators $`_j`$ and $`_j^{}`$ act on the degree $`N`$ subspaces as follows:
$$_j(_N)_{N1},N1,_j^{}(_N)_{N+1},N0.$$
The boundedness of these operators follows from the observation that
$$\frac{1}{(N!)^N}\underset{I_N}{inf}H_I^2\underset{I_N}{sup}H_I^21,$$
and from the existence of two strictly positive constants, $`C_1^N`$, $`C_2^N`$, that depend just on $`N`$ such that
(7)
$$C_1^N\underset{I_N}{}\widehat{\xi }_I^2\xi ^2C_2^N\underset{I_N}{}\widehat{\xi }_I^2,C_1^N\underset{I_N}{}\widehat{\xi }_I^2_j(\xi )^2C_2^N\underset{I_N}{}\widehat{\xi }_I^2.$$
Indeed
$$_j\xi ^2=\underset{I_N}{}\widehat{\xi }_I^2H_{I\delta _j}^2\underset{I_N}{}\widehat{\xi }_I^2(N!)^N\underset{I_N}{}\widehat{\xi }_I^2H_I^2(N!)^N\xi ^2.$$
and hence the norm of the operator $`_j:_N_{N1}`$ is bounded above by $`(N!)^N`$.
The vector field $`D_0`$ with constant coefficients has similar properties:
$$D_0(_N)_{N1},N1,D_0^{}(_N)_{N+1},N0.$$
For any function $`\xi L^2(^{},d\nu _0^{gc})`$ and any test function $`\varphi _{N1}`$ we have
$$<D_n(\tau ^m\xi ),\varphi >=<\xi ,\tau ^m(D_n^{}\varphi )>=<\mathrm{Proj}__N\xi ,\tau ^m(D_n^{}\varphi )>=$$
(8)
$$=<D_n(\tau ^m\mathrm{Proj}__N\xi ),\varphi >,$$
$$<D_m(\tau ^n\xi ),\varphi >=<\xi ,\tau ^n(D_m^{}\varphi )>=<\mathrm{Proj}__N\xi ,\tau ^n(D_m^{}\varphi )>=$$
(9)
$$=<D_m(\tau ^n\mathrm{Proj}__N\xi ),\varphi >.$$
It follows that $`D_n(\tau ^m\xi )=D_m(\tau ^n\xi )`$ in the weak sense if and only if $`D_n(\tau ^m\mathrm{Proj}__N\xi )=D_m(\tau ^n\mathrm{Proj}__N\xi )`$ in the strong sense for all $`N0`$.
We recall that a function $`\xi `$ is closed if and only if $`D_n(\tau ^m\xi )=D_m(\tau ^n\xi )`$ for all $`m,n`$, which, by the previous equalities (8) and (9), is equivalent to
$$D_n(\tau ^m\mathrm{Proj}__N\xi )=D_m(\tau ^n\mathrm{Proj}__N\xi )m,nN0.$$
Therefore, a function $`\xi `$ is closed if and only if $`\mathrm{Proj}__N\xi `$ is closed for all $`N0`$. โ
Note. If $`\xi =_{I_N}\widehat{\xi }_IH_I`$ is a function inside the space $`_N`$, two norms can be defined for $`\xi `$: the $`L^2`$ norm $`\xi `$ and the sum of squared Fourier coefficients $`_{I_N}\widehat{\xi }_N^2`$. It is important to note the inequality (7) implies that these two norms define the same topology on the space $`_N`$.
Note. Assume $`\xi _N`$ is a $`D_0`$-closed function, with Fourier series expansion $`\xi =_I\widehat{\xi }_IH_I`$. We calculate,
$$D_n\xi =\underset{I}{}\left[\underset{k}{}a_k\widehat{\xi }_{I+\delta _{(n+k)}}\right]H_I,D_0(\tau ^n\xi )=\underset{I}{}\left[\underset{k}{}a_k\widehat{\xi }_{\tau ^n(I+\delta _k)}\right]H_I.$$
Therefore a function is closed if and only if its Fourier coefficients satisfy the relations:
(10)
$$\underset{k}{}a_k\widehat{\xi }_{I+\delta _{(n+k)}}=\underset{k}{}a_k\widehat{\xi }_{\tau ^n(I+\delta _k)}n,I.$$
Construction of exact functions. It is important to have some examples of functions that are exact. The functions that will be constructed next will be used in the proof of the decomposition theorem 2.1, to approximate closed functions with exact ones.
###### Lemma 4.2.
Let $`c`$ be a function defined on the set of orbits with finite support (i.e., c(o)=0 except for finitely many orbits $`o`$). The function
$$\xi =\underset{o๐ช}{}c(o)D_0\left[\underset{n}{}\tau ^nH_{R(o)+\delta _0}\right]$$
is $`D_0`$-exact and the Fourier coefficients of $`\xi `$ are
(11)
$$\widehat{\xi }_I=\underset{k}{}a_kc(o(\tau ^kI)).$$
Proof. The function $`\xi `$, that has been introduced is well-defined since the sum is over a finite set, and is exact as a sum of exact functions. To conclude the lemma we need to calculate the Fourier coefficients of $`\xi `$. We have,
(12) $`\xi `$ $`=`$ $`{\displaystyle \underset{o๐ช}{}}c(o)D_0\left[{\displaystyle \underset{n}{}}H_{\tau ^n(R(o)+\delta _0)}\right]={\displaystyle \underset{o๐ช,n}{}}c(o){\displaystyle \underset{k}{}}a_kH_{\tau ^n(R(o)+\delta _0\delta _{n+k})}=`$
$`={\displaystyle \underset{o๐ช,n}{}}c(o){\displaystyle \underset{k}{}}a_kH_{\tau ^k[(n+k)R(o)]}={\displaystyle \underset{I}{}}{\displaystyle \underset{k}{}}a_kc(o(\tau ^kI))H_I.`$
To justify the integration by parts in the calculation above (12) we make the following observation. For any multi-index $`I`$ there exists a unique orbit $`o๐ช`$ and a unique integer $`n`$ such that $`I=\tau ^k[(n+k)R(o)]`$. This is a consequence of the freeness of the action (4). Moreover, the orbit $`o`$ is the same as $`o(\tau ^kI)`$. We stress again that the sums in (12) are over finite sets as $`c`$ has finite support. Actually all computations that we carried out to prove this lemma are valid because $`c`$ is a function with finite support and the sums are finite, although this wasnโt emphasized each time we used it. Also we have made use of the convention that $`H_I=0`$ if $`I`$ is a multi-index with negative entries. โ
###### Lemma 4.3.
Let $`N1`$ be a natural number and $`e=(1,\mathrm{},1)^N`$. In addition if $`\stackrel{~}{c}`$ is a real-valued function defined on $`^N`$, with finite support and $`\stackrel{~}{S}_N`$-invariant then the function
(13)
$$\xi _{\stackrel{~}{c}}=\underset{I_N}{}\left(\underset{k}{}a_k\stackrel{~}{c}(z_Ike)\right)H_I$$
is a well-defined $`D_0`$-exact function in the degree $`N`$ subspace $`_N`$.
Proof. This lemma follows from lemma 4.2. Since $`\stackrel{~}{c}:^N`$ is $`\stackrel{~}{S}_N`$-invariant, it makes sense to introduce $`c:๐ช`$, where $`c(o)=\stackrel{~}{c}(z_I)`$ if $`I`$ is a multi-index in the orbit $`o`$ of degree $`N`$, and $`c(o)=0`$ otherwise. We should note that if $`I`$ is a multi-index in the orbit $`o`$ then $`z_I+ke=z_{\tau ^kI}`$ and $`\stackrel{~}{c}(z_Ike)=c(o(\tau ^kI))`$. Hence the Fourier coefficients of the function $`\xi _{\stackrel{~}{c}}`$ are of the form $`_ka_kc(o(\tau ^kI))`$, and the function $`\xi _{\stackrel{~}{c}}`$ is $`D_0`$-exact. โ
In the previous lemma an operator has come out in a natural way in our construction of exact functions. Below we provide the exact definition of this operator.
###### Definition 4.1.
Let $`D_0=_ka_k_k`$ be a vector field with constant coefficients, all the coefficients being zero except finitely many. The vector field $`D_0`$ defines an operator $`T_{D_0}`$ that acts on functions $`c:^N`$ and produces a function $`T_{D_0}c:^N`$, where
$$(T_{D_0}c)(z)=\underset{k}{}a_kc(zke),z^N.$$
Above $`e`$ is the vector $`(1,\mathrm{},1)`$ of the lattice $`^N`$.
## 5. Proof of the decomposition theorem 2.1
We start by listing two important properties of the operator $`T_{D_0}`$ introduced at the end of the previous section.
###### Lemma 5.1.
Let $`c`$ be a real-valued function defined on the lattice $`^N`$, $`N1`$. We assume that the function $`c`$ is square-summable and $`\stackrel{~}{S}_N`$\- invariant. Then, there exists a sequence $`(c_n)_{n1}`$ of real-valued, finitely supported, $`\stackrel{~}{S}_N`$-invariant functions such that $`T_{D_0}c_nT_{D_0}c`$ as $`n\mathrm{}`$ and the convergence is in the Hilbert space topology of $`L^2(^N)`$.
Proof. We define a sequence of $`\stackrel{~}{S}_N`$-invariant regions of the lattice $`^N`$, namely
(14)
$$P_i=_{\gamma \stackrel{~}{S}_N}\gamma \{z=(z_1,\mathrm{},z_N)^N|0z_1\mathrm{}z_Ni1\},i1.$$
For the reader convenience we add two pictures of the region $`P_i`$ in dimension $`N=1`$, respectively $`N=2`$. In dimension $`N=1`$ the region $`P_i`$ contains the lattice points inside the segment $`[i+1,i1]`$, whereas in dimension $`N=2`$ the region $`P_i`$ contains the lattice points inside the hexagon shown below.
Figure 1. The region $`P_i`$ in dimension $`N=1`$.
Figure 2. The region $`P_i`$ in dimension $`N=2`$.
Beside being $`\stackrel{~}{S}_N`$-invariant, the sequence of regions $`(P_i)_{i1}`$ defined above, grows to cover the entire lattice $`^N`$ as $`i\mathrm{}`$. Define $`c_n`$ to be $`c\mathrm{๐}_{P_n}`$, for $`n1`$. Since $`\mathrm{๐}_{P_n}`$ is the characteristic function of the region $`P_n`$ we have immediately that $`c_n`$ is a finitely supported, $`\stackrel{~}{S}_N`$-invariant function. Square-summability of $`c`$ implies that $`c_nc`$ as $`n\mathrm{}`$ in the topology of $`L^2(^N)`$ (the norm $`cc_n^2=_{zP_n}c^2(z)`$ involves only the values of $`c`$ outside the region $`P_n`$, and these values decay to zero as $`n\mathrm{}`$ since $`c`$ is square-summable). Then, obviously, $`c_nc`$ and $`Tc_nTc`$ as $`n\mathrm{}`$ in the topology of $`L^2(^N)`$. โ
Below we discuss certain facts about the Fourier transform of functions defined on the lattice $`^N`$. The Fourier transform of a function $`c:^N`$ is formally defined to be
$$c:[\pi ,\pi )^N,c(\alpha )=\frac{1}{\sqrt{2\pi }}\underset{z^N}{}c(z)e^{iz\alpha }.$$
In the exponent above $`z\alpha `$ stands for the dot product $`z_1\alpha _1+\mathrm{}+z_N\alpha _N`$. The reader may consult Rudin for an extended treatment of Fourier transform of functions defined on lattice. We remind the reader that $``$ is an isometry between the spaces $`L^2(^N)`$ and $`L^2([\pi ,\pi )^N)`$. The space $`L^2([\pi ,\pi )^N)`$ is considered with respect to the Lebesgue measure on $`[\pi ,\pi )^N`$. Also if $`c`$ is invariant under a certain group of transformations then $`c`$ is invariant, as well. Though the symmetry group of $`c`$ might not coincide with the symmetry group of $`c`$. Indeed if $`c`$ is symmetric, or $`S_N`$-invariant then $`c`$ is symmetric. Now suppose that $`c`$ is $`\stackrel{~}{S}_N`$-invariant then $`c`$ is invariant under the action of the group $`\stackrel{~}{\mathrm{\Sigma }}_N`$, generated by the transformations:
$$s_{ii+1}:[\pi ,\pi )^N[\pi ,\pi )^N,1iN1$$
$$s_{ii+1}(\alpha _1,\mathrm{},\alpha _N)=(\alpha _1,\mathrm{},\alpha _{i+1},\alpha _i,\mathrm{},\alpha _N),$$
and
$$g:[\pi ,\pi )^N[\pi ,\pi )^N,g(\alpha _1,\mathrm{},\alpha _N)=(\mathrm{mod}_{2\pi }(\alpha _1\mathrm{}\alpha _N),\alpha _2,\mathrm{},\alpha _N).$$
On the line above we used the notation $`\mathrm{mod}_{2\pi }(t)`$. Any real number $`t`$ can be written uniquely as $`2\pi a+b`$, where $`a`$ is an integer number and $`b`$ is a real number in the interval $`[\pi ,\pi )`$. By $`\mathrm{mod}_{2\pi }(t)`$ we denote the remainder $`b`$. It is also true that if the Fourier transform $`c`$ is $`\stackrel{~}{\mathrm{\Sigma }}_N`$-invariant then $`c`$ is $`\stackrel{~}{S}_N`$-invariant.
In section 4 we have established that a function $`\xi =_{I_N}\widehat{\xi }_IH_I`$ is $`D_0`$-closed if and only if the following holds:
(15)
$$\underset{k}{}a_k\widehat{\xi }_{I+\delta _{(n+k)}}=\underset{k}{}a_k\widehat{\xi }_{\tau ^n(I+\delta _k)}n,I.$$
Obviously we can use the Fourier coefficients of $`\xi `$ to construct a $`S_N`$-invariant function $`\widehat{\xi }:^N`$, $`\widehat{\xi }(z)=\widehat{\xi }_I`$ if $`z\stackrel{S_N}{}z_I`$. The relations (15) force our function $`\xi `$ to satisfy
(16)
$$\underset{k}{}a_k\widehat{\xi }(z+ke_1)=\underset{k}{}a_k\widehat{\xi }(zz_1e(z_1+k)e_1),z=(z_1,\mathrm{},z_N)^N.$$
The vectors $`e`$ and $`e_1`$ of the lattice $`^N`$ are $`(1,\mathrm{},1)`$ and $`(1,0,\mathrm{},0)`$, respectively. After applying the Fourier transform in both sides of the equation (16) we find that $`\widehat{\xi }`$ satisfies
(17)
$$p(e^{i\alpha _1})(\widehat{\xi })(\alpha )=p(e^{i(\alpha _1+\mathrm{}+\alpha _N)})(\widehat{\xi })(g(\alpha )),\alpha =(\alpha _1,\mathrm{},\alpha _N)[\pi ,\pi )^N,$$
where $`p`$ is the rational function $`p(x)=_ka_kx^k`$. To wrap up our argument we can say that the $`D_0`$-closedness condition implies property (17).
Next we shall establish a crucial fact about functions that satisfy relation (17).
###### Lemma 5.2.
Let $`\widehat{\xi }`$ be a real-valued, $`S_N`$-invariant function defined on the lattice $`^N`$, $`N1`$ that satisfies (17). Then, there exists a sequence $`(c_n)_{n1}`$ of real-valued, square-summable, $`\stackrel{~}{S}_N`$-invariant functions defined on the lattice $`^N`$ such that $`T_{D_0}c_n\widehat{\xi }`$ as $`n\mathrm{}`$ in the topology of $`L^2(^N)`$.
Proof. Examples of functions $`\widehat{\xi }`$ satisfying the properties listed in the hypothesis of this lemma, are functions constructed, as explained before in this section, from the Fourier coefficients of closed functions.
The first step towards establishing our result is to solve, at least on a formal level, the equation $`T_{D_0}c=\widehat{\xi }`$. The Fourier transform for functions defined on the lattice will help us to make our guess.
After applying the Fourier transform in each side of the equation $`T_{D_0}c=\widehat{\xi }`$, we get
$$p(e^{i(\alpha _1+\mathrm{}+\alpha _N)})(c)(\alpha )=(\widehat{\xi })(\alpha ),\alpha =(\alpha _1,\mathrm{},\alpha _N)^N.$$
where $`p`$ is the rational function $`_ja_jx^j`$ canonically associated to $`T_{D_0}`$.
After multiplying with a high enough power of $`x`$ the equation $`p(x)=0`$, we see that any solution $`x`$ of $`p(x)=0`$ has to be the root of a certain polynomial. Since the number of roots of any polynomial is finite, the number of solutions of $`p(x)=0`$ is finite, as well. If the equation $`p(x)=0`$ has no solutions on the unit circle then the equation $`T_{D_0}c=\widehat{\xi }`$ can be solved in the space $`L^2(^N)`$. Indeed the unique, square-summable solution of $`T_{D_0}c=\widehat{\xi }`$ is
$$c=^1\left(\frac{1}{p(e^{i(\alpha _1+\mathrm{}+\alpha _N)})}(\widehat{\xi })\right).$$
That $`\widehat{\xi }`$ has been built out of the Fourier coefficients of a closed function implies that $`c`$ is $`\stackrel{~}{S}_N`$-invariant. Hence the lemma is proved in this case. We can choose $`c_n=c`$, for any $`n1`$.
A more involved case is when the equation $`p(x)=0`$ has solutions on the unit circle. If the sum of the coefficients of $`p`$ is equal to $`0`$ then the number $`1`$ is a solution of $`p(x)=0`$. The cases arising from interacting particle models are of this kind. The difficulty in this case arises because the equation $`Tc=\widehat{\xi }`$ can not be solved in $`L^2(^N)`$. To get around this problem we shall consider a slightly modified equation $`(T_{D_0}c)=(\widehat{\xi })1_{A_n}`$. The function $`\widehat{\xi }`$ is multiplied with the characteristic function of a set $`A_n`$ that in $`\stackrel{~}{\mathrm{\Sigma }}_N`$-invariant and carefully chosen to avoid the unit roots of $`p`$. More precisely,
$$A_n=_{\gamma \stackrel{~}{\mathrm{\Sigma }}_N}\left\{\gamma (\alpha )\right|\alpha =(\alpha _1,\mathrm{},\alpha _N)[\pi ,\pi ]^N,$$
$$|\mathrm{mod}_{2\pi }(\alpha _1+\mathrm{}+\alpha _N)r_k|>\frac{1}{n},e^{ir_k}\mathrm{unit}\mathrm{root}\mathrm{of}p\}.$$
The next table contains the roots of the rational function $`p(x)`$ in four particular cases.
Vector field $`D_0`$ Rational function $`p(x)`$ Solutions of $`p(x)=0`$ $`_0`$ $`1`$ none $`_1_0`$ $`x1`$ $`1`$ $`_12_0+_1`$ $`x2+x^1`$ $`1`$, $`1`$ $`_3_0`$ $`x^31`$ $`1`$, $`\frac{1+\sqrt{3}}{2}`$, $`\frac{1\sqrt{3}}{2}`$
If the vector field $`D_0`$ is $`_1_0`$ or $`_12_0+_1`$ then the region $`A_n`$ in dimension $`N=1`$ is just $`A_n=\{\alpha [\pi ,\pi )||\alpha |>\frac{1}{n}\}`$,
Suppose we are given a $`Y_0`$-closed function $`\xi `$ with Fourier expansion $`\xi =_I\widehat{\xi }_IH_I`$. We can actually turn our graph into a weighted graph by assigning to each directed edge $`(o(I),o(\tau I))`$ the weight $`\widehat{\xi }_I`$.
Figure 3. The region $`A_n`$ in dimension $`N=1`$.
Figure 4. The region $`A_n`$ in dimension $`N=2`$, as a subset of the square $`[\pi ,\pi ]^2`$.
Let $`c_n`$ be the unique $`L^2(^N)`$ solution of the equation $`(T_{D_0}c_n)=(\widehat{\xi })1_{A_n}`$, $`n1`$. The solution $`c_n`$ is defined through
$$c_n=^1\left(\frac{1}{p(e^{i(\alpha _1+\mathrm{}+\alpha _N)})}(\widehat{\xi })(\alpha )1_{A_n}(\alpha )\right).$$
The $`\stackrel{~}{\mathrm{\Sigma }}_N`$-invariance of the set $`A_n`$ and the fact that $`\widehat{\xi }`$ is constructed from the Fourier coefficients of a closed function implies that $`c_n`$ is $`\stackrel{~}{S}_N`$-invariant, $`n1`$.
Obviously we have the convergence $`(T_{D_0}c_n)\widehat{\xi }`$ as $`n\mathrm{}`$ in the topology of $`L^2([\pi ,\pi ]^N)`$, hence $`T_{D_0}c_n\widehat{\xi }`$ as $`n\mathrm{}`$ in the topology of $`L^2(^N)`$. โ
Proof of the decomposition theorem 2.1. Let $`\xi L^2(^{},d\nu _0^{gc})`$ be a $`D_0`$-closed function. From lemma 4.1 we know that $`\mathrm{Proj}__N\xi `$ is $`D_0`$-closed, for any $`N0`$. $`\mathrm{Proj}__0\xi `$ is a constant function. If the sum of the coefficients of $`D_0`$ is not equal to zero then $`\mathrm{Proj}__0\xi `$ is $`D_0`$-exact, otherwise $`\mathrm{Proj}__0\xi `$ is orthogonal on any $`D_0`$-exact function. Therefore, the decomposition theorem follows as long as we establish that any $`D_0`$-closed function $`\xi _N`$ can be approximated by $`D_0`$-exact functions, for any $`N1`$.
Assume $`\xi =_{I_N}\widehat{\xi }_IH_I_N`$, $`N1`$. Define the $`S_N`$-invariant function $`\widehat{\xi }:^N`$ through $`\widehat{\xi }(z)=\widehat{\xi }_I`$ if $`z\stackrel{S_N}{}z_I`$, see (3). We use lemmas 5.1 and 5.2 to find a sequence of finitely supported, $`\stackrel{~}{S}_N`$-invariant functions $`(c_n)_{n1}`$, such that $`T_{D_0}c_n\widehat{\xi }`$ as $`n\mathrm{}`$ in the topology of $`L^2(^N)`$. But lemmas 4.2 and 4.3 tell us that each of $`T_{D_0}c_n`$, $`n1`$ defines a $`D_0`$-exact function $`\xi _{c_n}`$, see (13). At the end of lemma 4.1 we noticed that the topology of $`_N`$ and $`L^2(^N)`$ are equivalent, hence we can claim that $`\xi _{c_n}\xi `$ as $`n\mathrm{}`$ in the Hilbert space topology of $`_N`$, or $`L^2(^{},d\nu _0^{gc})`$. โ
## 6. Second-order Ginzburg-Landau field and algebraic topology
We conclude with some remarks about the second-order Ginzburg-Landau field $`Y_0=_1_0`$ which has been studied in the work of S. R. S. Varadhan, . Our approach places Varadhanโs result in a new light by depicting an algebraic topologic aspect, to be explained below.
In section 3 we presented an extensive study of the set of multi-indices $``$. There we partitioned the set of multi-indices $``$ into disjoint orbits, and we denoted by $`๐ช`$ the space of orbits. Below we exhibit a procedure to construct a directed graph that has as vertices the orbits of the set of multi-indices.
A directed graph is a pair $`(V,E)`$ of two sets, where $`V`$ is the set of vertices of the graph and $`E`$ is the set of directed edges. A directed edge is a pair of two vertices $`(v_1,v_2)`$ where the first vertex indicates the starting point of the edge and the second vertex indicates the tip of the edge. For us we choose $`V`$ to be the set of orbits $`๐ช`$. Also we say that we have a directed edge $`(o_1,o_2)`$ if there exist a multi-index $`Io_1`$ such that $`\tau Io_2`$. Notice that if there exists an edge between two orbits then the orbits contain multi-indices with identical degrees. Hence our graph will have at least one connected component for each degree $`N0`$. We shall show that there exists precisely one connected component for each degree $`N0`$.
We would like to have a concrete or geometric presentation of the graph. For this purpose we use the identification of the set of orbits $`๐ช_N`$ containing the multi-indices of degree $`N`$, with the cone $`๐_N^+`$ of the lattice $`^N`$, $`N0`$.
Assume $`N=0`$, then $`๐ช_N`$ contains a single orbit, the orbit of the multi-index $`0`$ and this orbit contains a single multi-index. Since the multi-index $`0`$ has the property that $`0=\tau 0`$, we will have a directed edge going out of and returning to $`0`$; in other words we have a loop at $`0`$.
Assume that $`N1`$. It can be shown that a directed edge links $`z๐_N^+`$ to $`z^{}๐_N^+`$ if and only if either $`z^{}=ze_1\mathrm{}e_N`$ or $`z^{}=z+e_i`$ for some $`1iN`$. Here $`e_i`$ is lattice vector $`(0,\mathrm{},1,\mathrm{},0)`$ with the $`i`$th coordinate $`1`$. We shall indicate below how this presentation of the graph can be obtained.
Let $`o๐ช_N`$ be some orbit and $`R(o)=_{i=1}^ka_i\delta _{n_i}`$, with $`a_i1`$ and $`0n_1<n_2<\mathrm{}<n_k`$ be the representative of the orbit $`o`$. Given our rule, the orbit $`o`$ is connected by an edge going out of $`o`$ to each of the orbits $`o(\tau R(o))`$, $`o(\tau (n_1R(o)))`$, $`\mathrm{}`$, $`o(\tau (n_kR(o)))`$. For each of the orbits in the list before we can calculate the representatives and the corresponding point in the cone $`๐_N^+`$.
For example: the representative of the orbit $`o`$ is $`R(o)=_{i=1}^ka_i\delta _{n_i}`$ and the cone point is
$$z_{R(o)}=(\underset{a_1}{\underset{}{n_1,\mathrm{},n_1}},\mathrm{},\underset{a_k}{\underset{}{n_k,\mathrm{},n_k}}).$$
Assume $`n_11`$. The representative of the orbit of $`\tau R(o)`$ is $`\tau R(o)`$ and the corresponding cone point is
$$z_{\tau R(o)}=(\underset{a_1}{\underset{}{n_11,\mathrm{},n_11}},\mathrm{},\underset{a_k}{\underset{}{n_k1,\mathrm{},n_k1}}).$$
We notice that $`z_{\tau R(o)}=z_{R(o)}e_1\mathrm{}e_N`$. Also if $`n_1=0`$ the representative of the orbit of $`\tau R(o)`$ is $`(1)\tau R(o)`$ and the corresponding cone point is
$$z_{(1)\tau R(o)}=(\underset{a_11}{\underset{}{0,\mathrm{},0}},1,\mathrm{},\underset{a_k}{\underset{}{n_k1,\mathrm{},n_k1}}),$$
and $`z_{(1)\tau R(o)}=z_{R(o)}+e_{a_1}`$. Similarly, we can analyze the other orbits connected with $`o`$.
In particular our discussion proves that for any two given orbits $`o_1`$ and $`o_2`$ if there exists a multi-index $`Io_1`$ and $`\tau Io_2`$ then this multi-index is unique. As we will see later that this observation allows us to assign in a unique way a multi-index to any directed edge of our graph.
We include three pictures of the connected components of the directed graph for $`N=0`$, $`N=1`$ and $`N=2`$.
Figure 5. The connected component of the graph for $`N=0`$.
Figure 6. The connected component of the graph for $`N=1`$.
Figure 7. The connected component of the graph for $`N=2`$.
Suppose we are given a function $`\xi L^2(^{},d\nu _0^{gc})`$ with Fourier expansion $`\xi =_I\widehat{\xi }_IH_I`$. We can actually turn our directed graph into a weighted graph by assigning to each directed edge $`(o_1,o_2)`$ the Fourier coefficient $`\widehat{\xi }_I`$ corresponding to the unique multi-index $`I`$ such that $`Io_1`$ and $`\tau Io_2`$. Note that each Fourier coefficient will be assigned to one and only one edge and each edge will have assigned one and only one Fourier coefficient, since there is a one-to-one correspondence between the edges of our graph and the set $``$ of multi-indices. For example the edge $`(o(I),o(\tau I))`$ will have attached the weight $`\widehat{\xi }_I`$.
Figure 8. A directed edge and its attached weight.
It is interesting to note that if $`\xi `$ is $`Y_0`$-closed then the weights of the graph discussed before sum up to zero along any directed cycle of the graph except the loop of the connected component corresponding to $`N=0`$. Indeed the closedness condition of $`\xi `$ imposes no restriction on the coefficient $`\widehat{\xi }_0`$. Also note that the $`Y_0`$-closedness condition (10)
$$\widehat{\xi }_{I+\delta _{(n+1)}}\widehat{\xi }_{I+\delta _n}=\widehat{\xi }_{\tau ^n(I+\delta _1)}\widehat{\xi }_{\tau ^n(I+\delta _0)}n,I$$
plus the square-integrability of $`\xi `$ are equivalent to the property that the weights of the graph associated to $`\xi `$ sum up to zero around any directed cycle of any connected component corresponding to $`N1`$. However if $`\xi `$ is $`Y_0`$-exact then $`\widehat{\xi }_0=0`$ and hence the weights of the graph sum up to zero around any directed cycle of the graph. If $`\xi `$ is $`Y_0`$-exact and is constructed as in lemma 4.2 then we can say that in any connected component of the graph all but finitely many weights are zero.
The above can be explained from an algebraic topologic point of view. We can turn our directed graph into a $`2`$-dimensional $`\mathrm{\Delta }`$-complex (see A. Hatcherโs book on Algebraic Topology, ) by attaching enough discs to cycles of the graph such that each of the connected components $`N1`$ can be retracted to a single point. We do not attach a disc onto the loop of the connected component $`N=0`$. After the attaching process the $`2`$-dimensional $`\mathrm{\Delta }`$-complex can be retracted to the disjoint union of a circle with a countable number of points. The Fourier coefficients of a $`Y_0`$-exact function form a coboundary of our $`2`$-dimensional $`\mathrm{\Delta }`$-complex and the Fourier coefficients of a $`Y_0`$-closed function form a cocycle for our $`2`$-dimensional $`\mathrm{\Delta }`$-complex. Since the cohomology group $`H^1(C,)`$ of a circle $`C`$ is one-dimensional we expect the space of $`Y_0`$-exact functions to have codimension one inside the space of $`Y_0`$-closed functions.
## Acknowledgments
I would like to thank Professor George A. Elliott for his suggestion to use a Fourier analysis approach to solve the problem discussed in this paper. |
warning/0506/astro-ph0506550.html | ar5iv | text | # Diffraction Analysis of 2-D Pupil Mapping for High-Contrast Imaging
## 1 Introduction
Pupil mapping for the high-contrast imaging required by the problem of finding and imaging extra-solar terrestrial planets was first proposed by Guyon (2003). This idea has generated lots of excitement since it uses $`100\%`$ of the available light and exploits the full resolution of the optical system. Preliminary laboratory results were presented in Galicher et al. (2004).
In Traub and Vanderbei (2003) and Vanderbei and Traub (2005), we studied pupil mapping as a method for generating arbitrary pupil apodizations and, in particular, apodizations that provide the ultra-high contrast needed for terrestrial planet finding. By pupil mapping we mean a system of two lenses, or mirrors, that take a flat input field at the entrance pupil and produce an output field that is amplitude modified but still flat in phase (at least for on-axis sources).
Pupil mapping is easiest to describe in terms of ray-optics. An on-axis ray entering the first pupil at radius $`r`$ from the center is to be mapped to radius $`\stackrel{~}{r}=\stackrel{~}{R}(r)`$ at the exit pupil. Optical elements at the two pupils ensure that the exit ray is parallel to the entering ray. The function $`\stackrel{~}{R}(r)`$ is assumed to be positive and increasing or, sometimes, negative and decreasing. In either case, the function has an inverse that allows us to recapture $`r`$ as a function of $`\stackrel{~}{r}`$: $`r=R(\stackrel{~}{r})`$. The purpose of pupil mapping is to create nontrivial amplitude apodizations. An amplitude apodization function $`A(\stackrel{~}{r})`$ specifies the ratio between the output amplitude at $`\stackrel{~}{r}`$ to the input amplitude at $`r`$ (although we typically assume the input amplitude is a constant). We showed in Vanderbei and Traub (2005) that for any amplitude apodization function $`A(\stackrel{~}{r})`$ there is a pupil mapping function $`R(\stackrel{~}{r})`$ that achieves this amplitude profile. Specifically, the pupil mapping is given by
$$R(\stackrel{~}{r})=\pm \sqrt{_0^{\stackrel{~}{r}}2A^2(s)s๐s}.$$
(1)
Furthermore, if we consider the case of a pair of lenses that are plano on their outward-facing surfaces (as shown in Figure 1), then the inward-facing surface profiles, $`h(r)`$ and $`\stackrel{~}{h}(\stackrel{~}{r})`$, that are required to obtain the desired pupil mapping are given by the solutions to the following ordinary differential equations:
$$\frac{h}{r}(r)=\frac{r\stackrel{~}{R}(r)}{\sqrt{Q_0^2+(n^21)(r\stackrel{~}{R}(r))^2}},h(0)=z,$$
(2)
and
$$\frac{\stackrel{~}{h}}{\stackrel{~}{r}}(\stackrel{~}{r})=\frac{R(\stackrel{~}{r})\stackrel{~}{r}}{\sqrt{Q_0^2+(n^21)(R(\stackrel{~}{r})\stackrel{~}{r})^2}},\stackrel{~}{h}(0)=0.$$
(3)
Here, $`n`$ is the refractive index and $`Q_0`$ is a constant determined by the distance $`z`$ separating the centers ($`r=0`$, $`\stackrel{~}{r}=0`$) of the two lenses: $`Q_0=(n1)z`$.
Let $`S(r,\stackrel{~}{r})`$ denote the distance between a point on the first lens surface $`r`$ units from the center and the corresponding point on the second lens surface $`\stackrel{~}{r}`$ units from its center. Up to an additive constant, the optical path length of a ray that exits at radius $`\stackrel{~}{r}`$ after entering at radius $`r=R(\stackrel{~}{r})`$ is given by
$$Q_0(\stackrel{~}{r})=S(R(\stackrel{~}{r}),\stackrel{~}{r})+n(\stackrel{~}{h}(\stackrel{~}{r})h(R(\stackrel{~}{r}))).$$
(4)
In Vanderbei and Traub (2005), we showed that, for an on-axis source, $`Q_0(\stackrel{~}{r})`$ is constant and equal to $`Q_0`$.
## 2 High-Contrast Apodization
If we assume that an apodized beam with amplitude apodization profile $`A(\stackrel{~}{r})`$ such as one obtains as the output of a pupil mapping system is passed into an ideal imaging system with focal length $`f`$, the electric field $`E(\rho )`$ at the image plane is given by the Fourier transform of $`A(\stackrel{~}{r})`$:
$$E(\xi ,\eta )=\frac{E_0}{\lambda f}e^{2\pi i\frac{\stackrel{~}{x}\xi +\stackrel{~}{y}\eta }{\lambda f}}A(\sqrt{\stackrel{~}{x}^2+\stackrel{~}{y}^2})๐\stackrel{~}{y}๐\stackrel{~}{x}.$$
(5)
Here, $`E_0`$ is the input amplitude which, unless otherwise noted, we take to be unity. Since the optics are azimuthally symmetric, it is convenient to use polar coordinates. The apodization function $`A`$ is a function of $`\stackrel{~}{r}=\sqrt{\stackrel{~}{x}^2+\stackrel{~}{y}^2}`$ and the image-plane electric field depends only on image-plane radius $`\rho =\sqrt{\xi ^2+\eta ^2}`$:
$`E(\rho )`$ $`=`$ $`{\displaystyle \frac{1}{\lambda f}}{\displaystyle e^{2\pi i\frac{\stackrel{~}{r}\rho }{\lambda f}\mathrm{cos}(\theta \varphi )}A(\stackrel{~}{r})\stackrel{~}{r}๐\theta ๐\stackrel{~}{r}}`$ (6)
$`=`$ $`{\displaystyle \frac{2\pi }{\lambda f}}{\displaystyle J_0\left(2\pi \frac{\stackrel{~}{r}\rho }{\lambda f}\right)A(\stackrel{~}{r})\stackrel{~}{r}๐\stackrel{~}{r}}.`$ (7)
The point-spread function (PSF) is the square of the electric field:
$$\text{Psf}(\rho )=|E(\rho )|^2.$$
(8)
For the purpose of terrestrial planet finding, it is important to construct an apodization for which the PSF at small nonzero angles is ten orders of magnitude reduced from its value at zero. Figure 2 shows one such apodization function. In Vanderbei et al. (2003), we explain how these apodization functions are computed.
## 3 Huygens Wavelets
We have designed the pupil mapping system using simple ray optics but we have relied on diffraction theory to ensure that the apodization provides high contrast. This begs the question as to whether the desired high contrast will remain after a diffraction analysis of the entire system including the two-lens pupil mapping system or will diffraction effects in the pupil mapping system itself create โerrorsโ that are great enough to destroy the high-contrast that we seek. To answer this question, we need to do a diffraction analysis of the pupil mapping system itself.
If we assume that a flat, on-axis, electric field arrives at the entrance pupil, then the electric field at a particular point of the exit pupil can be well-approximated by superimposing the phase-shifted waves from each point across the entrance pupil (this is the well-known Huygens-Fresnel principleโsee, e.g., Section 8.2 in Born and Wolf (1999)). That is,
$$E_{\text{out}}(\stackrel{~}{x},\stackrel{~}{y})=\frac{1}{\lambda Q(\stackrel{~}{x},\stackrel{~}{y},x,y)}e^{2\pi iQ(\stackrel{~}{x},\stackrel{~}{y},x,y)/\lambda }๐y๐x,$$
(9)
where
$$Q(\stackrel{~}{x},\stackrel{~}{y},x,y)=\sqrt{(x\stackrel{~}{x})^2+(y\stackrel{~}{y})^2+(h(r)\stackrel{~}{h}(\stackrel{~}{r}))^2}+n(Zh(r)+\stackrel{~}{h}(\stackrel{~}{r}))$$
(10)
is the optical path length, $`Z`$ is the distance between the plano lens surfaces (i.e., a constant slightly larger than $`z`$), and where, of course, we have used $`r`$ and $`\stackrel{~}{r}`$ as shorthands for the radii in the entrance and exit pupils, respectively. As before, it is convenient to work in polar coordinates:
$$E_{\text{out}}(\stackrel{~}{r})=\frac{1}{\lambda Q(\stackrel{~}{r},r,\theta )}e^{2\pi iQ(\stackrel{~}{r},r,\theta )/\lambda }r๐\theta ๐r,$$
(11)
where
$$Q(\stackrel{~}{r},r,\theta )=\sqrt{r^22r\stackrel{~}{r}\mathrm{cos}\theta +\stackrel{~}{r}^2+(h(r)\stackrel{~}{h}(\stackrel{~}{r}))^2}+n(Zh(r)+\stackrel{~}{h}(\stackrel{~}{r})).$$
(12)
For numerical tractability, it is essential to make approximations so that the integral over $`\theta `$ can be carried out analytically. To this end, we need to make an appropriate approximation to the square root term:
$$S=\sqrt{r^22r\stackrel{~}{r}\mathrm{cos}\theta +\stackrel{~}{r}^2+(h(r)\stackrel{~}{h}(\stackrel{~}{r}))^2}.$$
(13)
## 4 Fresnel Propagation
In this section we consider approximations that lead to the so-called Fresnel propagation formula.
If we assume that the lens separation is fairly large, then the $`(h\stackrel{~}{h})^2`$ term dominates the rest and so we can use the first two terms of a Taylor series approximation (i.e., for $`u`$ small relative to $`a`$, $`\sqrt{u^2+a^2}a+u^2/2a`$) to get the following large separation approximation:
$$S(h(r)\stackrel{~}{h}(\stackrel{~}{r}))+\frac{r^22r\stackrel{~}{r}\mathrm{cos}\theta +\stackrel{~}{r}^2}{2(h(r)\stackrel{~}{h}(\stackrel{~}{r}))}.$$
(14)
If we assume further that the lenses are thin (i.e., that $`n`$ is large), then $`h\stackrel{~}{h}`$ in the denominators can be approximated simply by $`z`$:
$$S(h(r)\stackrel{~}{h}(\stackrel{~}{r}))+\frac{r^22r\stackrel{~}{r}\mathrm{cos}\theta +\stackrel{~}{r}^2}{2z}.$$
(15)
This is called the thin lens approximation.
Combining the large lens separation approximation with the thin lens approximation, we get
$$E_{\text{out}}(\stackrel{~}{r})=\frac{1}{\lambda Q_1(\stackrel{~}{r},r,\theta )}e^{2\pi iQ_1(\stackrel{~}{r},r,\theta )/\lambda }r๐\theta ๐r$$
(16)
where
$$Q_1(\stackrel{~}{r},r,\theta )=\frac{r^22r\stackrel{~}{r}\mathrm{cos}\theta +\stackrel{~}{r}^2}{2z}+Z+(n1)(Zh(r)+\stackrel{~}{h}(\stackrel{~}{r})).$$
(17)
Finally, we simplify the reciprocal of $`Q_1`$ by noting that the $`Z`$ term dominates the other terms (i.e., for $`u`$ small relative to $`a`$, $`1/(u+a)1/a`$) and so we get that:
$$\frac{1}{Q_1(\stackrel{~}{r},r,\theta )}\frac{1}{Z}.$$
(18)
This last approximation is called the paraxial approximation.
Combining all three approximations, we now arrive at the standard Fresnel approximation:
$$E_{\text{out}}(\stackrel{~}{r})=\frac{2\pi }{\lambda Z}e^{\pi i\frac{\stackrel{~}{r}^2}{z\lambda }+2\pi i\frac{(n1)\stackrel{~}{h}(\stackrel{~}{r})}{\lambda }}e^{\pi i\frac{r^2}{z\lambda }2\pi i\frac{(n1)h(r)}{\lambda }}J_0(2\pi r\stackrel{~}{r}/z\lambda )r๐r.$$
(19)
While the standard Fresnel approximation works very well in most conventional situations, it turns out (as well shall show) to be too crude of an approximation for high-contrast pupil mapping. It is inadequate because it does not honor the constancy of the optical path length $`Q(\stackrel{~}{r},r,\theta )`$ along the rays of ray-optics. That is, the fact that $`Q(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$ is constant has been lost in the approximations. We should have used the ray-tracing optical path length as the โlarge quantityโ in our large-lens-separation approximation instead of the simpler difference $`h(r)\stackrel{~}{h}(\stackrel{~}{r})`$. But, this seemingly simple adjustment quickly gets tedious and so we prefer to take a completely different (and simpler) approach, which is described in the next section.
## 5 An Alternative to Fresnel
As we have just explained and shall demonstrate later, the standard Fresnel approximation does not produce good results for high-contrast pupil mapping computations. In this section, we present an alternative approximation that is slightly more computationally demanding but is much closer to a direct calculation of the true Huygens wavelet propagation.
As with Fresnel, we approximate the $`1/Q(\stackrel{~}{r},r,\theta )`$ amplitude-reduction factor in (11) by the constant $`1/Z`$ (the paraxial approximation). The $`Q(\stackrel{~}{r},r,\theta )`$ appearing in the exponential must, on the other hand, be treated with care. Recall that $`Q(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$ is a constant. Since constant phase shifts are immaterial, we can subtract it from $`Q(\stackrel{~}{r},r,\theta )`$ in (11) to get
$$E_{\text{out}}(\stackrel{~}{r})\frac{1}{\lambda Z}e^{2\pi i(Q(\stackrel{~}{r},r,\theta )Q(\stackrel{~}{r},R(\stackrel{~}{r}),0))/\lambda }r๐\theta ๐r.$$
(20)
Next, we write the difference in $`Q`$โs as follows:
$`Q(\stackrel{~}{r},r,\theta )Q(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$ $`=`$ $`S(\stackrel{~}{r},r,\theta )S(\stackrel{~}{r},R(\stackrel{~}{r}),0)+n(h(R(\stackrel{~}{r}))h(r))`$ (21)
$`=`$ $`{\displaystyle \frac{S^2(\stackrel{~}{r},r,\theta )S^2(\stackrel{~}{r},R(\stackrel{~}{r}),0)}{S(\stackrel{~}{r},r,\theta )+S(\stackrel{~}{r},R(\stackrel{~}{r}),0)}}+n(h(R(\stackrel{~}{r}))h(r))`$ (22)
and then we expand out the numerator and cancel big terms that can be subtracted one from another to get
$`S^2(\stackrel{~}{r},r,\theta )S^2(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$ $`=`$ $`(rR(\stackrel{~}{r}))(r+R(\stackrel{~}{r}))2\stackrel{~}{r}\left(r\mathrm{cos}\theta R(\stackrel{~}{r})\right)`$ (23)
$`+\left(h(r)h(R(\stackrel{~}{r}))\right)\left(h(r)+h(R(\stackrel{~}{r}))2\stackrel{~}{h}(\stackrel{~}{r})\right).`$
When $`r=R(\stackrel{~}{r})`$ and $`\theta =0`$, the right-hand side clearly vanishes as it should. Furthermore, for $`r`$ close to $`R(\stackrel{~}{r})`$ and $`\theta `$ close to zero, the right-hand side gives an accurate formula for computing the deviation from zero. That is to say, the right-hand side is easy to program in such a manner as to avoid subtracting one large number from another, which is always the biggest danger in numerical computation.
So far, everything is exact (except for the paraxial approximation). The only further approximation we make is to replace $`S(\stackrel{~}{r},r,\theta )`$ in the denominator of (22) with $`S(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$ so that the denominator becomes just $`2S(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$. Putting this altogether and replacing the integral on $`\theta `$ with the appropriate Bessel function, we get a new approximation, which we refer to as the Huygens approximation:
$`E_{\text{out}}(\stackrel{~}{r})`$ $``$ $`{\displaystyle \frac{2\pi }{\lambda Z}}{\displaystyle e^{2\pi i\left(\frac{(rR(\stackrel{~}{r}))(r+R(\stackrel{~}{r}))+2\stackrel{~}{r}R(\stackrel{~}{r})+\left(h(r)h(R(\stackrel{~}{r}))\right)\left(h(r)+h(R(\stackrel{~}{r}))2\stackrel{~}{h}(\stackrel{~}{r})\right)}{2S(\stackrel{~}{r},R(\stackrel{~}{r}),0)}+n(h(R(\stackrel{~}{r}))h(r))\right)/\lambda }}`$ (24)
$`\times J_0\left(2\pi \stackrel{~}{r}r/\lambda S(\stackrel{~}{r},R(\stackrel{~}{r}),0)\right)rdr.`$
## 6 Sanity Checks
In this section we consider a number of examples.
### 6.1 Flat glass windows ($`A1`$)
We begin with the simplest example in which the apodization function is identically equal to one. Taking the positive root in (1), we get that $`R(\stackrel{~}{r})=\stackrel{~}{r}`$. That is, the ray-optic design is for the light to go straight through the system. The inverse map $`\stackrel{~}{R}(r)`$ is also trivial: $`\stackrel{~}{R}(r)=r`$. Hence, the right-hand sides in the differential equations (2) and (3) vanish and the lens figures become flat: $`h(r)z`$ and $`\stackrel{~}{h}(\stackrel{~}{r})0`$. In this case, $`Q_1(\stackrel{~}{r},R(\stackrel{~}{r}),0)`$ is a constant (independent of $`\stackrel{~}{r}`$) and so the Fresnel approximation is a good one. The Fresnel results, shown in Figure 3 should match textbook examples for simple open circular aperturesโand they do.
### 6.2 A Galilean Telescope ($`Aa>1`$)
In this case, we ask for a system in which the output pupil has been uniformly amplitude intensified by a factor $`a>1`$. If we choose the positive root in (1), then we get a Galilean-style refractor telescope consisting of a convex lens at the entrance pupil and a concave lens as an eyepiece. Specifically, we get $`R(\stackrel{~}{r})=a\stackrel{~}{r}`$ and $`\stackrel{~}{R}(r)=r/a`$. From these it follows that if the aperture of the first lens is $`D`$, then the aperture of the second is $`\stackrel{~}{D}=D/a`$. It is easy to compute the lens figures
$`h(r)`$ $`=`$ $`z+{\displaystyle \frac{\sqrt{Q_0^2+(n^21)(11/a)^2r^2}|Q_0|}{(n^21)(11/a)}}`$ (25)
$`\stackrel{~}{h}(\stackrel{~}{r})`$ $`=`$ $`{\displaystyle \frac{\sqrt{Q_0^2+(n^21)(a1)^2\stackrel{~}{r}^2}|Q_0|}{(n^21)(a1)}}.`$ (26)
If the relative index of refraction $`n`$ is greater than $`1`$ (as in air-spaced glass lenses), then these functions represent portions of a hyperbola. If, on the other hand, $`n<1`$ (as in a glass medium between the two surfaces), then the functions are ellipses. Fresnel results for $`a=3`$ and $`n=1.5`$ are shown in Figure 4. Note the large error in the phase map and the fact that the computed PSF does not follow the usual Airy pattern. This is strong evidence that the Fresnel approximation is too crude since real systems of this sort exist and exhibit the expected Airy pattern.
In Figure 5 we consider the same system but use the Huygens approximation. Note that the phase map, while not perfect, is now much flatter. Also, the computed PSF is closely matches the expected Airy pattern.
### 6.3 An Ideal Lens
If we let $`a`$ tend to infinity, we see that $`\stackrel{~}{D}`$ tends to zero and the system reduces to a convex lens focusing a collimated input beam to a point. In this case, the second lens vanishes; only the first lens is of interest. Its equation is
$$h(r)=z+\frac{\sqrt{Q_0^2+(n^21)r^2}|Q_0|}{(n^21)}.$$
(27)
The plane where the second lens was is now the image plane. To compute the electric field here, we put $`\stackrel{~}{h}(\stackrel{~}{r})0`$ and use either the Fresnel or the Huygens approximation. Since we would put a detector at this plane, we can ignore any final phase corrections and both approximations reduce to the same formula:
$$E_{\text{out}}(\stackrel{~}{r})=\frac{2\pi }{\lambda Z}e^{\pi i\frac{r^2}{z\lambda }2\pi i\frac{(n1)h(r)}{\lambda }}J_0(2\pi r\stackrel{~}{r}/z\lambda )r๐r.$$
(28)
Substituting (27) into (28) and dropping any unit complex numbers that factor out of the integral, we get
$$E_{\text{out}}(\stackrel{~}{r})=\frac{2\pi }{\lambda Z}e^{\pi i\frac{r^2}{z\lambda }\frac{2\pi i}{\lambda }\sqrt{\left(\frac{n1}{n+1}\right)^2z^2+\frac{n1}{n+1}r^2}}J_0(2\pi r\stackrel{~}{r}/z\lambda )r๐r.$$
(29)
Of course, this formula is for a uniform collimated input beam. If the input beam happens to be apodized by some upstream optical element, then the expression becomes
$$E_{\text{out}}(\stackrel{~}{r})=\frac{2\pi }{\lambda Z}e^{\pi i\frac{r^2}{z\lambda }\frac{2\pi i}{\lambda }\sqrt{\left(\frac{n1}{n+1}\right)^2z^2+\frac{n1}{n+1}r^2}}J_0(2\pi r\stackrel{~}{r}/z\lambda )A(r)r๐r$$
(30)
where $`A(r)`$ denotes the apodization function. This formula does not agree with the Fourier transform expression given earlier by equation (5). However, if the square root is approximated by the first two terms of its Taylor expansion,
$$\sqrt{\left(\frac{n1}{n+1}\right)^2z^2+\frac{n1}{n+1}r^2}=\frac{n1}{n+1}z+\frac{r^2}{2z},$$
(31)
then (30) reduces to a Fourier transform as in (5) (again dropping unit complex factors).
This raises an interesting question: if a high-contrast apodization is designed based on the assumption that the focusing element behaves like a Fourier transform (i.e., as in (5)), how well will the apodization work if the true expression for the electric field is closer to the one given by (30)? The answer is shown in Figure 6. The PSF degradation is very small.
## 7 High-Contrast Apodization
The purpose of the examples discussed in the previous section was to convince the reader that the Huygens approximation provides a reasonable estimate of the electric field at the exit pupil of a pupil mapping system. Assuming we were convincing, we now proceed to apply the Huygens approximation to compute the electric field and focused PSF of the pupil mapping system corresponding to the apodization function shown in Figure 2. As before, we assume a wavelength of $`0.6328\mu `$ and lenses with aperture $`D=25`$mm. In Figure 7, we show the results for $`z=15D`$ and a refraction index of $`1.5`$. For these simulations to be meaningful, it is critical that the integrals be represented by a sum over a sufficiently refined partition. The bigger the disparity between wavelength $`\lambda `$ and aperture $`D`$, the more refined the partition needs to be. For the parameters we have chosen, a partition into $`5000`$ parts proved to be adequate and is what we have chosen. The plot in the upper-left section of the figure shows in gray the target amplitude apodization profile and in black the amplitude profile computed using the Huygens approximation (i.e., equation (24)). The plot in the upper right shows the lens profiles. The first lens is shown in black and the second in gray. The plot in the lower left shows in gray the computed optical path length $`Q_0(\stackrel{~}{r})`$. If the numerical computation of the lens figures had been done with insufficient precision, this curve would not be flat. As we see, it is flat. The lower left also shows in black the phase map as computed by the Huygens approximation. Note that here there are high frequency oscillations everywhere and low frequency oscillation that has an amplitude that increases as one moves out to the rim of the lens. The lower-right plot shows in gray the PSF associated with the ideal amplitude apodization and in black the PSF computed by Huygens propagation.
The PSF in Figure 7 is disappointing. It is important to determine whether this is real or is a result of the approximations behind the Huygens propagation formula. As a check, we did a brute force computation of the Huygens integral (11). Because this integral is more difficult, we were forced to use only 500 $`r`$-values and 500 $`\theta `$-values. Hence, we had to increase the wavelength by a factor of $`10`$. With these changes, the result is shown in Figure 8. It too shows the same amplitude and phase oscillations. This sanity check convinces us that these effects are physical. We need to consider changes to the physical setup that might mollify these oscillations. Such changes are considered next.
There is no particular reason to make the two lenses have equal aperture. By scaling the apodization function, we can easily generate examples with unequal aperture. One such experiment we tried was to scale the apodization function so that its value at the center is one. This scaling results in the second lens having almost four times the aperture of the first lens. The results for this case are shown in Figure 9.
If an apodization designed for $`10^{10}`$ contrast only produces $`10^5`$, one wonders how well an apodization designed for $`10^5`$ will do. The answer is shown in Figure 11. In this case, the degradation due to diffraction effects is very small.
## 8 Hybrid Systems
If a pupil mapping system designed for $`10^5`$ works well, perhaps it could be followed by a conventional apodizer that attempts to bring the system from $`10^5`$ down to $`10^{10}`$. We tried this. The result is shown in Figure 12. As can be seen in the lower-right plot, the contrast achieved is limited to about $`10^{7.5}`$. Apparently the diffraction ripples going into the apodizer are enough to prevent the system from achieving the desired contrast.
As an alternative to a downstream apodizer, one could consider a pre-apodizer placed in front of the first lens. Since it is generally hard edges that create bad diffraction effects, we can imagine using the pre-apodizer to provide the near-outer-edge apodization and allow the pupil mapping system to provide the main body of the apodization. In this way, perhaps the diffraction effects can be minimized while at the same time maintaining a system with high throughput. The results of one such experiment are shown in Figure 13.
Acknowledgements. This research was partially performed for the Jet Propulsion Laboratory, California Institute of Technology, sponsored by the National Aeronautics and Space Administration as part of the TPF architecture studies and also under JPL subcontract number 1260535. The first author received support from the NSF (CCR-0098040) and the ONR (N00014-98-1-0036). |
warning/0506/math0506401.html | ar5iv | text | # An ergodic action of the outer automorphism group of a free group
## Introduction
Let $`F_n`$ be a free group of rank $`n>1`$ and let $`G`$ be a compact Lie group. Then $`\mathrm{๐ง๐๐}(F_n,G)`$ admits a natural volume form which is invariant under $`\mathrm{๐ ๐๐}(F_n)`$. This volume form descends to a finite measure on the character variety $`\mathrm{๐ง๐๐}(F_n,G)/G`$ which is invariant under $`\mathrm{๐ฎ๐๐}(F_n)`$. The purpose of this note is to prove:
###### Theorem.
Suppose that $`G`$ is a connected group locally isomorphic to a product of copies of $`\mathrm{๐ฒ๐ด}(2)`$ and $`๐ด(1)`$. If $`n>2`$, then the $`\mathrm{๐ฎ๐๐}(F_n)`$-action on $`\mathrm{๐ง๐๐}(F_n,G)/G`$ is ergodic.
We conjecture that $`\mathrm{๐ฎ๐๐}(F_n)`$ is ergodic on each connected component of $`\mathrm{๐ง๐๐}(F_n,G)/G`$ for every compact Lie group $`G`$ and $`n>2`$.
When $`G=๐ด(1)`$, then this action is just the action of $`\mathrm{๐ฆ๐ซ}(n,)`$ on the $`n`$-torus $`^n/^n`$, which is well known to be ergodic. In fact, certain cyclic subgroups of $`\mathrm{๐ฆ๐ซ}(n,)`$ act ergodicly.
The proof relies heavily on , both in its outline and a key result. When $`n=2`$, the action is not ergodic, since it preserves the function
$`\mathrm{๐ง๐๐}(F_n,G)/G`$ $`\stackrel{๐
}{}[2,2]`$
$`[\rho ]`$ $`\mathrm{๐๐}([X_1,X_2])`$
where $`X_1,X_2`$ are a pair of free generators for $`F_2`$. However, for each $`2t2`$, the action is ergodic on $`\kappa ^1(t)`$.
When $`\pi `$ is the fundamental group of a closed surface, then Pickrell and Xia have proved $`\mathrm{๐ฎ๐๐}(\pi )`$ is ergodic on $`\mathrm{๐ง๐๐}(\pi ,G)/G`$ for any compact Lie group $`G`$.
As in , the methods here apply when $`G`$ is any Lie group having simple factors $`๐ด(1)`$ and $`\mathrm{๐ฒ๐ด}(2)`$. In particular, since this class of groups is closed under the operation of taking direct products, the action of $`\mathrm{๐ฎ๐๐}(F_n)`$ is ergodic on
$$\mathrm{๐ง๐๐}(F_n,G\times G)/(G\times G)\mathrm{๐ง๐๐}(F_n,G)/G\times \mathrm{๐ง๐๐}(F_n,G)/G.$$
As in , the action of $`\mathrm{๐ฎ๐๐}(F_n)`$ on $`\mathrm{๐ง๐๐}(F_n,G)/G`$ is weak-mixing, that is:
###### Corollary.
The only invariant finite-dimensional subrepresentation of the induced unitary representation of $`\mathrm{๐ฎ๐๐}(F_n)`$ on $`L^2(\mathrm{๐ง๐๐}(F_n,G)/G)`$ consists of constants.
I would like to thank David Fisher for pointing out an error in the original proof of Lemma 3.1 and for many helpful suggestions.
## 1. Ergodic theory of the $`\mathrm{๐ฒ๐ด}(2)`$-character variety
Let $`\{X_1,\mathrm{},X_n\}`$ be a set of free generators for $`F_n`$ and let
$$X_0=X_n^1\mathrm{}X_1^1.$$
Then $`F_n`$ is the fundamental group of an $`n+1`$-holed sphere $`S_{n+1}`$, where the $`X_0,X_1,\mathrm{},X_n`$ correspond to components of $`S_{n+1}`$. The mapping class group $`\mathrm{\Gamma }_{n+1}`$ of $`S_{n+1}`$ embeds in $`\mathrm{๐ฎ๐๐}(F_n)`$ as the subgroup preserving the conjugacy classes of the cyclic subgroups $`X_i`$ for $`i=0,\mathrm{},n`$.
The proof proceeds as follows. Let
(1.1)
$$\mathrm{๐ง๐๐}(F_n,G)/G\stackrel{๐}{}$$
be an $`\mathrm{๐ฎ๐๐}(F_n)`$-invariant measurable function. We show that $`f`$ is constant almost everywhere.
The main result of applied to the surface $`S_{n+1}`$ gives the following:
###### Proposition 1.1.
The mapping
$`\mathrm{๐ง๐๐}(F_n,G)/G`$ $`\stackrel{t_{}}{}[2,2]^{n+1}`$
$`[\rho ]`$ $`\left[\begin{array}{c}\mathrm{๐๐}(\rho (X_0))\\ \mathrm{}\\ \mathrm{๐๐}(\rho (X_n))\end{array}\right]`$
is an ergodic decomposition for the action of $`\mathrm{\Gamma }_{n+1}`$. That is, for every $`\mathrm{\Gamma }_{n+1}`$-invariant measurable function
$$\mathrm{๐ง๐๐}(F_n,G)/G\stackrel{}{}$$
there exists a measurable function
$$[2,2]^{n+1}\stackrel{๐ป}{}$$
such that $`h=Ht_{}`$ almost everywhere.
Using the embedding of the mapping class group
$$\mathrm{\Gamma }_{n+1}\mathrm{๐ฎ๐๐}(F_n)$$
as above, the $`\mathrm{๐ฎ๐๐}(F_n)`$-invariant function $`f`$ is $`\mathrm{\Gamma }_{n+1}`$-invariant, and hence factors through $`t_{}`$.
By Proposition 1.1 there exists a function
(1.2)
$$[2,2]^{n+1}\stackrel{๐น}{}.$$
such that $`f=Ft_{}`$, where $`f`$ is the function discussed in (1.1).
## 2. The case of rank $`n=3`$
First consider the case $`n=3`$. Following the notation of , denote the generators by
$$A=X_1,B=X_2,C=X_3,D=X_0$$
so that $`A,B,C,D`$ are subject to the relation
(2.1)
$$ABCD=1.$$
A representation $`\rho `$ is determined by its values
$$(\rho (A),\rho (B),\rho (C))G^3$$
on the generators $`A,B,C`$ and
$$\rho (D)=\rho (C)^1\rho (B)^1\rho (A)^1.$$
### 2.1. Trace coordinates
The equivalence class $`[\rho ]`$ is determined by the seven functions
$`a`$ $`=\mathrm{๐๐}(\rho (A))`$
$`b`$ $`=\mathrm{๐๐}(\rho (B))`$
$`c`$ $`=\mathrm{๐๐}(\rho (C))`$
$`d`$ $`=\mathrm{๐๐}(\rho (D))=\mathrm{๐๐}(\rho (D^1))=\mathrm{๐๐}(\rho (ABC))`$
$`x`$ $`=\mathrm{๐๐}(\rho (AB))`$
$`y`$ $`=\mathrm{๐๐}(\rho (BC))`$
$`z`$ $`=\mathrm{๐๐}(\rho (CA))`$
subject to the polynomial relation
(2.2) $`x^2+`$ $`y^2+z^2+xyz=`$
$`(ab+cd)x+(ad+bc)y+(ac+bd)z`$
$`+(4a^2b^2c^2d^2abcd).`$
In other words, the $`\mathrm{SL}(2,)`$-character variety of $`F_3`$ is the hypersurface in $`^7`$ defined by (2.2).
When $`a,b,c,d`$ the topology of the set of $``$-points is analyzed in . In particular the $`\mathrm{๐ฒ๐ด}(2)`$-character variety is the union over the set $`V`$ of all
$$(a,b,c,d)[2,2]^4$$
satisfying
$`0\mathrm{\Delta }(a,b,c,d)`$ $`=\left(2(a^2+b^2+c^2+d^2)abcd16\right)^2`$
$`(4a^2)(4b^2)(4c^2)(4d^2)`$
of compact components of the cubic surface in $`^3`$ satisfying (2.2). Here is an alternate description with which it is easier to work.
### 2.2. Rank two free groups
The $`\mathrm{๐ฒ๐ด}(2)`$-character variety of $`F_2`$ is the subset $`V_3^3`$ defined by traces
$$(x_1,x_2,x_3)[2,2]^3$$
satisfying the inequality
(2.3)
$$x_1^2+x_2^2+x_3^2+x_1x_2x_34$$
which is depicted in Figure 1.
A quadruple $`(a,b,c,d)[2,2]^4`$ is the image of an $`\mathrm{๐ฒ๐ด}(2)`$-character if and only if there exists $`y`$ such that both triples $`(a,d,y)`$ and $`(b,c,y)`$ lie in $`V_3`$.
Using (2.2), we determine the condition that $`(a,d,y)V_3`$. Apply (2.2) to $`x_1=y,x_2=a,x_3=d`$ to see that $`(a,d,y)V_3`$ if and only if $`y`$ lies in the interval
$$Y(a,d):=[y_{}(a,d),y_+(a,d)]$$
with endpoints
$$y_\pm (a,d):=\frac{ad\pm \sqrt{(4a^2)(4d^2)}}{2}.$$
For any $`(a,d)[2,2]^2`$, the interval $`Y(a,d)`$ is nonempty, so that the restriction of the projection
$`V_3`$ $`\stackrel{\mathrm{\Pi }_{a,d}}{}[2,2]^2`$
$`\left[\begin{array}{c}a\\ b\\ c\\ d\\ x\\ y\\ z\end{array}\right]`$ $`\left[\begin{array}{c}a\\ d\end{array}\right]`$
is onto. Furthermore the fiber $`V_3(a,d)`$ of the surjection
$`t_{}(V_3)`$ $`\stackrel{\mathrm{\Pi }_{a,d}}{}[2,2]^2`$
$`\left[\begin{array}{c}a\\ b\\ c\\ d\end{array}\right]`$ $`\left[\begin{array}{c}a\\ d\end{array}\right]`$
consists of all $`(b,c)[2,2]^2`$ such that
$$Y(b,c)Y(a,d)\mathrm{}.$$
If $`2<y<2`$, then the set of $`(b,c)`$ such that $`yY(b,c)`$ is the closed elliptical region $`\overline{}_y`$ inscribed in the square $`[2,2]^2`$ at the four points
$$(2,y),(y,2),(2,y),(y,2),$$
depicted in Figure 3. If $`y=\pm 2`$, then the set of $`(b,c)`$ such that $`yY(b,c)`$ is the line segment
$$\{(b,b)2b2\}.$$
For fixed $`a,d[2,2]`$, the fiber $`\mathrm{\Pi }^1\left(V_3(a,d)\right)`$ equals
$$\underset{yY(a,d)}{}\{(b,c)yY(b,c)\},$$
depicted in Figure 4.
### 2.3. A non-geometric automorphism
The automorphism $`\alpha \mathrm{๐ ๐๐}(F_3)`$ defined by:
$`A`$ $`\stackrel{\alpha }{}A`$
$`B`$ $`\stackrel{\alpha }{}BA^1`$
$`C`$ $`\stackrel{\alpha }{}AC`$
$`D`$ $`\stackrel{\alpha }{}D`$
induces the following automorphism of the character variety:
$$\left[\begin{array}{c}a\\ b\\ c\\ d\\ x\\ y\\ z\end{array}\right]\stackrel{\alpha ^{}}{}\left[\begin{array}{c}a\\ x\\ acz\\ d\\ axb\\ y\\ c\end{array}\right]$$
Since $`\alpha ^{}`$ preserves the coordinates $`a,d,y`$, this automorphism restricts to a diffeomorphism on the level sets $`a=a_0,d=d_0`$ and $`y=y_0`$ and the restriction is linear:
$$\left[\begin{array}{c}x\\ b\\ c\\ z\end{array}\right]\stackrel{\alpha ^{}}{}\left[\begin{array}{cccc}a_0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& a_0& 1\\ 0& 0& 1& 0\end{array}\right]\left[\begin{array}{c}x\\ b\\ c\\ z\end{array}\right]$$
The following elementary fact (whose proof is omitted) is useful:
###### Lemma 2.1.
Let $`2<a<2`$. The linear transformation
$`^2`$ $`\stackrel{L_a}{}^2`$
$`\left[\begin{array}{c}x\\ y\end{array}\right]`$ $`\left[\begin{array}{c}axy\\ x\end{array}\right]`$
preserves the positive definite quadratic form
$$\left[\begin{array}{c}x\\ y\end{array}\right]\stackrel{Q}{}x^2axy+y^2$$
and lies in the linear flow generated by the vector field
$$\mathrm{{\rm Y}}:=\left(\frac{a}{2}xy\right)_x+\left(x\frac{a}{2}y\right)_y.$$
The trajectories of $`\mathrm{{\rm Y}}`$ are the level sets of $`Q`$, which are ellipses. $`L_a`$ is linearly conjugate to a rotation by angle
(2.4)
$$\theta =2\mathrm{cos}^1(a/2).$$
If $`\theta \pi `$, then $`L_a`$ has infinite order and is ergodic on each $`Q`$-level set.
The $`4`$-dimensional affine subspaces of $`^7`$ corresponding to levels of $`(a_0,d_0,y_0)`$ split as products of two affine $`2`$-planes (corresponding to levels of $`(x,b)`$ and $`(c,z)`$ respectively). Evidently the linear map $`\alpha ^{}`$ on these $`4`$-planes splits as a direct sum of two copies of the linear map $`L_{a_0}`$. It preserves the trajectories of the linear vector field
$`๐`$ $`=\left({\displaystyle \frac{a_0}{2}}xb\right)_x+\left(x{\displaystyle \frac{a_0}{2}}b\right)_b`$
$`+\left({\displaystyle \frac{a_0}{2}}cz\right)_c+\left(c{\displaystyle \frac{a_0}{2}}z\right)_z.`$
The zeroes of this vector field consist of the origin
$$\left[\begin{array}{c}x\\ b\\ c\\ z\end{array}\right]=\left[\begin{array}{c}0\\ 0\\ 0\\ 0\end{array}\right]$$
and, when $`a_0=\pm 2`$, the squares defined by
$$x=\pm b,z=\pm c.$$
All other trajectories are ellipses. The transformation $`\alpha ^{}`$ acts by rotation along these ellipses through angle $`\theta `$ given by (2.4). When $`\theta /\pi `$ is irrational, this action is ergodic. Thus, for almost every $`a_0[2,2]`$, the restriction of $`\alpha ^{}`$ to these ellipses is ergodic. On a set of full measure, the function $`f`$ is constant along the projections of trajectories of $`๐`$.
Fix $`(a_0,d_0)(2,2)^2`$ and consider the equivalence relation $``$ on $`V_3(a_0,d_0)`$ generated by the projections of trajectories of $`๐`$.
###### Lemma 2.2.
For $`a_0,d_0\pm 2`$, all points in $`V_3(a_0,d_0)`$ are $``$-equivalent.
###### Proof.
Since $`V_3(a_0,d_0)`$ is connected, it suffices to prove that each equivalence class is open. To this end, suppose that $`(b_0,c_0)V_3(a_0,d_0)`$; we show that every $`(b,c)`$ sufficiently close to $`(b_0,c_0)`$ is equivalent to $`(b_0,c_0)`$.
The image of the tangent vector $`๐`$ at $`(x,b,c,z)`$ under the differential of the coordinate projection
$$\left[\begin{array}{c}x\\ b\\ c\\ z\end{array}\right]\stackrel{\mathrm{\Pi }_{b,c}}{}\left[\begin{array}{c}b\\ c\end{array}\right]$$
is
$$\left(\mathrm{\Pi }_{b,c}\right)_{}๐=\left(x\frac{a_0}{2}b\right)_b+\left(\frac{a_0}{2}cz\right)_c$$
The fiber $`\mathrm{\Pi }_{b,c}^1(b_0,c_0)`$ is an interval. For some (and hence almost every) $`(x,z)\mathrm{\Pi }_{b,c}^1(b_0,c_0)`$, the vector
$$\left(\mathrm{\Pi }_{b,c}\right)_{}๐(b_0,x,z,c_0)$$
is nonzero. Choose such an $`(x_0,z_0)`$. For any open neighborhood $`U`$ of $`(x_0,z_0)`$, the values $`\left(\mathrm{\Pi }_{b,c}\right)_{}๐(b_0,x,z,c_0)`$ for $`(x,z)U`$, span $`^2`$.
Let $`\mathrm{\Phi }_t`$ denote the flow generated by $`๐`$ and choose an open neighborhood $`U_0`$ of $`(b_0,x_0,z_0,c_0)`$ in $`V_3(a_0,d_0)`$. The differential of the map
$`\times \left(\{b_0\}\times U_0\times \{c_0\}\right)`$ $`^2`$
$`(t,b_0,x,z,c_0)`$ $`\mathrm{\Pi }_{b,c}\left(\mathrm{\Phi }_t(b_0,x,z,c_0)\right)`$
at $`(0,b_0,x_0,z_0,c_0)`$ is onto. The inverse function theorem guarantees an open neighborhood of $`(0,b_0,x_0,z_0,c_0)`$ mapping onto an open neighborhood of $`(b_0,c_0)`$, as desired. โ
Thus, for almost every $`(a_0,d_0)[2,2]^2`$, the function $`F`$ of (1.2) is constant along the level surfaces $`\mathrm{\Pi }_{a,d}^1(a_0,d_0)`$, and hence factors through the projection $`\mathrm{\Pi }_{a,d}`$:
$$F(a,b,c,d)=F(a,d).$$
Applying the same argument to the automorphism
$`A`$ $`\stackrel{\gamma }{}CA`$
$`B`$ $`\stackrel{\gamma }{}B`$
$`C`$ $`\stackrel{\gamma }{}C`$
$`D`$ $`\stackrel{\gamma }{}DC^1`$
implies that $`F`$ factors through the projection $`\mathrm{\Pi }_{(b,c)}`$ and $`F`$ is almost everywhere constant.
Thus the function $`f`$, which is invariant under the mapping class group $`\mathrm{\Gamma }_4`$ and the automorphisms $`\alpha _{},\gamma _{}`$, must be constant almost everywhere. This completes the proof of the Theorem when $`n=3`$.
## 3. General rank
The case of rank $`n>3`$ follows easily from the case $`n=3`$. The following elementary lemma is useful:
###### Lemma 3.1.
Let $`G`$ be a compact Lie group. Then $`\mathrm{๐ฎ๐๐}(F_n)`$ is ergodic on $`\mathrm{๐ง๐๐}(F_n,G)/G`$ if and only if $`\mathrm{๐ ๐๐}(F_n)`$ is ergodic on $`\mathrm{๐ง๐๐}(F_n,G)`$. More generally, let
$$\mathrm{๐ง๐๐}(F_n,G)\stackrel{\Pi }{}\mathrm{๐ง๐๐}(F_n,G)/G$$
denote the quotient map and suppose $`\mathrm{\Gamma }\mathrm{๐ ๐๐}(F_n)`$ is a subgroup containing $`\mathrm{๐จ๐๐}(F_n)`$. A measurable $`\mathrm{\Gamma }`$-invariant mapping.
$$\mathrm{๐ง๐๐}(F_n,G)/G\stackrel{๐}{}W$$
is an ergodic decomposition for $`\mathrm{\Gamma }`$ if and only if
$$\mathrm{๐ง๐๐}(F_n,G)\stackrel{f\mathrm{\Pi }}{}W$$
is an ergodic decomposition for $`\mathrm{\Gamma }`$.
###### Proof.
For almost every $`\rho \mathrm{๐ง๐๐}(F_n,G)`$, the image $`\rho (F_n)`$ is dense in $`G`$. Thus $`\mathrm{๐จ๐๐}(F_n)`$ is ergodic on each $`G`$-orbit in $`\mathrm{๐ง๐๐}(F_n,G)`$, and the quotient map $`\mathrm{\Pi }`$ is an ergodic decomposition for the $`\mathrm{๐จ๐๐}(F_n)`$-action. Since $`\mathrm{๐จ๐๐}(F_n)\mathrm{\Gamma }`$, every $`\mathrm{\Gamma }`$-invariant function
$$\mathrm{๐ง๐๐}(F_n,G)\stackrel{}{}W$$
(where $`W`$ is a standard measure space) factors through $`\mathrm{\Pi }`$. Thus composition with $`\mathrm{\Pi }`$ induces an isomorphism between the set of equivalence classes of $`\mathrm{\Gamma }`$-invariant measurable maps on $`\mathrm{๐ง๐๐}(F_n,G)`$ and on $`\mathrm{๐ง๐๐}(F_n,G)/G`$. In particular the ergodic decompositons for the respective $`\mathrm{\Gamma }`$-actions are $`\mathrm{\Pi }`$-related, as desired. โ
In the following lemma, let $`\iota _j`$ (for $`j=1,\mathrm{},n`$) denote the monomorphism defined by
$`F_{n1}`$ $`\stackrel{\iota _j}{}F_n`$
$`X_i`$ $`\{\begin{array}{cc}X_i\text{ if }i<j\hfill & \\ X_{i+1}\text{ if }ij\hfill & \end{array}`$
for $`i=1,\mathrm{},n1`$. Denote by
$$\mathrm{๐ ๐๐}(F_{n1)}\stackrel{I_j}{}\mathrm{๐ ๐๐}(F_n)$$
the homomorphism defined by
$$X_i\stackrel{I_j(\varphi )}{}\{\begin{array}{cc}X_j\text{ if }i=j\hfill & \\ \iota _j(\varphi (\iota _j^1(X_i)))\text{ if }ij\hfill & \end{array}$$
for $`\varphi \mathrm{๐ ๐๐}(F_n)`$. Clearly $`\iota _j`$ is equivariant with respect to $`I_j`$.
###### Lemma 3.2.
Let $`n>3`$. Suppose that $`\mathrm{๐ ๐๐}(F_{n1})`$ is ergodic on $`\mathrm{๐ง๐๐}(F_{n1},G)`$. Let $`1jn`$. Then
$`\mathrm{๐ง๐๐}(F_n,G)`$ $`\stackrel{t_j}{}[2,2]`$
$`\rho `$ $`\mathrm{๐๐}\left(\rho (X_j)\right)`$
is an ergodic decomposition for the action of $`I_j(\mathrm{๐ ๐๐}(F_{n1}))\times \mathrm{๐จ๐๐}(G)`$.
###### Proof.
Clearly $`t_j`$ is $`I_j(\mathrm{๐ ๐๐}(F_{n1})`$-invariant. Let $`2<\tau <2`$. The fiber $`t_j^1(\tau )`$ identifies with the set of equivalence classes of
$$((\rho (X_1),\mathrm{},\rho (X_{n1}),\rho (X_n))$$
where $`\mathrm{๐๐}\left(\rho (X_j)\right)=\tau `$, that is, the set of $`\mathrm{๐จ๐๐}(G)`$-orbits of
$$\mathrm{๐ง๐๐}(F_{n1},G)\times \mathrm{๐๐}^1(\tau ).$$
By Lemma 3.1 and the ergodicity hypothesis, $`\mathrm{๐ ๐๐}(F_{n1})`$ is ergodic on $`\mathrm{๐ง๐๐}(F_{n1},G)`$. Transitivity of the $`\mathrm{๐จ๐๐}(G)`$-action on $`\mathrm{๐๐}^1(\tau )`$ implies ergodicity of $`\mathrm{๐ ๐๐}(F_{n1})\times \mathrm{๐จ๐๐}(G)`$ acting on the subset of $`\mathrm{๐ง๐๐}(F_n,G)`$ corresponding to $`t_j^1(\tau )`$, which is equivalent to ergodicity of $`\mathrm{๐ฎ๐๐}(F_{n1})`$ on $`t_j1(\tau )`$. โ
Suppose inductively that $`\mathrm{๐ ๐๐}(F_{n1})`$ is ergodic on $`\mathrm{๐ง๐๐}(F_{n1},G)`$. By Lemma 3.2, any $`I_j(\mathrm{๐ ๐๐}(F_{n1}))`$-invariant measurable function $`f`$ factors through $`t_j`$ for each $`j=1,\mathrm{},n`$. We need only consider $`j=1`$ and $`j=n`$. Since
$$\mathrm{๐ง๐๐}(F_n,G)/G\stackrel{(t_1,t_n)}{}[2,2]^2$$
is onto, any two points $`\rho ,\rho ^{}`$ in a full measure subset of $`\mathrm{๐ง๐๐}(F_n,G)`$ can be joined by a third $`\rho ^{\prime \prime }`$ such that
$`t_1(\rho ^{\prime \prime })`$ $`=t_1(\rho ),`$
$`t_n(\rho ^{\prime \prime })`$ $`=t_n(\rho ^{}).`$
Let $`f`$ be an $`\mathrm{๐ ๐๐}(F_n)`$-invariant measurable function on $`\mathrm{๐ง๐๐}(F_n,G)/G`$. Factorization of $`f`$ through $`t_1`$ implies that
$$f(\rho )=f(\rho ^{\prime \prime })$$
and factorization of $`f`$ through $`t_n`$ implies that
$$f(\rho ^{\prime \prime })=f(\rho ^{}).$$
Thus $`f(\rho )=f(\rho ^{\prime \prime })`$, and $`f`$ is constant almost everywhere, as desired. โ
## 4. Other Lie groups
The extension of the proof to products of $`\mathrm{๐ฒ๐ด}(2)`$ and $`๐ด(1)`$ proceeds exactly as in . Away from a null subset of $`\mathrm{๐ง๐๐}(F_n,G)/G`$, the action of the automorphisms of $`F_n`$ preserve tori (products of the $`๐ด(1)`$ factors and the invariant ellipses). Furthermore on almost every torus, the action is ergodic, and hence every invariant function factors through the trace functions of the generators. The rest of the proof is completely analogous. |
warning/0506/hep-lat0506026.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Rare decays mediated by flavour-changing neutral-currents (FCNC) are among the deepest probes to uncover the fundamental mechanism of quark flavour mixing. Within the Standard Model (SM), these rare decays are strongly suppressed both by the GIM mechanism and by the hierarchy of the CKM matrix , and are often dominated by short-distance dynamics. As a result, FCNC processes are very sensitive to possible new sources of flavour mixing, even if these occur well above the electroweak scale. The sensitivity to physics beyond the SM of these rare processes is closely related to the theoretical accuracy on which we are able to compute their amplitudes within the SM.
Within the family of FCNC decays, long-distance effects are not always negligible and, in most cases, they represent the dominant source of theoretical uncertainty. Long-distance contributions are typically relevant in: i) amplitudes where the GIM mechanism is only logarithmic; ii) amplitudes where the power-like GIM suppression of the long-distance component is partially compensated by a large CKM coefficient. So far, the evaluation of these non-perturbative contributions has been performed by means of effective theories. These analytic tools require the introduction of additional parameters, the knowledge of which constitutes a source of sizable theoretical uncertainty.
In this paper we show that for a class of very interesting processes, such as $`K^+\pi ^+\nu \overline{\nu }`$ and $`K\pi \mathrm{}^+\mathrm{}^{}`$, it is possible in principle to compute non-perturbatively the long-distance contribution to the physical amplitudes on the lattice. The physical information is encoded in the following $`T`$-products:
$`๐ฏ_{Q,J}^\mu (q^2)=๐ฉ_V{\displaystyle d^4xd^4ye^{iqy}\pi |T[Q(x)J^\mu (y)]|K},`$ (1)
where $`Q`$ denotes a generic four-quark operator of the effective weak Hamiltonian, $`J^\mu `$ is either the electromagnetic or the weak neutral current, and $`๐ฉ_V`$ is an appropriate volume factor. If the invariant mass of the lepton pair ($`q^2`$) is smaller than any physical hadronic threshold, the calculation proceeds as in the case of semileptonic form factors (see e.g. refs. ), and one obtains directly the relevant amplitude. When instead the leptonic invariant mass exceeds the pion threshold, the final state interaction induces problems similar to those encountered with non-leptonic kaon decays . However, the knowledge of the amplitude for $`q^2<m_\pi ^2`$ is sufficient to determine the leading unknown effective couplings of these amplitudes within the framework of chiral perturbation theory (CHPT) . Therefore, the combination of lattice calculations and CHPT should allow to reach an unprecedented level of precision for these rare decays.
When using a lattice action with explicit chiral symmetry breaking, such as Wilson, Clover or twisted mass fermions, further problems arise because of additional ultraviolet (power) divergences which may appear in the operator matrix elements or in the relevant $`T`$-products.<sup>1</sup><sup>1</sup>1 Alternative formulations which guarantee chiral symmetry in the physical matrix elements, such as overlap fermions, do not have this problem . However, these formulations are not mature yet to be used for unquenched calculations of these complicated matrix elements, for quark masses close to the physical values. We show that for the electromagnetic current, gauge invariance prevents the appearance of these divergences even if the most popular lattice actions are used. Consequently, when $`J^\mu `$ is the electromagnetic current, the $`T`$-products in eq. (1) are finite provided that a renormalized weak effective Hamiltonian is used. The situation is slightly more complicated when $`J^\mu `$ is the weak neutral current. In this case, simple power counting, related to the behavior of the $`T`$-product at short distances, shows that both quadratic and linear divergences may appear. We show that the quadratic divergence, which is not a peculiarity of the lattice regularization, is canceled by the GIM mechanism. Concerning the linear divergence, which is present only if there is an explicit chiral symmetry breaking term in the lattice action, we demonstrate that it can be avoided by using the maximally twisted mass fermion action .
There is a further subtlety concerning the ambiguity in the renormalization of the effective weak Hamiltonian out of the chiral limit . In ref. it has been shown that this ambiguity does not affect the physical $`K\pi \pi `$ amplitudes, but is present in โnon-physicalโ matrix elements, such as $`\pi |Q|K`$. This problem is present also in our case and implies an ambiguity in the $`T`$-products of eq. (1). By means of appropriate Ward Identities, we show that the physical amplitude, the extraction of which requires a specific spectral analysis discussed in the following, is instead free of ambiguities.
The paper is organized as follows: in sect. 2 we recall the basic ingredients of radiative decays in the framework of the effective Hamiltonian approach. In sect. 3 we describe the strategy for computing the relevant amplitudes from the Euclidean Green functions and discuss the structure of the divergences in both cases: when they cancel because of gauge invariance and when it is necessary to get rid of them using GIM mechanism and twisted mass. In sect. 4 we show how to extract the physical amplitude in spite of the ambiguity of the renormalized effective Hamiltonian. The results are summarized in the conclusions.
## 2 Effective Hamiltonian for $`K\pi \mathrm{}^+\mathrm{}^{}(\nu \overline{\nu })`$ decays
The dimension-six effective Hamiltonian relevant to evaluate $`sd\mathrm{}^+\mathrm{}^{}(\nu \overline{\nu })`$ amplitudes at next-to-leading order accuracy, renormalized at a scale $`M_W\mu >m_c`$, can be written as
$$_{eff}=_{eff}^{|\mathrm{\Delta }S|=1}+_{eff}^{\mathrm{FCNC}}+\frac{G_F}{\sqrt{2}}\underset{q=u,d,s,c}{}Q_q^{\mathrm{NC}}+\frac{G_F}{\sqrt{2}}\underset{q=u,cq^{}=d,s}{}V_{ij}Q_{qq^{}}^{\mathrm{CC}}+\text{h.c.},$$
(2)
where $`V_{ij}`$ denote the elements of the CKM matrix,
$`_{eff}^{|\mathrm{\Delta }S|=1}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{us}^{}V_{ud}\left[{\displaystyle \underset{i=1,2}{}}C_i\left(Q_i^uQ_i^c\right)+{\displaystyle \underset{i=3\mathrm{}8}{}}C_iQ_i+๐ช\left({\displaystyle \frac{V_{ts}^{}V_{td}}{V_{us}^{}V_{ud}}}\right)\right],`$ (3)
is the usual $`|\mathrm{\Delta }S|=1`$ weak Hamiltonian, for which the Wilson coefficients are known at the NLO , and
$`_{eff}^{\mathrm{FCNC}}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \frac{\alpha }{2\pi \mathrm{sin}^2\theta _W}}V_{us}^{}V_{ud}\left[{\displaystyle \underset{i=7V,7A,\nu }{}}C_iQ_i+๐ช\left({\displaystyle \frac{V_{ts}^{}V_{td}}{V_{us}^{}V_{ud}}}\right)\right].`$ (4)
Here
$`Q_{qq^{}}^{\mathrm{CC}}`$ $`=`$ $`\overline{q}\gamma ^\mu (1\gamma _5)q^{}\overline{\nu }\gamma _\mu (1\gamma _5)\mathrm{}`$
$`Q_q^{\mathrm{NC}}`$ $`=`$ $`\overline{q}\gamma ^\mu \left[2\widehat{T}(1\gamma _5)4\widehat{Q}\mathrm{sin}^2\theta _W\right]q`$ (5)
$`\times \left[\overline{\nu }\gamma _\mu (1\gamma _5)\nu \overline{\mathrm{}}\gamma _\mu (1\gamma _54\mathrm{sin}^2\theta _W)\mathrm{}\right]`$
are the charged-current and neutral-current effective interactions obtained by the integration of the heavy $`W`$ and $`Z`$ fields,
$`Q_{7V}`$ $`=`$ $`\overline{s}\gamma ^\mu (1\gamma _5)d\overline{\mathrm{}}\gamma _\mu \mathrm{},`$ (6)
$`Q_{7A}`$ $`=`$ $`\overline{s}\gamma ^\mu (1\gamma _5)d\overline{\mathrm{}}\gamma _\mu \gamma _5\mathrm{},`$ (7)
$`Q_\nu `$ $`=`$ $`\overline{s}\gamma ^\mu (1\gamma _5)d\overline{\nu }\gamma _\mu (1\gamma _5)\nu ,`$ (8)
are the leading FCNC operators, and
$`Q_1^q`$ $`=`$ $`\overline{s}_\alpha \gamma ^\mu (1\gamma _5)q_\beta \overline{q}_\beta \gamma _\mu (1\gamma _5)d_\alpha ,`$ (9)
$`Q_2^q`$ $`=`$ $`\overline{s}_\alpha \gamma ^\mu (1\gamma _5)q_\alpha \overline{q}_\beta \gamma _\mu (1\gamma _5)d_\beta `$ (10)
the leading four-quark operators. The four-quark operators originated by penguin contractions are denoted by $`Q_{1\mathrm{}6}`$, whereas $`Q_7`$ and $`Q_8`$ correspond to magnetic and chromomagnetic operators, respectively (see e.g. ref. ).
Thanks to both the GIM mechanism and the unitarity of the CKM matrix, the contributions to the FCNC amplitudes can be unambiguously decomposed into two parts, the first one proportional to the CKM combination $`V_{us}^{}V_{ud}`$, the second one proportional to $`V_{ts}^{}V_{td}`$. Since $`|V_{ts}^{}V_{td}||V_{us}^{}V_{ud}|`$, the contribution proportional to $`V_{ts}^{}V_{td}`$ is negligible but for cases where it is enhanced by the large top-quark mass (i.e. for amplitudes which exhibit a power-like GIM mechanism). In these cases, the amplitudes are completely dominated by short distances (top-quark loops) and can be evaluated in perturbation theory to an excellent degree of approximation. In this paper instead we are interested only in the long-distance components of the amplitudes, therefore we can safely work in the limit $`V_{td}=0`$.
We can seemingly neglect the matrix elements of $`Q_{1\mathrm{}6}`$ and $`Q_8`$ in the evaluation of $`K\pi \mathrm{}^+\mathrm{}^{}(\nu \overline{\nu })`$ amplitudes: these matrix elements vanish at the tree level and the corresponding Wilson coefficients are substantially smaller than those of $`Q_{12}^{u,c}`$. In this approximation, we only have to consider the contributions of the leading FCNC operators in eqs. (6)โ(8) and the non-trivial contractions of $`Q_{1,2}^{u,c}`$ with the electromagnetic current and the currents defined by $`Q_{qq^{}}^{\mathrm{CC}}`$ and $`Q_q^{\mathrm{NC}}`$. The $`K\pi `$ matrix elements of the FCNC operators in eqs. (6)โ(8) can be extracted from data on the leading $`K_\mathrm{}3`$ modes using isospin symmetry , or even computed directly on the lattice, with high accuracy, as recently shown in .<sup>2</sup><sup>2</sup>2 In principle, in the $`K\pi \mathrm{}^+\mathrm{}^{}`$ case one should also consider the tree-level matrix element of the magnetic operator $`Q_7=m_s\overline{s}\sigma ^{\mu \nu }(1\gamma _5)dF_{\mu \nu }`$, which cannot be directly extracted from $`K_\mathrm{}3`$ data. However, within the SM the smallness of the corresponding Wilson coefficient makes this contribution negligible for practical purposes. This matrix element can be computed on the lattice with standard techniques, as shown in . Concerning the contractions of $`Q_{1,2}^{u,c}`$, those with a charged current receive very small non-perturbative contributions (estimated to be below $`1\%`$ at the amplitude level in the $`K^+\pi ^+\nu \overline{\nu }`$ case and even smaller in all the other channels), which can be reliably estimated within CHPT . Thus the main problem are the contractions of $`Q_{1,2}^{u,c}`$ with a neutral current, as outlined in eq. (1).
So far, this problem has been addressed with the following two-step procedure: i) integrating out the charm as dynamical degree of freedom; ii) constructing the chiral realization of the corresponding effective Hamiltonian with light quarks only. This procedure suffers from two sources of theoretical errors: slow convergence of perturbation theory because of the low renormalization scale of the effective Hamiltonian ($`\mu <m_c`$); uncertainties associated to the new low-energy couplings appearing in the effective theory. Both these sources of uncertainties are naturally reduced in the lattice approach, where the effective Hamiltonian is renormalized above the charm scale and the $`T`$-products are evaluated in full QCD.
We now discuss separately electromagnetic and neutrino amplitudes in more detail.
### 2.1 $`K\pi \mathrm{}^+\mathrm{}^{}`$
The main non-perturbative correlators relevant for these decays are those with the electromagnetic current. In particular, the relevant $`T`$-product in Minkowski space is
$`\left(๐ฏ_i^j\right)_{\mathrm{em}}^\mu (q^2)`$ $`=`$ $`i{\displaystyle d^4xe^{iqx}\pi ^j(p)|T\left\{J_{\mathrm{em}}^\mu (x)\left[Q_i^u(0)Q_i^c(0)\right]\right\}|K^j(k)},`$ (11)
$`J_{\mathrm{em}}^\mu `$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \underset{q=u,c}{}}\overline{q}\gamma ^\mu q{\displaystyle \frac{1}{3}}{\displaystyle \underset{q=d,s}{}}\overline{q}\gamma ^\mu q`$ (12)
for $`i=1,2`$ and $`j=+,0`$. Thanks to gauge invariance we can write
$$\left(๐ฏ_i^j\right)_{\mathrm{em}}^\mu (q^2)=\frac{w_i^j(q^2)}{(4\pi )^2}\left[q^2(k+p)^\mu (m_k^2m_\pi ^2)q^\mu \right].$$
(13)
The normalization of (13) is such that the $`O(1)`$ scale-independent low-energy couplings $`a_{+,0}`$ defined in can be expressed as
$$a_j=\frac{1}{\sqrt{2}}V_{us}^{}V_{ud}\left[C_1w_1^j(0)+C_2w_2^j(0)+\frac{2N_j}{\mathrm{sin}^2\theta _W}f_+(0)C_{7V}\right].$$
(14)
where $`f_+`$ is the $`K\pi `$ vector form factor and $`\{N_+,N_0\}=\{1,2^{1/2}\}`$ . To a good approximation, the decay rates of the CP-conserving transitions $`K^+\pi ^+\mathrm{}^+\mathrm{}^{}`$ and $`K_S\pi ^0\mathrm{}^+\mathrm{}^{}`$ are proportional to the square of these effective couplings :
$$(K^+\pi ^+e^+e^{})6.6a_+^2\times 10^7,(K_S\pi ^0e^+e^{})10.4a_0^2\times 10^9.$$
(15)
At present, we are not able to predict $`a_{+,0}`$ with sufficient accuracy: we simply fit their $`๐ช(1)`$ values from the measured rates of the corresponding decay modes (an updated numerical analysis can be found in ). Being completely dominated by long distance contributions, these two CP-conserving processes would provide an excellent testing ground for the lattice technique.
On the other hand, the calculation of $`a_0`$ from first principles would have a very interesting phenomenological application in the $`K_L\pi ^0\mathrm{}^+\mathrm{}^{}`$ case, which proceeds via a CP-violating amplitude: the calculation of $`a_0`$ would allow to determine in a model-independent way the sign of the interference between the (long-distance) indirect-CP-violating component of the amplitude and the interesting (short-distance) direct-CP-violating term . This result would allow to perform a very precise test of direct-CP-violation in the kaon sector.
### 2.2 $`K\pi \nu \overline{\nu }`$
The power-like GIM mechanism of the leading electroweak amplitude, implies a severe suppression of long-distance effects in these modes. In the CP-violating channel, $`K_L\pi ^0\nu \overline{\nu }`$, long-distance contributions are negligible well below the $`1\%`$ level . However, this is not the case for the charged channel, $`K^+\pi ^+\nu \overline{\nu }`$, where the suppression of long-distance effects is partially compensated by a large CKM coefficient. The $`T`$-product which determines the size of non-perturbative effects in this mode is
$$\left(๐ฏ_i^+\right)_Z^\mu (q^2)=id^4xe^{iqx}\pi ^+(p)|T\left\{J_Z^\mu (x)\left[Q_i^u(0)Q_i^c(0)\right]\right\}|K^+(k),$$
(16)
where $`J_Z^\mu =\overline{q}\gamma ^\mu (2\widehat{T}(1\gamma _5)4\widehat{Q}\mathrm{sin}^2\theta _W)q`$ is the neutral current defined by $`Q_q^{\mathrm{NC}}`$ in (5). Separating the electromagnetic component, we can write $`\left(๐ฏ_i^+\right)_Z^\mu =\left(๐ฏ_i^+\right)_L^\mu 4\mathrm{sin}^2\theta _W\left(๐ฏ_i^+\right)_{\mathrm{em}}^\mu `$, where
$`\left(๐ฏ_i^+\right)_L^\mu (q^2)`$ $`=`$ $`i{\displaystyle d^4xe^{iqx}\pi ^+(p)|T\left\{J_L^\mu (x)\left[Q_i^u(0)Q_i^c(0)\right]\right\}|K^+(k)},`$ (17)
$`J_L^\mu `$ $`=`$ $`{\displaystyle \underset{q=u,c}{}}\overline{q}\gamma ^\mu (1\gamma _5)q{\displaystyle \underset{q=d,s}{}}\overline{q}\gamma ^\mu (1\gamma _5)q.`$ (18)
Contrary to $`\left(๐ฏ_i^+\right)_{\mathrm{em}}^\mu `$, the structure of $`\left(๐ฏ_i^+\right)_L^\mu `$ is not protected by gauge invariance and we can decompose it as
$$\left(๐ฏ_i^+\right)_\mathrm{L}^\mu (q^2)=\frac{m_K^2}{\pi ^2}\left[z_i^+(q^2)(k+p)^\mu +๐ช(q^\mu )\right],$$
(19)
where the normalization is such that the $`z_i^+(q^2)`$ are expected to be $`๐ช(1)`$ . The value of these form factors at $`q^2=0`$ is sufficient to control the long-distance contributions to the $`K^+\pi ^+\nu \overline{\nu }`$ amplitude down to the 1% level of precision .
Charm and, more generally, long-distance contributions to the $`K^+\pi ^+\nu \overline{\nu }`$ amplitude, are usually parametrized in terms of a scale-independent coefficient $`P_c`$ . According to the decomposition (19), this can be be written as
$$P_c=\frac{1}{|V_{us}|^4}\left\{\frac{m_K^2}{M_W^2}\left[C_1z_1^+(0)+C_2z_2^+(0)\right]+f_+(0)C_\nu \right\}.$$
(20)
The coefficient $`P_c`$ expresses the relative weight of the subleading terms relative to the top-quark amplitude, which is the leading contribution and is precisely determined in perturbation theory . As can be noted, the non-perturbative parameters $`z_i^+(0)`$ appear in (20) multiplied by a very small coefficient: $`m_K^2/M_W^2/|V_{us}|^40.015`$. Thus even a determination of these matrix elements at the 30โ50% level from lattice QCD would be sufficient to reduce the overall error on the $`K^+\pi ^+\nu \overline{\nu }`$ rate around or below the 1โ2% level.
## 3 $`T`$-products at short-distances on the lattice
In this section we discuss the properties of the Euclidean Green functions necessary to extract the physical amplitudes defined in eqs. (11) and (17) in a numerical simulation. Since the ultraviolet behavior is quite different in the two cases, we discuss them separately, starting from the $`T`$-product which involves the electromagnetic current. In both cases, we assume that the operators of the effective weak Hamiltonian have been renormalized, namely that all their physical matrix elements are finite as the lattice spacing goes to zero ($`a0`$). The renormalization of the effective Hamiltonian is discussed in the next section.
The starting point to extract the physical matrix elements is the following Euclidean Green function
$`\left(๐ฏ_i\right)_\mathrm{X}^\mu (q^2,t_\pi ,t_K)={\displaystyle d^4x\mathrm{\Phi }_\pi (t_\pi ,\stackrel{}{p})J_X^\mu (0)\left[Q_i^u(x)Q_i^c(x)\right]\mathrm{\Phi }_K^{}(t_K,\stackrel{}{k})},`$ (21)
$`t_\pi >0,t_K<0,`$
where the source (sink) for creating (annihilating) the pseudoscalar mesons at fixed space momentum are defined as
$$\mathrm{\Phi }_i(t_i,\stackrel{}{q}_i)=d^3ze^{i\stackrel{}{q}_i\stackrel{}{z}}\mathrm{\Phi }_i(t_i,\stackrel{}{z}),$$
(22)
and $`\mathrm{\Phi }_i(t_i,\stackrel{}{z})`$ is a suitable local operator with the quantum numbers of the pion or kaon, respectively. Note that, in order to simplify the notation and the comparison between continuum and lattice formulae, we use the symbol of integral also to indicate sums over the lattice sites.
If not for the presence of the weak four-fermion operator, the calculation would proceed as for the standard weak and electromagnetic form factors, by studying the behavior of the Green functions at large $`t_\pi `$ and $`|t_K|`$ . This would give the form factors computed at momentum transfer $`\stackrel{}{q}=\stackrel{}{k}\stackrel{}{p}`$ and with energy transfer $`q_0=E_KE_\pi `$. Since $`Q_i`$ is summed over the whole lattice volume and hence it carries zero momentum, this general strategy remains valid also for the Green function in eq. (21). As explained in the previous section, in order to extract the relevant low energy couplings, we are interested only to study the correlation function for $`q^2<m_\pi ^2`$. In this range no rescattering of intermediate states is possible and thus we do not have problems in relating the Minkowskian $`T`$ product to the Euclidean one.
The additional problem which arises in this case is the possibility that the Green function itself diverges because of the short distance behavior when $`x0`$. By dimensional arguments, this divergence can at most be quadratic. At fixed lattice spacing $`a`$, this would imply potential contributions to the Green function of $`๐ช(1/a^2)`$. Fortunately this never happens, since the strongest divergence associated to the diagram in figure 1 is independent of the quark masses and is canceled by the GIM mechanism. However, this cancellation does not guarantee the absence of linear divergences, which are naturally present when using lattice actions which break explicitly chiral invariance.
### 3.1 The electromagnetic current
Even if the chirality of the fermion action is explicitly broken, we are still able to define a conserved vector current on the lattice, which we can identify with the electromagnetic one. For example, with Wilson fermions we have
$`\widehat{J}_V^\mu ={\displaystyle \frac{1}{2}}\left[\overline{q}(x+\mu )U^\mu (x)(r+\gamma ^\mu )q(x)\overline{q}(x)U^\mu (x)(r\gamma ^\mu )q(x+\mu )\right],`$ (23)
where $`U^\mu `$ is the link variable. With a conserved current, gauge invariance is strong enough to protect the Green functions from the appearance of both quadratic and linear divergences. This remains true even when the Wick contractions correspond to a vacuum polarization diagram of the type in figure 1, where only one of the two currents is the lattice conserved one, and the other is a local vector current originating from the weak four-fermion operator. We have verified this argument by an explicit perturbative calculation using Wilson, Clover and twisted mass fermions. Since the results of this calculation (more precisely of the subdiagram in figure 2) could be useful for other applications, we give them below for the Wilson and Clover cases.
The amplitude we have considered is
$`\mathrm{\Pi }_{\mu \nu }(p)`$ $`=`$ $`{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{Tr}\left[\mathrm{\Gamma }_\nu ^{(1)}(q;p+q)\mathrm{\Delta }(q)\gamma _\mu \mathrm{\Delta }(q+p)\right]}`$ (24)
$`=`$ $`{\displaystyle \frac{8}{(4\pi )^2}}(\delta _{\mu \nu }p^2p_\mu p_\nu )\left\{(p^2a^2,m^2a^2)+L\right\},`$
where $`\mathrm{\Gamma }_\nu ^{(1)}(q;p+q)`$ is the vertex derived from eq. (23) and $`\mathrm{\Delta }(q)`$ the fermion propagator . Both in Wilson and Clover cases we can identify a universal infrared term, given by
$`(p^2a^2,m^2a^2)={\displaystyle _0^1}๐xx(1x)\mathrm{log}\left[m^2a^2+p^2a^2x(1x)\right],`$ (25)
while the finite constant $`L`$ depends from the details of the regularization. In the Wilson case we find
$`L_W={\displaystyle \frac{1}{6}}\mathrm{log}\left(m^2a^2\right)+(1\delta _{\rho \sigma }){\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4q}{2\pi ^2}}{\displaystyle \frac{\frac{1}{2}\mathrm{cos}q_\rho \mathrm{sin}^2q_\sigma \frac{1}{3}\mathrm{cos}q_\rho \mathrm{cos}2q_\sigma }{\mathrm{\Delta }_F^2(q;m)}}+`$
$`+{\displaystyle \frac{r^2\left(\frac{1}{3}\mathrm{cos}q_\rho \mathrm{cos}q_\sigma _\tau \mathrm{sin}^2\frac{q_\tau }{2}\frac{1}{6}\mathrm{cos}q_\sigma \mathrm{sin}^2q_\rho \frac{1}{3}\mathrm{cos}q_\rho \mathrm{sin}^2q_\sigma \right)}{\mathrm{\Delta }_F^2(q;m)}}`$
while in the Clover case
$`L_{\mathrm{Cl}}=L_\mathrm{W}+{\displaystyle \frac{r^2}{2\pi ^2}}(1\delta _{\rho \sigma }){\displaystyle _\pi ^\pi }d^4q{\displaystyle \frac{\frac{1}{2}\mathrm{sin}^2q_\rho \mathrm{cos}q_\sigma _\tau \mathrm{sin}^2\left(\frac{q_\tau }{2}\right)}{\mathrm{\Delta }_F^2(q)}}.`$ (26)
Note that in both cases the absence of power divergences holds independently from the GIM mechanism. Using the above results we could then match the lattice calculation with the continuum one even in an effective theory where the charm quark is integrated out. The comparison of the results obtained with or without dynamical charm quarks would provide a useful insight about the validity of the standard effective theory obtained by renormalizing $`_{eff}`$ below the charm mass. On the other hand, when the calculation is performed with a dynamical charm, the logarithmic divergence (and even the finite coefficient) in eq. (24) is cancelled by the GIM mechanism. For this reason no matching lattice to continuum is needed in this case.
Beside the possible singularities for $`x0`$, further divergences may arise from contact terms of $`Q_i`$ with the external sources, namely for $`xx_\pi `$ or $`xx_K`$. However, it is easy to show that these contact terms do not contribute to the physical amplitudes. Let us consider the Minkowski $`T`$-product
$$\left(๐ฏ_i^\mu \right)_X(q^2,t_\pi ,t_K)=id^4x0|T\left\{\mathrm{\Phi }_\pi (t_\pi ,\stackrel{}{p})J_X^\mu (0)\left[Q_i^u(x)Q_i^c(x)\right]\mathrm{\Phi }_K(t_K,\stackrel{}{k})\right\}|0,$$
(27)
corresponding to the Euclidean Green function of eq. (21). The contact terms are proportional to the following pole terms: $`(p^2m_\pi ^2)^1(k^2m_\pi ^2)^1`$ or $`(p^2m_K^2)^1(k^2m_K^2)^1`$, while the on-shell amplitudes are obtained form the coefficient of $`(p^2m_\pi ^2)^1(k^2m_K^2)^1`$. As we shall discuss in more detail in the next section, these different pole structures in the Minkowski space correspond to a different $`t_\pi \mathrm{}`$ and $`t_K\mathrm{}`$ behavior in the Euclidean case. As a result, we can eliminate the contact terms by an appropriate spectral analysis of the Green function computed in the numerical simulation.
### 3.2 The axial current
With the axial current appearing in the $`T`$-product (17), which is relevant for $`K\pi \nu \overline{\nu }`$ decays, we cannot invoke gauge invariance: it remains true that the quadratic divergence is canceled by GIM, but we must face the problem of the linear one. With power divergences, any subtraction procedure, though non-perturbative, would produce an irreducible (and thus unacceptable) ambiguity in the final result. This implies that the linear divergence can only be an artifact of the regularization procedure. This divergence is indeed absent in regularizations which preserve chirality.
With Wilson fermions the explicit breaking of chiral invariance leads to the appearence of such linear divergence. Since this problem is associated only to the contact term of the integrand (21) for $`x0`$, we can in principle obtain a finite subtracted $`T`$-product, with the correct chiral behaviour of the Green function, by an integration which avoids the region close to $`x=0`$.<sup>3</sup><sup>3</sup>3 We thank Massimo Testa for discussions on this point. Otherwise, one could introduce an appropriate set of counterterms and fix their values by imposing an appropriate set of Ward identities, to recover the correct chiral behaviour. However, both these procedures are technically very complicated to be implemented.
A much simpler and techincally feasible solution is obtained by means of maximally twisted mass terms . In this case, the additional symmetries of the action imply that the amplitude we are interested in is even in the Wilson parameter ($`r`$). This, in turn, implies the absence of the linear divergence which can only be odd in $`r`$, being associated to the breaking of chirality. We have verified this statement by an explicit perturbative calculation at the one loop level. As expected, the structure of the divergent terms is the same as in the continuum and the result is free from ambiguities. The discussion of the axial current can be repeated for a โnon-conservedโ vector current, such as the lattice local electromagnetic current, or the vector component of the weak left-handed current in eq. (17).
At this point we wish to comment about the possibility to determine the physical $`K\pi \pi `$ amplitudes by using information about the following $`T`$-product
$`{\displaystyle d^4xe^{ip_1x}\pi (p_2)|T[_{eff}(0)A^\mu (x)]|K(q_1)},`$ (28)
where the axial current $`A^\mu `$ has the quantum numbers of the pion, but the kinematical configuration does not correspond to an on-shell pion field . The Wick contractions for this $`T`$-product are similar to those considered for $`K\pi \nu \overline{\nu }`$. In particular, the quadratic divergence generated when $`x0`$ is present also in this case. However, the situation is worse than the case discussed in this work, since there is no GIM mechanism to cancel the leading singularity. Of course we can define a renormalized $`T^{}`$-product, but this would entail a finite ambiguity. The practical problems which need to be faced in order to avoid this ambiguity make this calculation very difficult (if not practically impossible) with Wilson-type fermions. For this reason, we do not believe that a lattice study of this $`T`$-product can provide a useful tool to simplify the problem of determining $`K\pi \pi `$ amplitudes.
## 4 $`_{eff}`$ ambiguities
In this section we address the problems arising by the renormalization of the lattice operators of the effective weak Hamiltonian. We first note that only the parity-even or parity-odd terms of the operators contribute to the vector or axial-vector cases, respectively. This observation is relevant since parity-even and parity-odd parts of the operators renormalize in a different way under regularizations which break chiral symmetry. On general grounds, whether chirality is broken or not, the mixing with operators of dimension five or six, in the presence of the GIM mechanism, does not introduce any ambiguity and the corresponding mixing coefficients can be computed in lattice perturbation theory. The problem arises from the mixing of the standard dimension-six operators with the scalar and pseudoscalar densities, which we now consider separately for the two cases.
Schematically, we can write the renormalized operator as
$`\widehat{Q}^\pm `$ $`=`$ $`Z^\pm (\mu a)\left[Q^\pm +C_P(m_cm_u)(m_sm_d)\overline{s}\gamma _5d+C_S(m_cm_u)\overline{s}d\right],`$ (29)
where $`Q^\pm =(Q_1\pm Q_2)/2+\mathrm{}`$ represents the ensemble of all the dimension six and five operators with mixing coefficients computed in perturbation theory. By dimensional arguments, it follows that the coefficients $`C_S`$ and $`C_P`$ are power divergent in the limit $`a0`$:
$$C_P\frac{1}{a},C_S\frac{1}{a^2}.$$
(30)
Using suitable Ward identities (subtraction conditions), we can cancel the divergent parts of $`C_P`$ and $`C_S`$; however, this leaves an ambiguity in their finite values out of the chiral limit . For physical $`K\pi \pi `$ amplitudes this ambiguity turns out to be irrelevant: the pseudoscalar density is proportional to the four divergence of the axial current and its matrix element vanishes for the on-shell $`K\pi \pi `$ transition . We stress that this conclusion does not hold for the $`K\pi `$ case: the off-shell matrix element $`\pi |\overline{s}d|K`$ is different from zero thus, in general, the $`\pi |\widehat{Q}|K`$ matrix element does suffer from this ambiguity.
Since we are interested in physical amplitudes, we must be able to demonstrate that also in the case of radiative decays the matrix elements of scalar and pseudoscalar densities do not contribute to the on-shell amplitudes. This can be done by means of suitable Ward identities and the spectral analysis of the relevant Euclidean Green functions.
In the vector case we can use the following Ward identity
$`{\displaystyle d^4x\left\{\mathrm{\Phi }_\pi (x_\pi )\left[_\mu \widehat{V}_\mu ^{sd}(x)+(m_dm_s)\overline{s}d(x)\right]\widehat{J}_V^\nu (y)\mathrm{\Phi }_K^{}(x_K)\right\}}=`$
$`=\mathrm{\Phi }_K(x_\pi )\widehat{J}_V^\nu (y)\mathrm{\Phi }_K^{}(x_K)+\mathrm{\Phi }_\pi (x_\pi )\widehat{J}_V^\nu (y)\mathrm{\Phi }_\pi ^{}(x_K),`$ (31)
where the term between square bracket is the rotation of the lattice action,
$`\widehat{V}_\mu ^{sd}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\overline{s}(x)U_\mu (x)(r\gamma _\mu )d(x+\mu )\overline{s}(x+\mu )U_\mu ^{}(x)(r+\gamma _\mu )d(x)],`$ (32)
and the last two terms in (4) correspond to the rotation of the pion sink and the kaon source, respectively. The term with the four divergence of $`\widehat{V}_\mu ^{sd}(x)`$, integrated over all space, vanishes. Thus on the left-hand side we are left with the term we are looking for, up to overall factors, namely the contribution of the scalar density to the Euclidean Green function (21), which enters when the bare weak operators are replaced with the renormalized ones. We need to show that this term does not contribute to the physical amplitude.
In the Minkowski space, the physical amplitude is identified by the coefficient of the physical pole for $`p^2m_\pi ^2`$ and $`k^2m_K^2`$. In the Euclidean space, this corresponds to a well-defined dependence on $`t_K`$ and $`t_\pi `$ (for $`t_\pi \mathrm{}`$ and $`t_K\mathrm{}`$), namely
$$\frac{1}{(p^2m_\pi ^2)(k^2m_K^2)}e^{E_K|t_K|}\times e^{E_\pi t_\pi }.$$
(33)
The Ward identity (4) tell us that the contribution of the scalar density give rise to a different pole structure:
$`{\displaystyle \frac{1}{(p^2m_\pi ^2)(k^2m_\pi ^2)}}`$ $``$ $`e^{E_\pi |t_K|}\times e^{E_\pi t_\pi }`$
$`{\displaystyle \frac{1}{(p^2m_K^2)(k^2m_K^2)}}`$ $``$ $`e^{E_K|t_K|}\times e^{E_Kt_\pi }`$ (34)
Thus the scalar density contribution can simply be eliminated by a study of the time dependence of the appropriate Green function. Incidentally, this procedure eliminates also the divergent contact terms mentioned at the end of the section 3.1.
In the axial case, we have a similar situation, up to terms which vanish (linearly or quadratically in the lattice spacing) and inessential numerical factors. In particular, we can use the following Ward identity
$`{\displaystyle d^4x\left\{\mathrm{\Phi }_\pi (x_\pi )\left[_\mu Z_A\widehat{A}_\mu ^{sd}(x)+(m_d+m_s)\overline{s}\gamma _5d(x)+๐ช(a)\right]J_A^\nu (y)\mathrm{\Phi }_K^{}(x_K)\right\}}=`$
$`=\mathrm{\Sigma }_K(x_\pi )J_A^\nu (y)\mathrm{\Phi }_K^{}(x_K)+\mathrm{\Phi }_\pi (x_\pi )J_A^\nu (y)\mathrm{\Sigma }_\pi ^{}(x_K),`$ (35)
where again the term between square bracket is the rotation of the lattice action (for the explicit expressions of $`Z_A`$ and the weak renormalized axial current see ) and $`\mathrm{\Sigma }_i`$ is a scalar particle source. This immediately shows that also the pseudoscalar density give rise to a time dependence different from the one in eq. (33) and thus does not contribute to the on-shell amplitude.
## 5 Conclusions
The potential of rare $`K`$ decays in performing precise tests of the SM and setting stringent bounds on physics beyond the SM depends, to a large extent, from our ability compute their amplitudes within the SM. In this paper we have shown that for a class of very interesting processes, such as $`K^+\pi ^+\nu \overline{\nu }`$ and $`K\pi \mathrm{}^+\mathrm{}^{}`$, the theoretical error associated to non-perturbative effects could be reduced by means of lattice calculations. In particular, the numerical study of the Euclidean Green functions in eq. (21), combined with CHPT, should allow to reach an unprecedented level of precision for these rare decays.
The main problem which needs to be addressed before starting a lattice calculation of these Euclidean Green functions is the absence of power divergences in the extraction of the physical amplitudes. These may originate from contact terms between the weak four-fermion operators and the external fields ($`\pi `$, $`K`$ and the lepton current), or from the mixing of the four fermion operators with operators of lower dimensionality. In this paper we have shown that both these problems can be solved.
As demonstrated in section 4, the spectral analysis necessary to extract the physical amplitudes eliminates both the power divergences due to the operator mixing and the contact terms with the external $`\pi `$ and $`K`$ fields. The only remaining issue is then the ultraviolet behavior associated to the contact terms between the weak operators and the lepton current. This point is different for weak and electromagnetic currents.
In the electromagnetic case, relevant for $`K\pi \mathrm{}^+\mathrm{}^{}`$ decays, gauge invariance prevents the appearance of power divergences for all the popular Wilson-type actions. The cancellation of power divergences is also independent of the GIM mechanism. We can thus match the lattice calculation with the continuum one also in an effective theory where the charm quark is integrated out. The perturbative expressions necessary for this matching at the one-loop level have been presented both for Wilson and Clover fermions. The situation is slightly more complicated for the weak (axial or vector) current, relevant for $`K^+\pi ^+\nu \overline{\nu }`$ decays, where we cannot invoke anymore gauge invariance. One can cancel power divergences also in this case with Wilson-type fermions, but only using maximally twisted mass terms and taking advantage of the GIM mechanism.
In summary, our analysis shows that the numerical study of the Green functions relevant for $`K\pi \mathrm{}^+\mathrm{}^{}`$ decays can be performed with any Wilson-type action, independently of the GIM mechanism. On the other hand, the study of $`K^+\pi ^+\nu \overline{\nu }`$ decays on the lattice requires a more sophisticated action: with Wilson-type fermions the only possibility is to use maximally twisted mass terms. We believe that these results opens a new field of interesting physical applications to the lattice community.
## Acknowledgments
We warmly thank Vittorio Lubicz, Giancarlo Rossi, Silvano Simula and Massimo Testa for illuminating discussions. This work was supported in part by the IHP-RTN program, EC contract No. HPRN-CT-2002-00311 (EURIDICE). |
warning/0506/hep-ph0506111.html | ar5iv | text | # References
Transverse positron polarization in the $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ decay in SM.
V.V.Braguta<sup>1</sup><sup>1</sup>1Electronic address: braguta@mail.ru
Institute for High Energy Physics, Protvino, Russia
## Abstract
In this paper transverse positron polarization in the $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ decay in the framework of SM is considered. It is shown that the final state interaction effect leads to nonzero transverse polarization. Numerical value of the considered effect proved to be negligible. Thus SM contribution to the transverse positron polarization in the $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ decay will not be the obstacle to a new physics searches.
The study of muon decay $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ can give valuable knowledge about lepton sector of weak interactions. The Standard Model successfully describes this decay by vector interaction of four left-handed fermions. To go beyond Standard Model(SM) one introduces into the lagrangian scalar, vector and tensor interactions of right- and left-handed particles . All these interactions can be parameterized by 10 complex constants or one common unphysical phase with nineteen independent real parameters. If some of this parameters are nonzero there exists $`CP`$โviolation in purely leptonic decays.
One possible way to study the phenomenon of $`CP`$โviolation in the muon decay $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ is the measurement of transverse positron polarization that proved to be sensitive to $`CP`$โviolation. Recently the measurement of this observable was improved and new upper bound on the parameters of $`CP`$โviolation has been obtained. The averaged value of transverse positron polarization obtained at this experiment is
$`<P_{T2}>=(3.7\pm 7.7\pm 3.4)\times 10^3`$ (1)
Further improvement of the accuracy can lead to the discovery of $`CP`$โviolation in leptonic decay or can put more strict bounds on parameters of different SM extensions.
In order to study possible SM extensions by the measurement of the transverse positron polarization one needs to find SM contribution to the observable. Transverse muon polarization in $`K^+\pi ^0\mu \nu _\mu `$, $`K^+\mu \nu _\mu \gamma `$, $`K^0\pi ^{}\mu ^+\nu `$ decays, $`T`$-odd correlation in $`K^+\pi ^0\mu \nu _\mu \gamma `$ decay etc. are the examples of similar physical observables that are very sensitive to the effect of $`CP`$โviolation. At tree level these observables are equal zero. In the framework of SM the nonzero contribution is caused by final state interaction effect. Though the effect is strongly suppressed it can be real obstacle in a new physics searches. In this paper SM contribution to the transverse positron polarization due to the final state interaction is considered
There are many ways to write general effective lagrangian that describes the decay $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$. In our paper the following form of the effective lagrangian will be used
$`L_{eff}={\displaystyle \frac{4G_f}{\sqrt{2}}}{\displaystyle \underset{\gamma ,\xi ,\eta ,n,m}{}}g_{\xi \eta }^\gamma \overline{\mu }_\xi \mathrm{\Gamma }^\gamma (\nu _\mu )_n(\overline{\nu _e})_m\mathrm{\Gamma }_\gamma e_\eta ,`$ (2)
where the following designations are used: $`\gamma =\text{scalar(S)},\text{vector(V)},\text{tensor(T)}`$ the type of interactions, $`\xi ,\eta ,n,m`$ are the chiral projections of spinors(left-handed(L), right-handed(R)). It should be noted that the chiral projections of neutrinos $`m,n`$ are uniquely determined if $`\xi ,\eta `$ are given. Usually lagrangian (2) parameterizes all possible SM extensions. We imply here in the frames of SM that only one constant $`g_{LL}^V`$ equals unity and all others are zero. The deviation from this form is caused by radiative corrections. It is shown below that one loop radiative correction leading to nonzero transverse positron polarization can be parameterized by formula (2).
The differential decay $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ probability in the framework of lagrangian (2) is given by the formula
$`{\displaystyle \frac{d^2\mathrm{\Gamma }}{dxdcos\theta }}={\displaystyle \frac{m_\mu }{4\pi ^3}}W_{e\mu }^4G_f^2\sqrt{x^2x_0^2}(F_{IS}(x)+P_\mu cos\theta F_{AS}(x))\left(1+๐ฌ_๐(P_{T1}๐_\mathrm{๐}+P_{T2}๐_\mathrm{๐}+P_L๐_\mathrm{๐})\right),`$ (3)
where $`W_{e\mu }=(m_\mu ^2+m_e^2)/2m_\mu `$, $`x=E_e/W_{e\mu }`$, $`x_0=m_e/W_{e\mu }`$, $`P_\mu `$ is the muon polarization, $`\theta `$ is the angle between muon polarization and direction of positron momentum, $`s_e`$ is the unit vector in the direction of positron spin, $`P_{T1},P_{T2},P_L`$ are the polarizations of the positron corresponding to the unit vectors
$`๐_\mathrm{๐}={\displaystyle \frac{๐ฉ_๐}{|๐ฉ_๐|}},๐_\mathrm{๐}={\displaystyle \frac{๐_\mathrm{๐}\times ๐_\mu }{|๐_\mathrm{๐}\times ๐_\mu |}},๐_\mathrm{๐}=๐_\mathrm{๐}\times ๐_\mathrm{๐}`$ (4)
The functions $`F_{IS},F_{AS},P_{T1},P_{T2},P_L`$ can be expressed through the parameters $`\rho ,\eta ,\xi ,\delta ,etc.`$. In turn this parameters are functions of coupling constants $`g_{\xi \eta }^\gamma `$.
If $`g_{LL}^V=1`$ and all other constants $`g_{\xi \eta }^\gamma `$ are equal zero, the functions $`F_{IS},F_{AS}`$ have the form
$`F_{IS}(x)=x(1x)+{\displaystyle \frac{1}{6}}(4x^23xx_0^2)`$
$`F_{AS}(x)={\displaystyle \frac{1}{3}}\sqrt{x^2x_0^2}\left(1x+{\displaystyle \frac{1}{2}}(4x3+(\sqrt{1x_0^2}1))\right)`$ (5)
In addition to this functions one needs only the expression for transverse positron polarization $`P_{T2}`$ that is sensitive to $`CP`$-violation in muon decay. If we suppose that all coupling constants except $`g_{LL}^V`$ is much less than unity and omit all terms of second order in coupling constants than the expression for $`P_{T2}`$ can be written as follows
$`P_{T2}={\displaystyle \frac{P_\mu sin\theta F_{T2}(x)}{F_{IS}+P_\mu cos\theta F_{AS}(x)}},`$ (6)
where $`F_{T2}`$ is given by
$`F_{T2}={\displaystyle \frac{1}{3}}\sqrt{x^2x_0^2}\left({\displaystyle \frac{\beta ^{}}{8}}\sqrt{1x_0^2}\right)`$ (7)
The constant $`\beta ^{}`$ can be expressed through coupling constants
$`\beta ^{}=4Im(g_{RR}^Vg_{LL}^Sg_{LL}^Vg_{RR}^S)=4Im(g_{RR}^S)`$ (8)
The last formula shows that radiative corrections lead to nonzero transverse polarization only if the constant $`g_{RR}^S`$ acquires nonzero phase. Thus, among all one loop correction diagrams one should consider only those, which contain an imaginary part. Moreover, there is no need to calculate full expression for radiative corrections since only imaginary parts are needed. Calculating imaginary parts of one loop diagrams we use unitarity of $`S`$-matrix in the form
$`ImT_{fi}={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}T_{fn}T_{ni}^{}`$ (9)
In many processes nonzero transverse polarization is caused by electromagnetic final state interaction. But it is not the case for $`\mu ^+e^+\overline{\nu _\mu }\nu _e`$ decay where QED corrections do not lead to nonzero effect. It is easy to prove the statement using formula (9). First it should be noticed that since muon decay is considered the intermediate particles denoted by $`n`$ must be lighter than muon or one gets the amplitude($`T_{ni}`$) for the process in which muon decays into a number of particle with center of mass energy greater than muon mass. Thus positron, some electron positron pairs, photons can be in the intermediate state and of course $`\overline{\nu _\mu }\nu _e`$. Taking into the account that in QED $`\overline{\nu _\mu }\nu _e`$ do not interact with other intermediate particles one gets the amplitude of the process($`T_{fn}`$) where positron, electron positron pairs and photons are in the intermediate state $`n`$ and one positron in the final state. In other words one positron absorbs many particles what are forbidden by the energy conservation law. This proves that in QED $`ImT_{fi}`$ is zero.
If one considers weak interaction in addition to electromagnetic interaction, then there appears diagrams that potentially give nonzero contribution to the transverse positron polarization. These diagrams are presented in Fig. 1 Using (9) one can write imaginary part for the diagram depicted in Fig. 1a.
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}\left({\displaystyle \frac{G_f}{\sqrt{2}}}{\displaystyle ๐\tau _2\overline{\mu }(1+\gamma _5)\gamma _\alpha \widehat{k}_1\gamma _\sigma \nu _\mu \overline{\nu _e}(1+\gamma _5)\gamma ^\sigma \widehat{k}_2\gamma ^\alpha e}\right),`$ (10)
where $`k_1,k_2`$โ 4-momentum of intermediate neutrinos $`\overline{\nu _\mu },\nu _e`$, $`d\tau _2`$โ two particle phase space. The expression (12) can be simplified by using of the formula
$`\gamma ^\mu \gamma ^\nu \gamma ^\lambda =iฯต^{\mu \nu \lambda \rho }\gamma _5\gamma _\rho +g^{\mu \nu }\gamma ^\lambda g^{\mu \lambda }\gamma ^\nu +g^{\lambda \nu }\gamma ^\mu `$ (11)
Then we get
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}\left({\displaystyle \frac{G_f}{\sqrt{2}}}{\displaystyle ๐\tau _2k_1^\rho k_2^\sigma \overline{\mu }(1+\gamma _5)\gamma _\rho \nu _\mu \overline{\nu _e}(1+\gamma _5)\gamma _\sigma e}\right),`$ (12)
The integration of $`k_1^\rho k_2^\sigma `$ over $`d\tau _2`$ results in tensor structures: $`g^{\mu \nu }`$ and $`(p_{\nu _e}^\mu +p_{\nu _\mu }^\mu )(p_{\nu _e}^\nu +p_{\nu _\mu }^\nu )`$. The former tensor structure gives contribution of the form $`\overline{\mu }(1+\gamma _5)\gamma ^\rho \nu _\mu \overline{\nu _e}(1+\gamma _5)\gamma _\rho e`$. It is easy to see that this term gives zero effect since it changes the phase of the coupling constant $`g_{LL}^V`$. The latter tensor structure gives contribution of the form $`\overline{\mu }(1+\gamma _5)\widehat{p}_{\nu _e}\nu _\mu \overline{\nu _e}(1+\gamma _5)\widehat{p}_{\nu _\mu }e`$. Direct calculation of this contribution to the transverse positron polarization gives zero effect. So diagram in fig 1a gives zero transverse polarization.
Letโs consider the diagram in fig. 1b. Imaginary part of the diagram can be written in the form
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}\left(\sqrt{2}G_f{\displaystyle ๐\tau _2\overline{\mu }(1+\gamma _5)\gamma _\alpha \widehat{k}_1\gamma _\sigma \nu _\mu \overline{\nu _e}(1+\gamma _5)\left((\frac{1}{2}+sin^2\theta _W)\gamma ^\alpha \widehat{k}_2\gamma ^\sigma +sin^2\theta _W\gamma ^\alpha \gamma ^\sigma \widehat{p_e}\right)e}\right),`$ (13)
where $`\theta _W`$ is Weinberg angle, $`k_1,k_2`$ are the 4-momentums of intermediate $`\nu _\mu ,e^+`$ respectively, $`p_e`$ is the final positron momentum. Taking into the account (11) one can show that formula (13) is proportional to $`\overline{\mu }(1+\gamma _5)\gamma ^\rho \nu _\mu \overline{\nu _e}(1+\gamma _5)\gamma _\rho e`$ what again gives zero contribution to transverse polarization.
First nonzero contribution comes from the diagram in fig. 1c. After necessary transformations the imaginary part of this diagram can be represented as follows
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}\left(\sqrt{2}G_f{\displaystyle ๐\tau _2\overline{\mu }(1+\gamma _5)\gamma _\alpha \nu _\mu \overline{\nu _e}(1+\gamma _5)\left((12sin^2\theta _W)\widehat{k}_1\gamma ^\alpha \widehat{k}_24m_esin^2\theta _Wk_2^\alpha \right)e}\right),`$ (14)
here $`k_1,k_2`$ are the 4-momentums of intermediate $`e^+,\nu _e`$ respectively. The integration over $`\tau _2`$ can be carried out using the formulae
$`{\displaystyle ๐\tau _2k_2^\alpha }`$ $`=`$ $`4\pi \tau _2^2P^\alpha `$
$`{\displaystyle ๐\tau _2k_1^\alpha k_2^\beta }`$ $`=`$ $`{\displaystyle \frac{2\pi }{3}}(P^2m_e^2)\tau _2^2g^{\alpha \beta }+{\displaystyle \frac{4\pi }{3}}(P^2+2m_e^2)\tau _2^2{\displaystyle \frac{P^\alpha P^\beta }{P^2}},`$ (15)
where $`P=p_e+p_{\nu _e}`$, $`\tau _2=(P^2m_e^2)/8\pi P^2`$ is $`(e^+,\nu _e)`$ phase space. Having made necessary transformations one gets the result
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}8\sqrt{2}\pi G_fm_em_\mu \tau _2^2\left({\displaystyle \frac{1}{3}}(1+2{\displaystyle \frac{m_e^2}{P^2}})(12sin^2\theta _W)+2sin^2\theta _W\right)\overline{\mu }(1\gamma _5)\nu _\mu \overline{\nu }_e(1+\gamma _5)e`$ (16)
So, the diagram in fig. 1c gives nonzero contribution to the $`Img_{RR}^S`$
$`Img_{RR}^S={\displaystyle \frac{\sqrt{2}}{8\pi }}G_fm_em_\mu (1{\displaystyle \frac{m_e^2}{P^2}})^2\left({\displaystyle \frac{1}{3}}(1+2{\displaystyle \frac{m_e^2}{P^2}})(12sin^2\theta _W)+2sin^2\theta _W\right)`$ (17)
Omitting the terms $`m_e^2/P^2`$ which are much less than unity formula (17) becomes much simpler
$`Img_{RR}^S={\displaystyle \frac{1}{12\sqrt{2}\pi }}G_fm_em_\mu (1+4sin^2\theta _W)`$ (18)
Last diagram that can give nonvanishing contribution is presented in fig 1d. Imaginary part of the diagram has the form
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}\left({\displaystyle \frac{G_f}{\sqrt{2}}}{\displaystyle ๐\tau _2Tr(\widehat{k}_2\gamma _\alpha \widehat{k}_1\gamma _\beta )\overline{\mu }(1+\gamma _5)\gamma ^\alpha \nu _\mu \overline{\nu _e}(1+\gamma _5)\gamma ^\beta e}\right)`$ (19)
Transforming expression (19) by the procedure which we used above, one gets:
$`ImT_{fi}={\displaystyle \frac{G_f}{\sqrt{2}}}\left(Img_{RR}^S\overline{\mu }(1\gamma _5)\nu _\mu \overline{\nu }_e(1+\gamma _5)e\right),`$ (20)
where
$`Img_{RR}^S={\displaystyle \frac{1}{6\sqrt{2}\pi }}G_fm_\mu m_e(1+2{\displaystyle \frac{m_e^2}{P^2}})(1m_e^2/P^2)^2={\displaystyle \frac{1}{6\sqrt{2}\pi }}G_fm_\mu m_e`$ (21)
In last equality the terms $`m_e^2/P^2`$ are omitted. Summing the contributions from fig 1c. and fig 1d. we get
$`Img_{RR}^S={\displaystyle \frac{1}{12\sqrt{2}\pi }}G_fm_\mu m_e(3+4sin^2\theta _W)`$ (22)
Numerical estimation of this quantity gives us the result
$`Img_{RR}^S=4\times 10^{11}`$ (23)
Now it is seen that the value obtained in the experiment
$`Img_{RR}^S=(5.2\pm 14.0\pm 2.4)\times 10^3`$ (24)
is much greater than the value of the effect predicted in the framework of SM. This fact allows us to state that the search of the effect of $`CP`$-violation in purely lepton decay by the measurement of transverse positron polarization is not obscured by SM contributions. So the measurement of the transverse positron polarization is very promising in the search of new physics since it either discovers $`CP`$-violation in muon decay or puts very strict bounds on parameters of different SM extensions.
The author thanks professor A.K. Likhoded, V.V. Bytev and A.S. Zhemchugov for useful discussions. This work was partially supported by Russian Foundation of Basic Research under grant 04-02-17530, Russian Education Ministry grant E02-31-96, CRDF grant MO-011-0, Scientific School grant SS-1303.2003.2 and Dynasty foundation. |
warning/0506/cond-mat0506117.html | ar5iv | text | # New perspectives on the Ising model
## I Introduction
It is really very hard to say something new on the Ising model. The model, originally proposed by LenzLenz (1920) in 1920, was exactly solved for the case of an infinite chain by IsingIsing (1925) in 1925. Since then, thousand and thousand of articles and several books have been published on the subject. The reason is that the model is very simple, but still can be considered as the prototype for systems subject to second order phase transitions and can be effectively used for studying critical phenomena. Moreover, the model and its generalizations and modifications can also be used for studying a large variety of physical systems. We do not attempt to summarize the enormous work done in these 80 years; it would go well beyond the purpose of this article. An excellent historical presentation of the Ising model can be found in Ref. Brush, 1967, although it is old and obviously not updated. With no pretension of being exhaustive and complete, we here summarize the principal approaches used in these 80 years.
A basic tool is the transfer matrix methodKramers and Wannier (1941, 1951); Montroll (1941, 1942). By means of this approach, OnsagerOnsager (1944) in 1944 succeed to give an exact solution of the model for a cubic two-dimensional lattice in absence of external magnetic field. The theory of spinors and Lie algebra was used to simplify the Onsager solution Kaufman (1949); Newell and Montroll (1953). Among the exact results for the two-dimensional case, the calculation of the magnetizationYang (1952) and the writing of the spin correlation function in the form of a Toeplitz determinantMontroll et al. (1963) have to be mentioned. Other simplifications of the Onsager solution have been obtained by means of the Jordan-Wigner transformation and fermionization methods Schultz et al. (1964); Lieb et al. (1961); Lieb and Mattis (1962). Different approaches are based on combinatorial methodsKac and Ward (1952); Vdovichenko (1965a, b) and pfaffian methodsHurst and Green (1991); Green and Hurst (1964); Kasteleyn (1963); McCoy and Wu (1973). More recent approaches have seen the Ising Hamiltonian expressed as a Gaussian Grasmannian actionSamuel (1980); Nojima (1998). Along this line, use of operatorial symmetries that simplify the algebra of the transfer matrix has led to the calculation of the partition function for a large class of latticesPlechko (1985, 1988).
Many approximation methods have been used with the goal of obtaining an expression for the partition function valid over a large temperature range: mean field theory, Bethe approximationBethe (1935), cluster variational methodsKikuchi (1951), Monte Carlo simulations, series expansions. The spin correlation functions have been studied at the critical temperatureKadanoff (1969) and in the asymptotic regionWu (1966); Au-Yang (1977). To study critical phenomena and critical indices, tools like series expansions Kirkwood (1938); Domb (1960); Domb and Hunter (1965), scalingPatashinskii and Pokrovskii (1966); Widom (1965); Kadanoff (1966), renormalization group theoryWilson (1971, 1975) have been used.
In spite of the tremendous work done, many problems remain unsolved. The exact partition function in a finite magnetic field is still unknown for dimensions higher than one. Very few exact results have been obtained for the three-dimensional model. There is no exact solution for the two-layer Ising model either. Most of all, a general approach which works in all dimensions and under general boundary conditions, although in some approximation, is needed.
In a recent workMancini (2005), we have shown that there is a large class of models which are exactly solvable in terms of a finite number of parameters that have to be self-consistently calculated. The purpose of the present paper is to apply the method proposed in Ref. Mancini, 2005 to the Ising model and to show that an exact solution of the model does exist for any dimension. In Section 2, we introduce the Ising model for a $`d`$ -dimensional cubic lattice and show that the model is isomorphic to a system of localized spinless interacting particles, satisfying the Fermi statistics. In Section 3, the Hamiltonian of the latter model is solved, that is, a complete finite set of eigenoperators and the relative eigenvalues are determined. Then, as shown in Section 4, the exact form of the retarded Greenโs function (GF) and of the correlation function (CF) can be obtained. In Section 5, we derive a set of equations for determining the charge/spin correlation functions. As the composite operators do not satisfy a canonical algebra, the GF, the CF and the charge/spin correlation functions depend on a set of internal parameters not calculable by the dynamics. For the one-dimensional case, by means of the composite operator methodMancini and Avella (2003); Mancini (2003); Mancini and Avella (2004), we calculate these internal parameters (Section 6) and the charge/spin correlation functions (Section 7). Although obvious, it is worth noticing that all the results reproduce the exact solution known in the literature.
What are the advantages of the present method and what is new in the context of the Ising model? We present a new scheme of calculations for treating the model. The scheme is general and can be applied to any dimension. In the framework of this scheme we show that the model is always solvable for all dimensions. The energy spectra of the system are known. In the one dimensional case we show that the energy scales determined by the spectra rule the behavior at the critical temperature. It is reasonable to expect that this is true also for higher dimensions. General relations among different spin correlation functions have been obtained. These are exact relations and might be used to check the consistency of some approximate treatments or numerical calculations. In order to get quantitative results for the cases of two and three dimensions we have to determine a finite small number of parameters. All the properties of the Ising model, the magnetization, the thermodynamical quantities, the spin correlation functions, depend on these parameters that have to be self-consistently determined. By using algebra and symmetry considerations we calculate these parameters for the case $`d=1`$. Extension of the calculations to higher dimensions is under investigation.
## II The Ising model
The Ising model, in presence of an uniform external magnetic field $`h`$, is described by the following Hamiltonian
$$H_{Ising}=\underset{\mathrm{๐ข๐ฃ}}{}J_{\mathrm{๐ข๐ฃ}}S(๐ข)S(๐ฃ)h\underset{๐ข}{}S(๐ข)$$
(1)
$`S(๐ข)`$ are spin variables, residing on a $`d`$-dimensional Bravais lattice of $`N`$ sites spanned by the vectors $`๐_i=๐ข`$ . The variables $`S(๐ข)`$ takes only two values: up or down, or more simply $`S(๐ข)=\pm 1`$. For a hypercubic lattice of lattice constant $`a`$ with nearest neighbor interactions, the exchange matrix $`J_{\mathrm{๐ข๐ฃ}}`$ is given by
$$\begin{array}{c}J_{\mathrm{๐ข๐ฃ}}=2dJ\alpha _{\mathrm{๐ข๐ฃ}}\alpha _{\mathrm{๐ข๐ฃ}}=\frac{1}{N}\underset{๐ค}{}e^{\text{i}๐ค(๐_i๐_j)}\alpha (๐ค)\hfill \\ \alpha (๐ค)=\frac{1}{d}\underset{n=1}{\overset{d}{}}\mathrm{cos}(๐ค_na)\hfill \end{array}$$
(2)
where $`d`$ is the dimensionality of the system and $`๐ค`$ runs over the vectors in the first Brillouin zone. The exchange constant $`J`$ can be positive or negative, and accordingly the coupling will be ferromagnetic or antiferromagnetic. According to (2), the Hamiltonian (1) can be rewritten as
$$H_{Ising}=dJ\underset{๐ข}{}S(๐ข)S^\alpha (๐ข)h\underset{๐ข}{}S(๐ข)$$
(3)
where
$$S^\alpha (๐ข)=\underset{๐ฃ}{}\alpha _{\mathrm{๐ข๐ฃ}}S(๐ฃ)$$
(4)
It is worth to recall that the Ising Hamiltonian (1) is invariant under the transformation
$$S(๐ข)S(๐ข)hh$$
(5)
Let us consider a system of $`N_e`$ interacting spinless fermions residing on the same lattice and let $`c(i)`$ and $`c^{}(i)`$ be the related annihilation and creation operators. These operators are Heisenberg fields $`[i=(๐ข,t)]`$ satisfying canonical anticommutation relations
$$\begin{array}{c}\{c(๐ข,t),c^{}(๐ฃ,t)\}=\delta _{\mathrm{๐ข๐ฃ}}\hfill \\ \{c(๐ข,t),c(๐ฃ,t)\}=\{c^{}(๐ข,t),c^{}(๐ฃ,t)\}=0\hfill \end{array}$$
(6)
As a consequence of the algebra (6), each site can be occupied at most by a single particle. The occupation number of the site $`๐ข`$, $`\nu (i)=c^{}(i)c(i)`$, takes only the values 0 and 1. By taking into account two-body interactions, the Hamiltonian for such a system reads as
$$H=\underset{๐ข}{}\mu \nu (i)+\frac{1}{2}\underset{\mathrm{๐ข๐ฃ}}{}V(๐ข,๐ฃ)\nu (i)\nu (j)$$
(7)
where $`\mu `$ is the chemical potential and $`V(๐ข,๐ฃ)`$ is the potential. This model Hamiltonian can be connected to the Ising model by defining
$$\nu (i)=\frac{1}{2}[1+S(i)]$$
(8)
It is clear that
$$\begin{array}{c}\nu (i)=0S(i)=1\hfill \\ \nu (i)=1S(i)=+1\hfill \end{array}$$
(9)
By substituting (8) into (7) and by considering only a nearest-neighbor potential we can rewrite the Hamiltonian (7) in the following form
$$H=E_0h\underset{๐ข}{}S(๐ข)dJ\underset{๐ข}{}S(๐ข)S^\alpha (๐ข)$$
(10)
where we defined
$$\begin{array}{c}E_0=(\frac{1}{2}\mu +\frac{1}{4}Vd)N\hfill \\ h=\frac{1}{2}(\mu Vd)\hfill \\ J=\frac{1}{4}Vd\hfill \end{array}$$
(11)
Hamiltonian (10) is just the Ising Hamiltonian (3) as we have the equivalence
$$H_{Ising}=HE_0$$
(12)
The relation between the partition functions is
$$Z_H=e^{\beta E_0}Z_{Ising}$$
(13)
Then, the thermal average of any operator $`A`$ assumes the same value in both models
$$A(\nu )_H=A(S)_{Ising}$$
(14)
According to this, we can choose to study either one or the other model and get both solutions at once. We decide to put attention to the model Hamiltonian (7), which for a nearest-neighbor potential reads as
$$H=\mu \underset{๐ข}{}\nu (i)+Vd\underset{๐ข}{}\nu (i)\nu ^\alpha (i)$$
(15)
where
$$\nu ^\alpha (๐ข,t)=\underset{๐ฃ}{}\alpha _{\mathrm{๐ข๐ฃ}}\nu (๐ฃ,t)$$
(16)
The spin-inversion symmetry (5) of the Ising model (3) corresponds to the particle-hole symmetry exhibited by the Hamiltonian (15). In particular, we have that the chemical potential as a function of $`\nu =\nu (i)`$ scales as
$$\mu (1\nu )=2dV\mu (\nu )$$
(17)
## III Composite operators and equations of motion
It is immediate to see that the charge density operator $`\nu (i)`$ satisfies the equation of motion
$$\text{i}\frac{\nu (i)}{t}=[\nu (i),H]=0$$
(18)
Then, standard methods based on the use of equations of motion and Greenโs function (GF) formalism are not immediately applicable. Indeed, it is easy to check that the causal propagator $`T[\nu (i)\nu (j)]`$ \[$`T`$ is the chronological operator\] and the correlation function $`\nu (i)\nu (j)`$ assume the form
$$T[\nu (i)\nu (j)]=\nu (i)\nu (j)=\frac{1}{N}\underset{๐ค}{}e^{\text{i}๐ค(๐_i๐_j)}\mathrm{\Gamma }(๐ค)$$
(19)
where $`\mathrm{\Gamma }(๐ค)`$ is the zero frequency functionMancini and Avella (2003) which cannot be calculated by means of the dynamics<sup>1</sup><sup>1</sup>1 Use of the formula $`\mathrm{\Gamma }(๐ค)=\frac{1}{2}lim_{\omega 0}\omega G^{(+1)}(๐ค,\omega )`$, where $`G^{(+1)}(๐ค,\omega )`$ is the causal propagator defined in terms of fermionic algebra Mancini and Avella (2003) would lead just to an identity..
Then, in order to solve the Hamiltonian (15) let us consider the composite operator
$$\psi _p(i)=c(i)[\nu ^\alpha (i)]^{p1}\{p=1,2\mathrm{}\mathrm{}\}$$
(20)
This field satisfies the equation of motion
$$\text{i}\frac{}{t}\psi _p(i)=[\psi _p(i),H]=\mu \psi _p(i)+2dV\psi _{p+1}(i)$$
(21)
By taking higher-order time derivatives we generate a hierarchy of composite operators. However, we observe that for $`p1`$ the number operator $`\nu (i)=c^{}(i)c(i)`$ satisfies the following algebra
$$[\nu (i)]^p=[c^{}(i)c(i)]^p=\nu (i)$$
(22)
Therefore, the hierarchy of composite operators (20) must close for a certain value of $`p`$ and we should be able to derive a finite closed set of eigenoperators of the Hamiltonian. To this purpose, on the basis of (22) the following fundamental property of the field $`[\nu ^\alpha (i)]^p`$ can be established
$$[\nu ^\alpha (i)]^p=\underset{m=1}{\overset{2d}{}}A_m^{(p)}[\nu ^\alpha (i)]^m$$
(23)
where the coefficients $`A_m^{(p)}`$ satisfy the relation
$$\underset{m=1}{\overset{2d}{}}A_m^{(p)}=1$$
(24)
The proof of Eq. (23) and the explicit expressions of the coefficients $`A_m^{(p)}`$ are given in Appendix A for the cases $`d=1,2,3`$. We now define the composite operator
$$\psi ^{(d)}(i)=\left(\begin{array}{c}\psi _1(i)\\ \psi _2(i)\\ \mathrm{}\\ \psi _{2d+1}(i)\end{array}\right)=\left(\begin{array}{c}c(i)\\ c(i)\nu ^\alpha (i)\\ \mathrm{}\\ c(i)[\nu ^\alpha (i)]^{2d}\end{array}\right)$$
(25)
After (23), this field is an eigenoperator of the Hamiltonian (15)
$$\text{i}\frac{}{t}\psi ^{(d)}(i)=[\psi ^{(d)}(i),H]=\epsilon ^{(d)}\psi ^{(d)}(i)$$
(26)
where the $`(2d+1)\times (2d+1)`$ matrix $`\epsilon ^{(d)}`$, the energy matrix, is defined in Appendix B. It is easy to see that the eigenvalues $`E_n^{(d)}`$ of the energy matrix are given by
$$E_n^{(d)}=\mu +(n1)Vn=1,2,\mathrm{},2d+1$$
(27)
The Hamiltonian (15) has been solved since we know a complete set of eigenoperators and eigenvalues, and we can proceed to the calculations of observable quantities. This will be done in the next Sections by using the Greenโs function formalism .
## IV Retarded and correlation functions
We define now the thermal retarded Greenโs function
$`G^{(d)}(i,j)`$ $`=`$ $`R[\psi ^{(d)}(i)\psi ^{(d)}{}_{}{}^{}(j)]`$ (28)
$`=`$ $`\theta (t_it_j)\{\psi ^{(d)}(i),\psi ^{(d)}{}_{}{}^{}(j)\}`$
where $`\mathrm{}`$ denotes the quantum-statistical average over the grand canonical ensemble. By introducing the Fourier transform
$`G^{(d)}(i,j)`$ $`={\displaystyle \frac{1}{N}}{\displaystyle \underset{๐ค}{}}{\displaystyle \frac{\text{i}}{(2\pi )}}`$
$`\times {\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}d\omega e^{\text{i}๐ค(๐_i๐_j)\text{i}\omega (t_it_j)}G^{(d)}(๐ค,\omega )`$ (29)
and by means of the Heisenberg equation (26) we obtain the equation
$$[\omega \epsilon ^{(d)}]G^{(d)}(๐ค,\omega )=I^{(d)}(๐ค)$$
(30)
where $`I^{(d)}(๐ค)`$ is the Fourier transform of the normalization matrix, defined as
$`I^{(d)}(๐ข,๐ฃ)`$ $`=`$ $`\{\psi ^{(d)}(๐ข,t),\psi ^{(d)}{}_{}{}^{}(๐ฃ,t)\}`$ (31)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{๐ค}{}}e^{\text{i}๐ค(๐_i๐_j)}I^{(d)}(๐ค)`$
The solution of Eq. (30) is
$$G^{(d)}(๐ค,\omega )=\underset{n=1}{\overset{2d+1}{}}\frac{\sigma ^{(d,n)}(๐ค)}{\omega E_n^{(d)}+\text{i}\delta }$$
(32)
The spectral density matrices $`\sigma _{ab}^{(d,n)}(๐ค)`$ are calculated by means of the formulaMancini and Avella (2003, 2004)
$$\sigma _{ab}^{(d,n)}(๐ค)=\mathrm{\Omega }_{an}^{(d)}\underset{c}{}[\mathrm{\Omega }_{nc}^{(d)}]^1I_{cb}^{(d)}(๐ค)$$
(33)
where $`\mathrm{\Omega }^{(d)}`$is the $`(2d+1)\times (2d+1)`$ matrix whose columns are the eigenvectors of the matrix $`\epsilon ^{(d)}`$. The explicit expressions of $`\mathrm{\Omega }^{(d)}`$are given in Appendix B. The spectral density matrices $`\sigma ^{(d,n)}(๐ค)`$ satisfy the sum rule
$$\underset{n=1}{\overset{2d+1}{}}[E_n^{(d)}]^p\sigma ^{(d,n)}(๐ค)=M^{(d,p)}(๐ค)$$
(34)
where $`M^{(d,p)}(๐ค)`$ are the spectral moments defined as
$$M^{(d,p)}(๐ค)=F.T.\{(i/t)^p\psi ^{(d)}(๐ข,t),\psi ^{(d)}{}_{}{}^{}(๐ฃ,t)\}$$
(35)
$`F.T.`$ stays for the Fourier transform. It is a consequence of the theorem proved in Ref. Mancini, 1998 \[see also pag. 572 in Ref. Mancini and Avella, 2003\] that the spectral density matrices, for $`d=1,2,3`$, satisfy the sum rule (34). The explicit expressions of $`I^{(d)}(๐ค)`$ and $`\sigma ^{(d,n)}(๐ค)`$ are given in Appendices C and D, respectively, for the cases $`d=1,2,3`$. The correlation function
$$C^{(d)}(i,j)=\psi ^{(d)}(i)\psi ^{(d)}{}_{}{}^{}(j)$$
(36)
can be immediately calculated from (32) by using the spectral theorem and one obtains
$`C^{(d)}(i,j)`$
$`={\displaystyle \frac{1}{N}}{\displaystyle \underset{๐ค}{}}{\displaystyle \frac{1}{(2\pi )}}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}๐\omega e^{\text{i}(๐_i๐_j)\text{i}\omega (t_it_j)}C^{(d)}(๐ค,\omega )`$ (37)
$$C^{(d)}(๐ค,\omega )=\pi \underset{n=1}{\overset{2d+1}{}}\delta [\omega E_n^{(d)}]T_n^{(d)}\sigma ^{(d,n)}(๐ค)$$
(38)
with
$$T_n^{(d)}=1+\mathrm{tanh}\left(\frac{E_n^{(d)}}{2k{}_{B}{}^{}T}\right)$$
(39)
Equations (32) and (38) are an exact solution of the model Hamiltonian (15). One is able to obtain an exact solution as the composite operators $`\psi _p(i)=c(i)[\nu ^\alpha (i)]^{p1}`$ constitute a closed set of eigenoperators of the Hamiltonian. However, as stressed in Ref. Mancini and Avella, 2003, the knowledge of the GF is not fully achieved yet. The algebra of the field $`\psi ^{(d)}(i)`$ is not canonical: as a consequence, the normalization matrix $`I^{(d)}(๐ค)`$ in the equation ( 30) contains some unknown static correlation functions, correlators (see Appendix C for explicit calculations), that have to be self-consistently calculated. According to the scheme of calculations proposed by the composite operator method Mancini and Avella (2003); Mancini (2003); Mancini and Avella (2004)(COM), one way of calculating these unknown correlators is by specifying the representation where the GF are realized. The knowledge of the Hamiltonian and of the operatorial algebra is not sufficient to completely determine the GF. The GF refer to a specific representation (i.e., to a specific choice of the Hilbert space) and this information must be supplied to the equations of motion that alone are not sufficient to completely determine the GF. Usually, the use of composite operators leads to an enlargement of the Hilbert space by the inclusion of some unphysical states. Since the GF depend on the unknown correlators, it is clear that the value of these parameters and the representation are intimately related. The procedure is the following. We set up some requirements on the representation and determine the correlators in order that these conditions be satisfied. From the algebra it is possible to derive several relations among the operators. We will call algebra constraints (AC) all possible relations among the operators dictated by the algebra. This set of relations valid at microscopic level must be satisfied also at macroscopic level, when expectations values are considered. Use of these considerations leads to some self-consistent equations which will be used to fix the unknown correlators appearing in the normalization matrix. An immediate set of rules is given by the equation
$$\psi ^{(d)}(i)\psi ^{(d)}(i)=\frac{1}{N}\underset{๐ค}{}\frac{1}{2\pi }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}๐\omega C^{(d)}(๐ค,\omega )$$
(40)
where the l.h.s. is fixed by the AC and the boundary conditions compatible with the phase under investigation, while in the r.h.s. the correlation function $`C(๐ค,\omega )`$ is computed by means of equation of motion \[cfr. Eq. (38)\].
Another important set of AC can be derived by observing that there exist some operators, $`O`$, which project out of the Hamiltonian a reduced part
$$OH=OH_0$$
(41)
When $`H_0`$ and $`H_I=HH_0`$ commute, the quantum statistical average of the operator O over the complete Hamiltonian $`H`$ must coincide with the average over the reduced Hamiltonian $`H_0`$
$$Tr\{Oe^{\beta H}\}=Tr\{Oe^{\beta H_0}\}$$
(42)
Another relation is the requirement of time translational invariance which leads to the condition that the spectral moments, defined by Eq. (35 ), must satisfy the following relation
$$M_{ab}^{(d,p)}(๐ค)=[M_{ba}^{(d,p)}(๐ค)]^{}$$
(43)
It can be shown that if (43) is violated, then states with a negative norm appear in the Hilbert space. Of course the above rules are not exhaustive and more conditions might be needed.
According to the calculations given in appendices C and D, the GF and the correlation functions depend on the following parameters: external parameters $`(\mu ,T,V)`$, internal parameters $`(C_{1,1}^{(d)\alpha },C_{1,2}^{(d)\alpha },\mathrm{}C_{1,2d}^{(d)\alpha })`$, and $`(\kappa ^{(1)},\kappa ^{(1)},\mathrm{}\kappa ^{(2d)})`$, defined as
$$C_{\mu ,\nu }^{(d)\alpha }=\psi _\mu ^{(d)\alpha }(i)\psi _\nu ^{(d)}(i)$$
(44)
$$\kappa ^{(p)}=[v^\alpha (i)]^p$$
(45)
The parameters $`C_{\mu ,\nu }^{(d)\alpha }`$ are determined by means of their own definitions (44), where the r.h.s. is calculated by means of ( 37)-(38). This equation gives
$$C^{(d)\alpha }=\frac{1}{2}\underset{n=1}{\overset{2d+1}{}}T_n^{(d)}\frac{1}{N}\underset{๐ค}{}\alpha (๐ค)\sigma ^{(d,n)}(๐ค)$$
(46)
From the results given in the Appendices C and D, we see that the spectral density matrices have the form
$$\sigma ^{(d,n)}(๐ค)=\mathrm{\Lambda }_0^{(d,n)}+\alpha (๐ค)\mathrm{\Lambda }_1^{(d,n)}$$
(47)
where the matrices $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_1`$ do not depend on momentum $`๐ค`$. Putting (47) into (46) we obtain
$$C^{(d)\alpha }=\frac{1}{4d}\underset{n=1}{\overset{2d+1}{}}T_n^{(d)}\mathrm{\Lambda }_1^{(d,n)}$$
(48)
Calculations given in the Appendices C and D show that the matrices $`\mathrm{\Lambda }_1^{(d,n)}`$ are linear combinations of the matrix elements $`C_{1,p}^{(d)\alpha }`$. Then, Eq. (48) gives a system of homogeneous linear equations. The determinant of this system is only function of the external parameters $`\mu ,T,V`$. This function will vanish only if there is a particular relation among these parameters. Since these parameters are independent variables the only solution is that all the matrix elements must vanish
$$C_{1,p}^{(d)\alpha }=0$$
(49)
The matrices $`\mathrm{\Lambda }_1^{(d,n)}`$ are zero and the correlation function $`C^{(d)}(๐ค,\omega )`$ does not depend on momentum, as we expected. In the coordinate space the CF takes the expression
$$C^{(d)}(i,j)=\delta _{\mathrm{๐ข๐ฃ}}\frac{1}{2}\underset{n=1}{\overset{2d+1}{}}T_n\mathrm{\Lambda }_0^{(d,n)}e^{\text{i}E_n^{(d)}(t_it_j)}$$
(50)
The correlation function depends on $`2d`$ internal parameters: $`\kappa ^{(1)},\mathrm{},\kappa ^{(2d)}`$. In order to determine these parameters, we use the Pauli condition (40) which gives the self-consistent equations
$$\kappa ^{(p)}\lambda ^{(p)}=C_{1,p+1}^{(d)}(p=0,1,\mathrm{}2d)$$
(51)
where $`C_{1,p+1}^{(d)}=\psi _1^{(d)}(i)\psi _p^{(d)}(i)`$ is calculated by means of (50). New correlation functions
$$\lambda ^{(p)}=\nu (i)[\nu ^\alpha (i)]^p$$
(52)
appear and the set of self-consistent equations (51) is not sufficient to determine all unknown parameters. One needs more conditions. In the case of one-dimensional systems these extra conditions can be obtained by using the property (42).
## V Charge correlations functions
In Sections 3 and 4, we have solved the problem of the Ising model in terms of a set of local parameters, defined by (45) and (52). In this Section, we want to show how we can calculate non-local correlation functions. Let us define the causal Greenโs function (for simplicity in this Section we drop the superindex $`(d)`$)
$`F^C(i,l,j)`$ $`=`$ $`T[\psi (i)\psi ^{}(l)]\nu (j)`$ (53)
$`=`$ $`\theta (t_it_l)\psi (i)\psi ^{}(l)\nu (j)`$
$`\theta (t_lt_i)\psi ^{}(l)\psi (i)\nu (j)`$
the retarded and advanced functions
$`F^{R,A}(i,l,j)`$ $`=`$ $`R,A[\psi (i)\psi ^{}(l)]\nu (j)`$ (54)
$`=`$ $`\pm \theta [\pm (t_it_l)]\{\psi (i),\psi ^{}(l)\}\nu (j)`$
the correlation functions
$$\begin{array}{c}D^{\psi \psi ^{}}(i,l,j)=\psi (i)\psi ^{}(l)\nu (j)\hfill \\ D^{\psi ^{}\psi }(i,l,j)=\psi ^{}(l)\psi (i)\nu (j)\hfill \end{array}$$
(55)
where $`\psi (i)`$ is the composite field defined in (25) and we used the fact the field operator $`\nu (j)`$ does not depend on time. The Fourier transforms of these quantities read as
$$\begin{array}{c}F^Q(i,l,j)=\frac{\text{i}}{2\pi }๐\omega e^{\text{i}\omega (t_it_l)}F(๐ข,๐ฅ,๐ฃ;\omega )\hfill \\ D^{\psi \psi ^{}}(i,l,j)=\frac{1}{2\pi }๐\omega e^{\text{i}\omega (t_it_l)}D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ;\omega )\hfill \end{array}$$
(56)
where $`Q=C,R,A.`$ By means of the equation of motion (26) we have
$$(\omega \epsilon )F^Q(๐ข,๐ฅ,๐ฃ;\omega )=J(๐ข,๐ฅ,๐ฃ)$$
(57)
$$\begin{array}{c}(\omega \epsilon )D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ;\omega )=0\hfill \\ (\omega \epsilon )D^{\psi ^{}\psi }(๐ข,๐ฅ,๐ฃ;\omega )=0\hfill \end{array}$$
(58)
where the matrix $`J(๐ข,๐ฅ,๐ฃ)`$ is defined as
$$J(๐ข,๐ฅ,๐ฃ)=\{\psi (๐ข,t),\psi ^{}(๐ฅ,t)\}n(๐ฃ)$$
(59)
The most general solution of Eq. (57) is
$`F^Q(๐ข,๐ฅ,๐ฃ;\omega )`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{2d+1}{}}}[P{\displaystyle \frac{\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)}{\omega E_n}}`$ (60)
$`\text{i}\pi \delta (\omega E_n)g^{Q(n)}(๐ข,๐ฅ,๐ฃ)]`$
where
$$\tau _{ab}^{(n)}(๐ข,๐ฅ,๐ฃ)=\mathrm{\Omega }_{an}\underset{c=1}{\overset{2d+1}{}}\mathrm{\Omega }_{nc}^1J_{cb}(๐ข,๐ฅ,๐ฃ)$$
(61)
while the function $`g^{Q(n)}(๐ข,๐ฅ,๐ฃ)`$ must be determined. $`P`$ denotes the principal value. By recalling the retarded and advanced nature of $`F^{R,A}(i,l,j)`$, it is immediate to see that
$$g^{R(n)}(๐ข,๐ฅ,๐ฃ)=g^{A(n)}(๐ข,๐ฅ,๐ฃ)=\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)$$
(62)
Therefore
$$F^{R,A}(๐ข,๐ฅ,๐ฃ;\omega )=\underset{n=1}{\overset{2d+1}{}}\frac{\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)}{\omega E_n\pm \text{i}\delta }$$
(63)
The solution of (58) is
$$\begin{array}{c}D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ;\omega )=\underset{n=1}{\overset{2d+1}{}}\delta (\omega E_n)d^{\psi \psi ^{}(n)}(๐ข,๐ฅ,๐ฃ)\hfill \\ D^{\psi ^{}\psi }(๐ข,๐ฅ,๐ฃ;\omega )=\underset{n=1}{\overset{2d+1}{}}\delta (\omega E_n)d^{\psi ^{}\psi (n)}(๐ข,๐ฅ,๐ฃ)\hfill \end{array}$$
(64)
where the matrices $`d^{\psi \psi ^{}(n)}(๐ข,๐ฅ,๐ฃ)`$ and $`d^{\psi ^{}\psi (n)}(๐ข,๐ฅ,๐ฃ)`$ have to be determined. From the definitions (53)-(55) we can derive the following exact relations
$$\begin{array}{c}F^R(i,l,j)+F^A(i,l,j)=2F^C(i,l,j)[\psi (i),\psi ^{}(l)]\nu (j)\hfill \\ F^R(i,l,j)F^A(i,l,j)=\{\psi (i),\psi ^{}(l)\}\nu (j)\hfill \end{array}$$
(65)
A relation between the two correlation functions $`D^{\psi \psi ^{}}(i,l,j)`$ and $`D^{\psi ^{}\psi }(i,l,j)`$ can be established by means of trace properties. Indeed, it is straightforward to derive a KMS-like relation
$`\psi ^{}(l)\psi (i)\nu (j)`$ $`=`$ $`\psi (i,t_ii\beta )\psi ^{}(l)\nu (j)`$ (66)
$`+\delta _{\mathrm{๐ฅ๐ฃ}}\psi (i,t_ii\beta )\psi ^{}(l)`$
By recalling the definitions (55), this last equations can be written as
$$D^{\psi ^{}\psi }(๐ข,๐ฅ,๐ฃ;\omega )=e^{\beta \omega }[D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ;\omega )+\delta _{\mathrm{๐ฅ๐ฃ}}C(๐ข,๐ฅ;\omega )]$$
(67)
where $`C(๐ข,๐ฅ;\omega )`$ is the fermionic correlation function \[see Eq. ( 37)-(38)\]. Therefore, the anticommutator $`\{\psi (i),\psi ^{}(l)\}\nu (j)`$ in (65) can be expressed in terms of the correlation functions as
$$\begin{array}{c}\{\psi (i),\psi ^{}(l)\}\nu (j)=\frac{1}{2\pi }๐\omega e^{i\omega (t_it_l)}\hfill \\ \times [(1+e^{\beta \omega })D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ;\omega )+\delta _{\mathrm{๐ฅ๐ฃ}}e^{\beta \omega }C(๐ข,๐ฅ;\omega )]\hfill \end{array}$$
(68)
Analogous expression holds for the commutator. By means of (68) and by recalling that \[see Eqs. (37)-(38)\]
$$\begin{array}{c}(\omega \epsilon )C(๐ข,๐ฅ;\omega )=0\hfill \\ C(๐ข,๐ฅ;\omega )=\underset{n=1}{\overset{2d+1}{}}\delta (\omega E_n)c^{(n)}(๐ข,๐ฅ)\hfill \end{array}$$
(69)
we find that equations (65) have the following form
$$\begin{array}{c}\underset{n=1}{\overset{2d+1}{}}\delta (\omega E_n)\{g^{C(n)}(๐ข,๐ฅ,๐ฃ)\frac{1}{2\pi }[(1e^{\beta \omega })\hfill \\ \times d^{\psi \psi ^{}(n)}(๐ข,๐ฅ,๐ฃ)\delta _{\mathrm{๐ฅ๐ฃ}}e^{\beta \omega }c^{(n)}(๐ข,๐ฅ)]\}=0\hfill \end{array}$$
(70)
$$\begin{array}{c}\underset{n=1}{\overset{2d+1}{}}\delta (\omega E_n)\{\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)\frac{1}{2\pi }[(1+e^{\beta \omega })\hfill \\ \times d^{\psi \psi ^{}(n)}(๐ข,๐ฅ,๐ฃ)+\delta _{\mathrm{๐ฅ๐ฃ}}e^{\beta \omega }c^{(n)}(๐ข,๐ฅ)]\}=0\hfill \end{array}$$
(71)
By recalling that $`E_n^{(d)}=\mu +(n1)V`$, the solution of (70) and (71,) is:
$`d^{\psi \psi (n)}(๐ข,๐ฅ,๐ฃ)`$ $`=`$ $`{\displaystyle \frac{2\pi }{1+e^{\beta E_n}}}\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)`$ (72)
$`{\displaystyle \frac{\delta _{\mathrm{๐ฅ๐ฃ}}e^{\beta E_n}}{1+e^{\beta E_n}}}c^{(n)}(๐ข,๐ฅ)`$
$`g^{C(n)}(๐ข,๐ฅ,๐ฃ)`$ $`=`$ $`{\displaystyle \frac{1e^{\beta E_n}}{1+e^{\beta E_n}}}\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)`$ (73)
$`{\displaystyle \frac{\delta _{\mathrm{๐ฅ๐ฃ}}}{2\pi }}{\displaystyle \frac{2e^{\beta E_n}}{1+e^{\beta E_n}}}c^{(n)}(๐ข,๐ฅ)`$
By putting (72) and (73) into (60) and (64) we have
$`D^{\psi \psi ^{}}(i,l,j)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{2d+1}{}}}{\displaystyle \frac{e^{\text{ i}E_n(t_it_l)}}{1+e^{\beta E_n}}}[\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)`$ (74)
$`{\displaystyle \frac{1}{2\pi }}\delta _{\mathrm{๐ฅ๐ฃ}}e^{\beta E_n}c^{(n)}(๐ข,๐ฅ)]`$
$`D^{\psi ^{}\psi }(i,l,j)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{2d+1}{}}}{\displaystyle \frac{e^{\text{ i}E_n(t_it_l)}e^{\beta E_n}}{1+e^{\beta E_n}}}[\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)`$ (75)
$`+{\displaystyle \frac{1}{2\pi }}\delta _{\mathrm{๐ฅ๐ฃ}}c^{(n)}(๐ข,๐ฅ)]`$
$`F^C(i,l,j)`$ $`=`$ $`{\displaystyle \frac{\text{i}}{2\pi }}{\displaystyle ๐\omega e^{\text{i}\omega (t_it_l)}\underset{n=1}{\overset{2d+1}{}}\frac{\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)}{1+e^{\beta E_n}}}`$ (76)
$`\times \left[{\displaystyle \frac{1}{\omega E_n+\text{i}\delta }}+{\displaystyle \frac{e^{\beta E_n}}{\omega E_n\text{i}\delta }}\right]`$
$`{\displaystyle \frac{\text{i}}{2\pi }}\delta _{๐ฅ,๐ฃ}{\displaystyle ๐\omega e^{i\omega (t_it_l)}\underset{n=1}{\overset{2d+1}{}}\frac{c^{(n)}(๐ข,๐ฅ)e^{\beta E_n}}{1+e^{\beta E_n}}}`$
$`\left[{\displaystyle \frac{1}{\omega E_n+\text{i}\delta }}{\displaystyle \frac{1}{\omega E_n\text{i}\delta }}\right]`$
From the study of the fermionic sector we have
$$c^{(n)}(๐ข,๐ฅ)=\frac{2\pi }{1+e^{\beta E_n}}\delta _{\mathrm{๐ข๐ฅ}}\sigma ^{(n)}$$
(77)
where $`\sigma ^{(n)}`$are the spectral functions given in Appendix D. Then, at equal time (74) becomes
$`D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ)={\displaystyle \underset{n=1}{\overset{2d+1}{}}}{\displaystyle \frac{1}{1+e^{\beta E_n}}}`$
$`\times [\tau ^{(n)}(๐ข,๐ฅ,๐ฃ)\delta _{\mathrm{๐ข๐ฅ}}\delta _{\mathrm{๐ฅ๐ฃ}}{\displaystyle \frac{1}{1+e^{\beta E_n}}}\sigma ^{(n)}]`$ (78)
The system (78) gives a system of linear equations for the quantities $`D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ)`$. Since the inhomogeneous terms in this system are proportional to $`\delta _{\mathrm{๐ข๐ฅ}}`$, it is clear that $`D^{\psi \psi ^{}}(๐ข,๐ฅ,๐ฃ)\delta _{\mathrm{๐ข๐ฅ}}`$. Then, we will take $`๐ข=๐ฅ`$ and we write
$$D(๐ข,๐ฃ)=\psi (๐ข)\psi ^{}(๐ข)n(๐ฃ)$$
(79)
From (78) we have the system of equations
$$D(๐ข,๐ฃ)=\underset{n=1}{\overset{2d+1}{}}\frac{1}{1+e^{\beta E_n}}[\tau ^{(n)}(๐ข,๐ฃ)\delta _{\mathrm{๐ข๐ฃ}}\frac{1}{1+e^{\beta E_n}}\sigma ^{(n)}]$$
(80)
where
$$\begin{array}{c}\tau _{ab}^{(n)}(๐ข,๐ฃ)=\mathrm{\Omega }_{an}\underset{c=1}{\overset{2d+1}{}}\mathrm{\Omega }_{nc}^1J_{cb}(๐ข,๐ฃ)\hfill \\ J(๐ข,๐ฃ)=\psi (๐ข)\psi ^{}(๐ข)\nu (๐ฃ)\hfill \end{array}$$
(81)
From its own definition (79) and by using the recurrence relation ( 23), the matrix $`D(๐ข,๐ฃ)`$ has the following structure.
(i) One dimension
$$D^{(1)}(๐ข,๐ฃ)=\left(\begin{array}{ccc}D_{1,1}& D_{1,2}& D_{1,3}\\ D_{1,2}& D_{1,3}& D_{2,3}\\ D_{1,3}& D_{2,3}& D_{3,3}\end{array}\right)$$
(82)
$$\begin{array}{c}D_{1,p}(๐ข,๐ฃ)=K^{(p1)}(๐ข,๐ฃ)\mathrm{\Lambda }^{(p1)}(๐ข,๐ฃ)p=1,2,3\hfill \\ D_{p,3}(๐ข,๐ฃ)=\underset{m=1}{\overset{2}{}}A_m^{(p+1)}D_{1,m+1}(๐ข,๐ฃ)p=2,3\hfill \end{array}$$
(83)
(ii) Two dimensions
$$D^{(2)}(๐ข,๐ฃ)=\left(\begin{array}{ccccc}D_{1,1}& D_{1,2}& D_{1,3}& D_{1,4}& D_{1,5}\\ D_{1,2}& D_{1,3}& D_{1,4}& D_{1,5}& D_{2,5}\\ D_{1,3}& D_{1,4}& D_{1,5}& D_{2,5}& D_{3,5}\\ D_{1,4}& D_{1,5}& D_{2,5}& D_{3,5}& D_{4,5}\\ D_{1,5}& D_{2,5}& D_{3,5}& D_{4,5}& D_{5,5}\end{array}\right)$$
(84)
$$\begin{array}{c}D_{1,p}(๐ข,๐ฃ)=K^{(p1)}(๐ข,๐ฃ)\mathrm{\Lambda }^{(p1)}(๐ข,๐ฃ)p=1,2,\mathrm{},5\hfill \\ D_{p,5}(๐ข,๐ฃ)=\underset{m=1}{\overset{4}{}}A_m^{(p+3)}D_{1,m+1}(๐ข,๐ฃ)p=2,3,\mathrm{},5\hfill \end{array}$$
(85)
(iii) Three dimensions
$$D^{(3)}(๐ข,๐ฃ)=\left(\begin{array}{ccccc}D_{1,1}& D_{1,2}& \mathrm{}& D_{1,6}& D_{1,7}\\ D_{1,2}& D_{1,3}& \mathrm{}& D_{1,5}& D_{2,7}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ D_{1,6}& D_{1,7}& \mathrm{}& D_{5,7}& D_{6,7}\\ D_{1,7}& D_{2,7}& \mathrm{}& D_{6,7}& D_{7,7}\end{array}\right)$$
(86)
$$\begin{array}{c}D_{1,p}(๐ข,๐ฃ)=K^{(p1)}(๐ข,๐ฃ)\mathrm{\Lambda }^{(p1)}(๐ข,๐ฃ)p=1,2,\mathrm{},7\hfill \\ D_{p,7}(๐ข,๐ฃ)=\underset{m=1}{\overset{6}{}}A_m^{(p+5)}D_{1,m+1}(๐ข,๐ฃ)p=2,3,\mathrm{},7\hfill \end{array}$$
(87)
With the definitions
$$\begin{array}{c}K^{(p)}(๐ข,๐ฃ)=<[\nu ^\alpha (๐ข)]^p\nu (๐ฃ)>\hfill \\ \mathrm{\Lambda }^{(p)}(๐ข,๐ฃ)=<\nu (๐ข)[\nu ^\alpha (๐ข)]^p\nu (๐ฃ)>\hfill \end{array}$$
(88)
Then, we only need to calculate the matrix elements $`D_{1,p}(๐ข,๐ฃ)(p=1,2,\mathrm{}2d+1)`$. The matrix $`J(๐ข,๐ฃ)`$ can be obtained from the normalization matrix $`I(๐ข,๐ฃ)=\{\psi (๐ข,t),\psi ^{}(๐ฃ,t)\}`$, calculated in Appendix C, by means of the following substitution
$$\kappa ^{(p)}K^{(p)}(๐ข,๐ฃ)$$
(89)
Then, the matrices $`\tau ^{(n)}(๐ข,๐ฃ)`$ have the same expressions of the spectral matrices $`\sigma ^{(n)}`$ when the following substitution
$$I_{ab}J_{ab}(๐ข,๐ฃ)$$
(90)
is made. It can be seen that for $`๐ฃ=๐ข`$ and $`๐ฃ=๐ข^\alpha `$ the system (80) is exactly equivalent to the system (51). Then, it is enough to consider the case $`๐ฃ๐ข,๐ข^\alpha `$. In this case, the system (80) becomes
$$D(๐ข,๐ฃ)=\frac{1}{2}\underset{n=1}{\overset{2d+1}{}}T_n\tau ^{(n)}(๐ข,๐ฃ)$$
(91)
with $`T_n`$ given by (39). The system (91) gives a set of exact relations among the correlation functions. We might think to solve this system by induction method, since some of the first correlation functions can be expressed in terms of the basic parameters $`\kappa ^{(p)}`$ and $`\lambda ^{(p)}`$. However, when we do this, we immediately see that the number of equations is not sufficient to determine all the correlation functions and we need more equations. Once again, this can be done for the one dimensional system, as we shall see in the next Sections.
## VI Self-consistent equations for one-dimensional systems
Until now the analysis has been carried on in complete generality for a cubic lattice of $`d`$ dimensions. We now consider one-dimensional systems, and in particular we will study an infinite chain in the homogeneous phase. For simplicity of notation we shall drop the superindex $`(d)`$. By means of ( 207) and (222)-(223) the set of equations (51) gives the linear system
$$\begin{array}{c}T_12+(23T_1+4T_2T_3)\nu \hfill \\ +2(T_12T_2+T_3)\kappa ^{(2)}=0\hfill \\ (2T_2T_32)\nu 2(T_2T_3)\kappa ^{(2)}+2\lambda ^{(1)}=0\hfill \\ (T_2T_3)\nu (2+T_22T_3)\kappa ^{(2)}+2\lambda ^{(2)}=0\hfill \end{array}$$
(92)
where, because of translational invariance, we put
$$\nu =\nu (i)=\kappa ^{(1)}$$
(93)
It is immediate to see that for $`\mu =V`$, the solution of the first equation in (92) for $`T>0`$ is
$$\nu =\frac{1}{2}\text{for}\mu =V$$
(94)
This is in agreement with the particle-hole symmetry enjoyed by the model \[see (17)\]. Recalling (8) and (11), this situation corresponds to the zero magnetization of the Ising model in absence of external magnetic field. Coming back to general value of $`\mu `$, it is clear that Eqs. (92) are not sufficient to specify completely the 4 parameters $`\nu ,\kappa ^{(2)},\lambda ^{(1)},\lambda ^{(2)}`$ and we need another equation. A fourth equation can be easily obtained by means of the algebra. We observe that
$$c^{}(i)\nu (i)=0$$
(95)
This relation leads to
$$c^{}(i)e^{\beta H}=c^{}(i)e^{\beta H_0}$$
(96)
where
$$H_0=H2V\nu (i)\nu ^\alpha (i)$$
(97)
By means of the requirement (42) the correlation function $`C_{1,k}=c(i)c^{}(i)[\nu ^\alpha (i)]^{k1}`$ can be written as
$$\frac{C_{1,k}}{C_{1,1}}=\frac{C_{1,k}^{(0)}}{C_{1,1}^{(0)}}$$
(98)
where
$$C_{1,k}^{(0)}=c(i)c^{}(i)[v^\alpha (i)]^{k1}_0$$
(99)
and $`\mathrm{}_0`$denotes the thermal average with respect to $`H_0`$. Let us define the retarded GF
$`G_{1,k}^{(0)}(tt^{})`$ $`=`$ $`R[c(i,t)c^{}(i,t^{})][v^\alpha (i)]^{k1}_0`$ (100)
$`=`$ $`{\displaystyle \frac{\text{i}}{2\pi }}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}๐\omega e^{\text{i}\omega (tt)}G_{1,k}^{(0)}(\omega )`$
By means of the equation of motion
$$[c(i),H_0]=\mu c(i)$$
(101)
we have
$$G_{1,k}^{(0)}(\omega )=\frac{[\nu ^\alpha (i)]^{k1}_0}{\omega +\mu +\text{i}\delta }$$
(102)
Recalling the relation between retarded and correlation functions, from (102) we obtain
$$C_{1,k}^{(0)}=\frac{[\nu ^\alpha (i)]^{k1}_0}{1+e^{\beta \mu }}$$
(103)
By putting this result into (98) we have
$$C_{1,k}=C_{1,1}[\nu ^\alpha (i)]^{k1}_0$$
(104)
By noting that $`[\nu ^\alpha (i)]^2`$can be expressed as \[cfr. (151 )\]
$$[\nu ^\alpha (i)]^2=\frac{1}{2}[\nu ^\alpha (i)+\nu (i_1)\nu (i_2)]$$
(105)
we obtain from (104) the relations
$`C_{1,2}`$ $`=`$ $`C_{1,1}[\nu ^\alpha (i)]_0`$ (106)
$`C_{1,3}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[C_{1,2}+C_{1,1}\nu (i_1)\nu (i_2)_0]`$ (107)
Now, we observeFedro (1976) that $`H_0`$ describes a system where the original lattice is divided in two disconnected sublattices (the chains to the left and to the right of the site $`๐ข`$). Then, in $`H_0`$ -representation, the correlation function which relates sites belonging to different sublattices can be decoupled:
$$a(j)b(m)_0=a(j)_0b(m)_0$$
(108)
for $`j`$ and $`m`$ belonging to different sublattices. By using this property, invariance of $`H_0`$ under axis reflection and (106) we can write
$`\nu (i_1)\nu (i_2)_0`$ $`=`$ $`\nu (i_1)_0\nu (i_2)_0`$ (109)
$`=`$ $`[\nu ^\alpha (i)_0]^2=\left[{\displaystyle \frac{C_{1,2}}{C_{1,1}}}\right]^2`$
By putting (109) into (107), we obtain the following self-consistent equation among the correlation functions
$$C_{1,3}=\frac{1}{2}C_{1,2}\left(1+\frac{C_{1,2}}{C_{1,1}}\right)$$
(110)
By means of (37) and (38) and the results given in Appendices C and D, Eq. (110) takes the expression
$$\begin{array}{c}(4T_2^23T_1T_3)\nu ^2+[T_1T_38\kappa ^{(2)}(T_2^2T_1T_3)]\nu \hfill \\ +2\kappa ^{(2)}[2\kappa ^{(2)}(T_2^2T_1T_3)T_1T_3]=0\hfill \end{array}$$
(111)
This equation together with equations (92) gives a system of 4 self-consistent equations for the 4 parameters $`\nu ,\kappa ^{(2)},\lambda ^{(1)},\lambda ^{(2)}`$ as functions of $`\mu ,T,V`$. By solving the set of linear equations (92) with respect to $`\kappa ^{(2)},\lambda ^{(1)},\lambda ^{(2)}`$ as functions of $`\nu `$, we have
$`\kappa ^{(2)}`$ $`=`$ $`{\displaystyle \frac{2T_1+\nu (2+3T_14T_2+T_3)}{2(T_12T_2+T_3)}}`$ (112)
$`\lambda ^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2(T_12T_2+T_3)}}\{(2T_1)(T_2T_3)`$
$`+\nu [T_1(2+T_22T_3)+T_2(T_36)+4T_3]\}`$
$`\lambda ^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{4(T_12T_2+T_3)}}\{(2T_1)(2+T_22T_3)`$
$`+\nu [410T_2+T_1(6+T_24T_3)+6T_3+3T_2T_3]\}`$
To calculate the parameter $`\nu `$ let us put (112) into (111) and solve with respect to $`\nu `$. We have two roots. One solution corresponds to an unstable state with negative compressibility and must be disregarded. By picking up the right root, and by using the relation
$$T_3=\frac{2T_2^2(2T_1)}{T_1(2T_2)^2+T_2^2(2T_1)}$$
(115)
we find
$$\nu =\frac{1}{2}\left[1+(1T_2)\sqrt{\frac{T_1}{T_12T_1T_2+2T_2^2}}\right]$$
(116)
As shown in Appendix E, the solutions (112)-(116) exactly correspond to the well-known solution of the 1D Ising model, obtained by means of the transfer matrix method. We could manipulate the expression (116) and the ones for $`\kappa ^{(2)},\lambda ^{(1)}`$ and $`\lambda ^{(2)}`$ , obtained by substituting (116) into (112)-(VI), in order to reproduce the expressions of the Ising model, given in Appendix E. However, we prefer to maintain the present expressions as the following discussion will be more transparent. In Section 3, we have seen that in the present model, in the one-dimensional case, there are three energy scales
$$\begin{array}{c}E_1=\mu \hfill \\ E_2=\mu +V\hfill \\ E_3=\mu +2V\hfill \end{array}$$
(117)
At zero temperature, the three functions $`T_1`$, $`T_2`$ and $`T_3`$ are not analytical functions at the points $`\mu =0,`$ $`\mu =V,`$ and $`\mu =2V`$, respectively, and we expect that the parameters $`\nu ,\kappa ^{(2)},\lambda ^{(1)},\lambda ^{(2)}`$ exhibit some discontinuous behavior at these points. As shown in Fig. 1, in the limit $`T0`$ the particle density $`\nu `$ has a discontinuity at $`\mu =V`$ for the case of negative $`V`$ (i.e. $`J>0`$, ferromagnetic coupling) and two discontinuities at $`\mu =0`$ and $`\mu =2V`$ for the case of positive $`V`$ (i.e. $`J<0`$, antiferromagnetic coupling). Here and in the following, we take $`\left|V\right|=1`$: all energies are measured in units of $`\left|V\right|`$. In particular, the particle density increases by increasing $`\mu `$ from zero to one. At zero temperature, in the ferromagnetic case $`\nu `$ is zero for $`\mu <\left|V\right|`$ and equal to one for $`\mu >\left|V\right|`$; in the antiferromagnetic case $`\nu `$ is zero for $`\mu <0`$, jumps to $`1/2`$ and exhibits a plateau, centered at $`\mu =V`$, in the region $`0<\mu <2V`$, jumps to the value $`1`$ for $`\mu >2V`$. The parameter $`\kappa ^{(2)}`$ has a behavior similar to $`\nu `$.
In Fig. 2 we give the parameter $`\lambda ^{(1)}`$ as a function of $`\mu `$. For the ferromagnetic case the behavior is similar to that of $`\nu `$ . Instead, in the antiferromagnetic case $`\lambda ^{(1)}`$, at $`T=0`$, exhibits only one discontinuity point at $`\mu =2V`$, where jumps from zero to one. The parameter $`\lambda ^{(2)}`$ has a behavior similar to $`\lambda ^{(1)}`$. The different behavior exhibited by the pairs $`(\nu ,\kappa ^{(2)})`$ and $`(\lambda ^{(1)},\lambda ^{(2)})`$ for $`V<0`$ is naturally due to the antiferromagnetic correlations, when we recall that the pair $`(\lambda ^{(1)},\lambda ^{(2)})`$ describes a correlation between two first neighboring sites, while $`\kappa ^{(2)}`$ describes correlations between two second neighboring sites. Of course in the point of discontinuity the two limits $`T0`$ and $`\mu \mu _c`$ are not interchangeable. As we shall see in the next Section, the 4 local parameters $`\nu ,\kappa ^{(2)},\lambda ^{(1)},\lambda ^{(2)}`$ are really basic since all the properties of the model are described in terms of them. It is worthwhile to note that some simple relations can be established among the parameters
$$2\kappa ^{(2)}\nu \nu ^2=\frac{(\nu ^2\lambda ^{(1)})^2}{\nu (1\nu )}$$
(118)
$$2(\kappa ^{(2)}\lambda ^{(2)})=(\nu \lambda ^{(1)})+\frac{(\nu \lambda ^{(1)})^2}{1\nu }$$
(119)
The Ising model in one dimension can be described in terms of only two parameters: $`\nu `$ and $`\lambda ^{(1)}`$.
## VII Charge correlation functions for one-dimensional systems
The system of equations (91) establishes some relations among the non-local charge correlation functions $`K^{(p)}(๐ข,๐ฃ)=<[v^\alpha (๐ข)]^pv(๐ฃ)>`$ and $`\mathrm{\Lambda }^{(p)}(๐ข,๐ฃ)=<v(๐ข)[v^\alpha (๐ข)]^pv(๐ฃ)>`$. As already discussed, the number of equations is not sufficient to determine completely the charge correlation functions, and one needs more equations to close the system. In the one-dimensional case a fourth equation can be easily obtained by algebraic considerations. Recalling that $`c^{}(i)e^{\beta H}=c^{}(i)e^{\beta H_0}`$ we can easily derive the following result
$$\mathrm{\Lambda }^{(p)}(i,j)=K^{(p)}(i,j)C_{1,1}[v^\alpha (i)]^pv(j)_0$$
(120)
Now, for $`j>i+a`$ (because of invariance under axis reflection we could choose $`j<ia`$ as well), using the property (108)
$`\nu ^\alpha (i)\nu (j)_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\nu ^\alpha (i)_0`$ (121)
$`\times \nu (j)_0+{\displaystyle \frac{1}{2}}\nu (i+a)\nu (j)_0`$
$$\nu (i+a)\nu (ia)\nu (j)_0=\nu ^\alpha (i)_0\nu (i+a)\nu (j)_0$$
(122)
Therefore, from (120)
$$\mathrm{\Lambda }^{(0)}(i,j)=K^{(0)}(i,j)C_{1,1}\nu (j)_0$$
(123)
$`\mathrm{\Lambda }^{(1)}(i,j)`$ $`=`$ $`K^{(1)}(i,j){\displaystyle \frac{1}{2}}C_{1,2}\nu (j)_0`$ (124)
$`{\displaystyle \frac{1}{2}}C_{1,1}\nu (i+a)\nu (j)_0`$
$`\mathrm{\Lambda }^{(2)}(i,j)`$ $`=`$ $`K^{(2)}(i,j){\displaystyle \frac{1}{2}}[K^{(1)}(i,j)\mathrm{\Lambda }^{(1)}(i,j)]`$ (125)
$`{\displaystyle \frac{1}{2}}C_{1,2}\nu (i+a)\nu (j)_0`$
where we used Eq. (106). Equations (123) and (124) give
$$\nu (j)_0=\frac{1}{C_{1,1}}[K^{(0)}(i,j)\mathrm{\Lambda }^{(0)}(i,j)]$$
(126)
$`\nu (i+a)\nu (j)_0`$ $`={\displaystyle \frac{2}{C_{1,1}}}[K^{(1)}(i,j)\mathrm{\Lambda }^{(1)}(i,j)]`$
$`{\displaystyle \frac{C_{12}}{C_{1,1}^2}}[K^{(0)}(i,j)\mathrm{\Lambda }^{(0)}(i,j)]`$ (127)
Putting (127) into (125) we get the fourth self-consistent equation
$`\mathrm{\Lambda }^{(2)}(i,j)`$ $`=`$ $`K^{(2)}(i,j)`$ (128)
$`[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{C_{1,2}}{C_{1,1}}}][K^{(1)}(i,j)\mathrm{\Lambda }^{(1)}(i,j)]`$
$`+{\displaystyle \frac{C_{1,2}^2}{2C_{1,1}^2}}[K^{(0)}(i,j)\mathrm{\Lambda }^{(0)}(i,j)]`$
It is straightforward to verify that (128) is identically satisfied for $`j=i`$, while for $`j=i^\alpha `$ coincides with equation (119).
By adding (128) to (91), we have the following system of equations for the non-local spin correlation functions $`K^{(p)}(i,j)`$ and $`\mathrm{\Lambda }^{(p)}(i,j)`$
$$2\mathrm{\Lambda }^{(0)}(i,j)a_0K^{(1)}(i,j)+a_1K^{(2)}(i,j)=a_5\nu $$
(129)
$$2\mathrm{\Lambda }^{(1)}(i,j)+a_2K^{(1)}(i,j)2a_3K^{(2)}(i,j)=0$$
(130)
$$2\mathrm{\Lambda }^{(2)}(i,j)+a_3K^{(1)}(i,j)a_4K^{(2)}(i,j)=0$$
(131)
$$\begin{array}{c}\mathrm{\Lambda }^{(2)}(i,j)+p_0[K^{(1)}(i,j)\mathrm{\Lambda }^{(1)}(i,j)]\hfill \\ +p_1\mathrm{\Lambda }^{(0)}(i,j)K^{(2)}(i,j)=p_1\nu \hfill \end{array}$$
(132)
where we put
$$\begin{array}{c}a_0=(3T_14T_2+T_3)\hfill \\ a_1=2(T_12T_2+T_3)\hfill \\ a_2=(2T_2T_32)\hfill \end{array}\begin{array}{c}a_3=(T_2T_3)\hfill \\ a_4=(2+T_22T_3)\hfill \\ a_5=(2T_1)\hfill \end{array}$$
(133)
$$\begin{array}{c}p_0=(\frac{1}{2}+\frac{C_{1,2}}{C_{1,1}})\hfill \\ p_1=\frac{C_{1,2}^2}{2C_{1,1}^2}\hfill \end{array}$$
(134)
and we used $`K^{(0)}(i,j)=\nu (j)=\nu `$. To calculate for general $`\left|ij\right|`$ we proceed by induction.
Let us put $`ij=ma`$ and concentrate the attention on the spin CF $`\mathrm{\Lambda }^{(0)}(m)=\nu (m)\nu (0)`$. We start by observing that
$$\mathrm{\Lambda }^{(0)}(0)=\nu \mathrm{\Lambda }^{(0)}(1)=\lambda ^{(1)}$$
(135)
where the two parameters $`\nu ,`$ $`\lambda ^{(1)}`$ have been calculated in the fermionic sector. By taking $`m=2`$ we can calculate from the system (129)-(132) that
$$\mathrm{\Lambda }^{(0)}(2)=2\kappa ^{(2)}\nu $$
(136)
The results (135)-(136) and use of the relation (118) show that $`\mathrm{\Lambda }^{(0)}(m)`$ for $`m=0,1,2`$ can be cast in the form
$$\mathrm{\Lambda }^{(0)}(m)=\nu ^2+\nu (1\nu )p^m$$
(137)
where the parameter $`p`$ is defined as
$$p=\frac{2\kappa ^{(2)}\nu \nu ^2}{\lambda ^{(1)}\nu ^2}=\frac{\lambda ^{(1)}\nu ^2}{\nu \nu ^2}$$
(138)
By using the expressions of the basic parameters given in Section 5, it is possible to check that $`\left|p\right|<1`$. Then, we can introduce the Fourier transform and reexpress (137) as
$$\mathrm{\Lambda }^{(0)}(m)=\nu ^2+A\nu (1\nu )\frac{a}{2\pi }\underset{\pi /a}{\overset{\pi /a}{}}๐k\frac{e^{ikma}}{1+B\mathrm{cos}(ka)}$$
(139)
where
$$A=\frac{\nu \kappa ^{(2)}}{\kappa ^{(2)}\nu ^2}B=\frac{\lambda ^{(1)}\nu ^2}{\kappa ^{(2)}\nu ^2}$$
(140)
Then, the CF $`\mathrm{\Lambda }^{(0)}(m+1)`$ can be calculated as
$`\mathrm{\Lambda }^{(0)}(m+1)`$ $`=`$ $`\nu {}_{}{}^{2}+A\nu (1\nu ){\displaystyle \frac{a}{2\pi }}{\displaystyle \underset{\pi /a}{\overset{\pi /a}{}}}dk{\displaystyle \frac{e^{ika(m+1)}}{1+B\mathrm{cos}(ka)}}`$ (141)
$`=`$ $`\nu {}_{}{}^{2}+\nu (1\nu )p^{m+1}`$
Therefore, Eq. (137) is valid for any $`m`$. Recalling the definition of $`\mathrm{\Lambda }^{(0)}(m)`$, we can rewrite (137) under the form
$$\frac{<\nu (m)\nu (0>\nu ^2}{\nu \nu ^2}=p^m$$
(142)
Also, from (139) we see that the zero frequency function \[cfr. (19)\] has the expression
$$\mathrm{\Gamma }(k)=\nu ^2(2\pi /a)\delta (k)+\frac{A\nu (1\nu )}{1+B\mathrm{cos}(ka)}$$
(143)
By putting the obtained expression of $`\mathrm{\Lambda }^{(0)}(m)`$ in Eqs. (129 )-(132), we can solve the system. The solution gives
$$K^{(1)}(m)=\nu ^2+\frac{1}{2}\nu (1\nu )(p^{m1}+p^{m+1})$$
(144)
$`K^{(2)}(m)`$ $`=`$ $`\kappa ^{(2)}\nu +\nu (1\nu )`$ (145)
$`\left[{\displaystyle \frac{a_0}{2a_1}}(p^{m1}+p^{m+1}){\displaystyle \frac{2}{a_1}}p^m\right]`$
$`\mathrm{\Lambda }^{(1)}(m)`$ $`=\lambda ^{(1)}\nu +\nu (1\nu )\times `$
$`\times \left[{\displaystyle \frac{2a_0a_3a_1a_2}{4a_1}}(p^{m1}+p^{m+1}){\displaystyle \frac{2a_3}{a_1}}p^m\right]`$ (146)
$`\mathrm{\Lambda }^{(2)}(m)`$ $`=\lambda ^{(2)}\nu +\nu (1\nu )\times `$
$`\times \left[{\displaystyle \frac{a_0a_4a_1a_3}{4a_1}}(p^{m1}+p^{m+1}){\displaystyle \frac{a_4}{a_1}}p^m\right]`$ (147)
In Fig. 3 we give $`p`$ as a function of $`\mu `$ for $`V=1`$ and $`V=1`$ at various temperatures. We see that for negative $`V`$ (ferromagnetic case), $`p`$ is positive and various between zero and 1. For positive $`V`$ (antiferromagnetic case), $`p`$ is negative and various between $`1`$ and zero. In particular, for negative $`V`$, $`p`$ tends to 1 at $`\mu =V`$ in the limit $`T0`$. Instead, for positive $`V`$, $`p`$ tends to $`1`$ at $`\mu =V`$ in the limit $`T0`$. This is seen in Fig. 4 where $`p`$ is plotted versus $`T`$ at $`\mu =V=1`$ and at $`\mu =V=1`$.
Let us now discuss the correlation functions. $`\mathrm{\Lambda }^{(0)}(m)`$ is plotted against $`m`$ for $`\mu =V=1`$ \[Fig. 5 (top)\] and for $`\mu =V=1`$ \[Fig. 5 (bottom)\] at various temperatures. We see that when zero temperature is approached a long-range order of ferromagnetic and antiferromagnetic type is established, respectively. Also, we can see from ( 141)-(147) that for $`m\mathrm{}`$ and $`T0`$ (i.e. not at the critical temperature) the spin correlation functions assume the ergodic value. At the critical temperature $`T=0`$ we have breakdown of the ergodicity.
## VIII Conclusions
The Ising model in presence of an external magnetic field is isomorphic to a model of localized spinless interacting particles, satisfying the Fermi statistics. The latter model belongs to a class of models always solvable, as shown in Ref. Mancini, 2005. On this basis, we have constructed a general solution of the Ising model which holds for any dimensionality of the system. The Hamiltonian of the model has been solved in terms of a complete finite set of eigenoperators and eigenvalues. The Greenโs function and the correlation functions of the fermionic model are exactly known and are expressed in terms of a finite small number of parameters that have to be self-consistently determined. By using the equation of the motion method, we have derived a set of equations which connect different spin correlation functions. The scheme that emerges is that it is possible to describe the Ising model from a unified point of view where all the properties are connected to a small number of local parameters, and where the critical behavior is controlled by the energy scales fixed by the eigenvalues of the Hamiltonian. The latter considerations have been proved from a quantitative point of view in the one-dimensional case, where the equations which determine the self-consistent parameters and the spin correlation function have been solved. For $`d=1`$ all the properties of the system have been calculated and obviously agree with the exact results reported in the literature. Extension of the calculations to higher dimensions is under investigation.
After the paper was completed and submitted for publication, the author learned that approaches to the spin-1/2 Ising model based on a fermionization of the model have been previously reported in Refs. Tyablikov and Fedyanin, 1967 and Kalashnikov and Fradkin, 1969a. The author wishes to thank the referee and Prof. L. De Cesare for putting these papers to his attention. In particular, in Ref. Tyablikov and Fedyanin, 1967, Tyablikov and Fedyanin showed that the chain of equations for the double-time GF closes and the number of equations is determined only by the co-ordination number, independently by the dimensionality. This conclusion agrees with the results given in Sections 3 and 4. In order to close the set of equations for the fermionic correlation functions in the one-dimensional case, the authors of Ref. Tyablikov and Fedyanin, 1967 assumed ergodicity and solved the system, obtaining the exact solution of the 1D Ising model for an infinite chain. It should be remarked that ergodicity breaks down for finite systems and at the critical points. Kalashnikov and Fradkin in Ref. Kalashnikov and Fradkin, 1969a used the spectral density method Kalashnikov and Fradkin (1969b) to derive a system of equations for the correlation functions; also in this case the approach is valid for any dimension. However, the number of equations is less than the number of correlation functions.
## Appendix A Algebraic relations
As mentioned in Section 3, the number density operator $`\nu (i)=c^{}(i)c(i)`$ satisfies the algebra
$$\nu ^p(i)=\nu (i)p1$$
(148)
From this algebra an important relation can be derived. for the operator
$$\nu ^\alpha (i)=\underset{๐ฃ}{}\alpha _{\mathrm{๐ข๐ฃ}}\nu (j)=\frac{1}{2d}\underset{m=1}{\overset{2d}{}}\nu (i_m)$$
(149)
where $`i_m`$ are the first neighbors of the site $`๐ข`$. We shall discuss separately the cases of different dimensions.
### A.1 One dimension
We start from the equation
$$[\nu ^\alpha (i)]^p=\frac{1}{2^p}\underset{m=0}{\overset{p}{}}\left(\begin{array}{c}p\\ m\end{array}\right)\nu (i_1)^{pm}\nu (i_2)^m$$
(150)
After subtracting the terms $`m=0`$ and $`m=p`$, we can use the algebraic relation (148) to obtain
$$[\nu ^\alpha (i)]^p=\frac{1}{2^p}[2\nu ^\alpha (i)+a_p\nu (i_1)\nu (i_2)]$$
(151)
with
$$a_p=\underset{m=1}{\overset{p1}{}}\left(\begin{array}{c}p\\ m\end{array}\right)=2^p2$$
(152)
From (151), by putting $`p=2`$ we obtain
$$\nu (i_1)\nu (i_2)=2[\nu ^\alpha (i)]^2\nu ^\alpha (i)$$
(153)
By substituting (153) into (151) we have the recurrence rule
$$[\nu ^\alpha (i)]^p=\underset{m=1}{\overset{2}{}}A_m^{(p)}[\nu ^\alpha (i)]^m$$
(154)
where
$`A_1^{(p)}`$ $`=`$ $`{\displaystyle \frac{1}{2^p}}(2a_p)=2^{2p}1`$ (155)
$`A_2^{(p)}`$ $`=`$ $`{\displaystyle \frac{1}{2^p}}2a_p=2(12^{1p})`$ (156)
We note that the coefficients $`A_m^{(p)}`$ satisfy the relation
$$\underset{m=1}{\overset{2}{}}A_m^{(p)}=1$$
(157)
In table 1 we give the values of the $`A_m^{(p)}`$โs for $`1p6`$.
| p | $`A_1^{(p)}`$ | $`A_2^{(p)}`$ |
| --- | --- | --- |
| 1 | 1 | 0 |
| 2 | 0 | 1 |
| 3 | $`\frac{1}{2}`$ | $`\frac{3}{2}`$ |
| 4 | $`\frac{3}{4}`$ | $`\frac{7}{4}`$ |
| 5 | $`\frac{7}{8}`$ | $`\frac{15}{8}`$ |
| 6 | $`\frac{15}{16}`$ | $`\frac{31}{16}`$ |
### A.2 Two dimensions
We start from the equation
$`[\nu ^\alpha (i)]^p`$ $`={\displaystyle \frac{1}{4^p}}{\displaystyle \underset{m=0}{\overset{p}{}}}\left(\begin{array}{c}p\\ m\end{array}\right)\nu (i_1)^{pm}{\displaystyle \underset{n=0}{\overset{m}{}}}\left(\begin{array}{c}m\\ n\end{array}\right)`$ (162)
$`\times \nu (i_2)^{mn}{\displaystyle \underset{l=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ l\end{array}\right)\nu (i_3)^{nl}\nu (i_4)^l`$ (165)
By proceeding as in the case of one dimension, use of the algebraic relation (148) leads to
$$[\nu ^\alpha (i)]^p=\frac{1}{4^p}\underset{m=1}{\overset{4}{}}b_m^{(p)}Z_m$$
(166)
where the operators $`Z_m`$ are defined as
$$\begin{array}{c}Z_1=4\nu ^\alpha (i)\hfill \\ Z_2=\nu (i_1)\nu (i_2)+\nu (i_1)\nu (i_3)+\nu (i_1)\nu (i_4)\hfill \\ +\nu (i_2)\nu (i_3)+\nu (i_2)\nu (i_4)+\nu (i_3)\nu (i_4)\hfill \\ Z_3=\nu (i_1)\nu (i_2)\nu (i_3)+\nu (i_1)\nu (i_2)\nu (i_4)\hfill \\ +\nu (i_1)\nu (i_3)\nu (i_4)+\nu (i_2)\nu (i_3)\nu (i_4)\hfill \\ Z_4=\nu (i_1)\nu (i_2)\nu (i_3)\nu (i_4)\hfill \end{array}$$
(167)
and the coefficients $`b_m^{(p)}`$ have the expressions
$$\begin{array}{c}b_1^{(p)}=4\hfill \\ b_2^{(p)}=\underset{m=1}{\overset{p1}{}}\left(\begin{array}{c}p\hfill \\ m\hfill \end{array}\right)=2^p2\hfill \\ b_3^{(p)}=\underset{m=2}{\overset{p1}{}}\left(\begin{array}{c}p\hfill \\ m\hfill \end{array}\right)\underset{n=1}{\overset{m1}{}}\left(\begin{array}{c}m\hfill \\ n\hfill \end{array}\right)=3(12^p+3^{p1})\hfill \\ b_4^{(p)}=\underset{m=3}{\overset{p1}{}}\left(\begin{array}{c}p\hfill \\ m\hfill \end{array}\right)\underset{n=2}{\overset{m1}{}}\left(\begin{array}{c}m\hfill \\ n\hfill \end{array}\right)\underset{l=1}{\overset{n1}{}}\left(\begin{array}{c}n\hfill \\ n\hfill \end{array}\right)\hfill \\ =4(1+32^{p1}3^p+4^{p1})\hfill \end{array}$$
(168)
By solving the system (166) with respect to variables $`Z_m`$, we can obtain the recurrence rule
$$[\nu ^\alpha (i)]^p=\underset{m=1}{\overset{4}{}}A_m^{(p)}[\nu ^\alpha (i)]^m$$
(169)
where the coefficients $`A_m^{(p)}`$are defined as
$$\begin{array}{c}A_1^{(p)}=4^{1p}2^{12p}b_2^{(p)}+\frac{1}{3}4^{1p}b_3^{(p)}4^pb_4^{(p)}\hfill \\ A_2^{(p)}=2^{32p}b_2^{(p)}2^{32p}b_3^{(p)}+\frac{11}{3}2^{12p}b_4^{(p)}\hfill \\ A_3^{(p)}=\frac{1}{3}2^{52p}b_3^{(p)}4^{2p}b_4^{(p)}\hfill \\ A_4^{(p)}=\frac{1}{3}2^{52p}b_4^{(p)}\hfill \end{array}$$
(170)
We note that for all p
$$\underset{m=1}{\overset{4}{}}A_m^{(p)}=1$$
(171)
In table 2 we give the values of the $`A_m^{(p)}`$โs for $`1p8`$
| $`p`$ | $`A_1^{(p)}`$ | $`A_2^{(p)}`$ | $`A_3^{(p)}`$ | $`A_4^{(p)}`$ |
| --- | --- | --- | --- | --- |
| $`1`$ | $`1`$ | $`0`$ | $`0`$ | $`0`$ |
| $`2`$ | $`0`$ | $`1`$ | $`0`$ | $`0`$ |
| $`3`$ | $`0`$ | $`0`$ | $`1`$ | $`0`$ |
| $`4`$ | $`0`$ | $`0`$ | $`0`$ | $`1`$ |
| $`5`$ | $`\frac{3}{32}`$ | $`\frac{25}{32}`$ | $`\frac{35}{16}`$ | $`\frac{5}{2}`$ |
| $`6`$ | $`\frac{15}{64}`$ | $`\frac{119}{64}`$ | $`\frac{75}{16}`$ | $`\frac{65}{16}`$ |
| $`7`$ | $`\frac{195}{512}`$ | $`\frac{1505}{512}`$ | $`\frac{1799}{256}`$ | $`\frac{175}{32}`$ |
| $`8`$ | $`\frac{525}{1024}`$ | $`\frac{3985}{1024}`$ | $`\frac{1155}{128}`$ | $`\frac{1701}{256}`$ |
### A.3 Three dimensions
We start from the equation
$`[\nu ^\alpha (i)]^p`$ $`={\displaystyle \frac{1}{6^p}}{\displaystyle \underset{m=0}{\overset{p}{}}}\left(\begin{array}{c}p\hfill \\ m\hfill \end{array}\right)\nu (i_1)^{pm}{\displaystyle \underset{n=0}{\overset{m}{}}}\left(\begin{array}{c}m\hfill \\ n\hfill \end{array}\right)`$ (176)
$`\times \nu (i_2)^{mn}{\displaystyle \underset{l=0}{\overset{n}{}}}\left(\begin{array}{c}n\hfill \\ l\hfill \end{array}\right)\nu (i_3)^{nl}`$ (179)
$`\times {\displaystyle \underset{k=0}{\overset{l}{}}}\left(\begin{array}{c}l\hfill \\ k\hfill \end{array}\right)\nu (i_4)^{lk}{\displaystyle \underset{q=0}{\overset{k}{}}}\left(\begin{array}{c}k\hfill \\ q\hfill \end{array}\right)\nu (i_5)^{kq}\nu (i_6)^q`$ (184)
Because of the algebraic relations (148) we obtain
$$[\nu ^\alpha (i)]^p=\frac{1}{4^p}\underset{m=1}{\overset{6}{}}b_m^{(p)}Z_m$$
(185)
where the operators $`Z_m`$ are defined as
$`Z_1`$ $`=6\nu ^\alpha (i)`$ (186)
$`Z_2`$ $`=\nu _1\nu _2+\nu _1\nu _3+\nu _2\nu _3+\nu _1\nu _4+\nu _2\nu _4+\nu _3\nu _4`$
$`+\nu _1\nu _5+\nu _2\nu _5+\nu _3\nu _5+\nu _4\nu _5+\nu _1\nu _6`$
$`+\nu _2\nu _6+\nu _3\nu _6+\nu _4\nu _6+\nu _5\nu _6`$ (187)
$`Z_3`$ $`=\nu _1\nu _2\nu _3+\nu _1\nu _2\nu _4+\nu _1\nu _3\nu _4+\nu _2\nu _3\nu _4+\nu _1\nu _2\nu _5`$
$`+\nu _1\nu _3\nu _5+\nu _2\nu _3\nu _5+\nu _1\nu _4\nu _5+\nu _2\nu _4\nu _5+\nu _3\nu _4\nu _5`$
$`+\nu _1\nu _2\nu _6+\nu _1\nu _3\nu _6+\nu _2\nu _3\nu _6+\nu _1\nu _4\nu _6+\nu _2\nu _4\nu _6`$
$`+\nu _3\nu _4\nu _6+\nu _1\nu _5\nu _6+\nu _2\nu _5\nu _6+\nu _3\nu _5\nu _6+\nu _4\nu _5\nu _6`$ (188)
$`Z_4`$ $`=\nu _1\nu _2\nu _3\nu _4+\nu _1\nu _2\nu _3\nu _5+\nu _1\nu _2\nu _4\nu _5+\nu _1\nu _3\nu _4\nu _5`$
$`+\nu _2\nu _3\nu _4\nu _5+\nu _1\nu _2\nu _3\nu _6+\nu _1\nu _2\nu _4\nu _6+\nu _1\nu _3\nu _4\nu _6`$
$`+\nu _2\nu _3\nu _4\nu _6+\nu _1\nu _2\nu _5\nu _6+\nu _1\nu _3\nu _5\nu _6+\nu _2\nu _3\nu _5\nu _6`$
$`+\nu _1\nu _4\nu _5\nu _6+\nu _2\nu _4\nu _5\nu _6+\nu _3\nu _4\nu _5\nu _6`$ (189)
$`Z_5`$ $`=\nu _1\nu _2\nu _3\nu _4\nu _5+\nu _1\nu _2\nu _3\nu _4\nu _6+\nu _1\nu _2\nu _3\nu _5\nu _6`$
$`+\nu _1\nu _2\nu _4\nu _5\nu _6+\nu _1\nu _3\nu _4\nu _5\nu _6+\nu _2\nu _3\nu _4\nu _5\nu _6`$
$`Z_6`$ $`=\nu _1\nu _2\nu _3\nu _4\nu _5\nu _6`$ (190)
and the new coefficients $`b_m^{(p)}(p=5,6)`$ have the expressions
$$\begin{array}{c}b_5^{(p)}=\underset{m=2}{\overset{p1}{}}\left(\begin{array}{c}p\hfill \\ m\hfill \end{array}\right)\underset{n=1}{\overset{m1}{}}\left(\begin{array}{c}m\hfill \\ n\hfill \end{array}\right)a_{pm}a_{mn}\hfill \\ =5(12^{p+1}+23^p4^p+5^{p1})\hfill \\ b_6^{(p)}=\underset{m=2}{\overset{p1}{}}\left(\begin{array}{c}p\hfill \\ m\hfill \end{array}\right)\underset{n=1}{\overset{m1}{}}\left(\begin{array}{c}m\hfill \\ n\hfill \end{array}\right)a_{pm}a_{mn}a_n\hfill \\ =6+152^p+2^{3+2p}203^p+74^p65^p+6^p\hfill \end{array}$$
(191)
By solving the system (185) with respect to variables $`Z_m`$, we can obtain the recursion rule
$$[\nu ^\alpha (i)]^p=\underset{m=1}{\overset{6}{}}A_m^{(p)}[\nu ^\alpha (i)]^m$$
(192)
where the coefficients $`A_m^{(p)}`$ are defined as
$$\begin{array}{c}A_1^{(p)}=\frac{1}{6^p}[63b_2^{(p)}+2b_3^{(p)}\frac{3}{2}b_4^{(p)}+\frac{6}{5}b_5^{(p)}b_6^{(p)}]\hfill \\ A_2^{(p)}=\frac{1}{6^p}[18b_2^{(p)}18b_3^{(p)}+\frac{33}{2}b_4^{(p)}15b_5^{(p)}+\frac{137}{10}b_6^{(p)}]\hfill \\ A_3^{(p)}=\frac{1}{6^p}[36b_3^{(p)}54b_4^{(p)}+63b_5^{(p)}\frac{135}{2}b_6^{(p)}]\hfill \\ A_4^{(p)}=\frac{1}{6^p}[54b_4^{(p)}108b_5^{(p)}+153b_6^{(p)}]\hfill \\ A_5^{(p)}=\frac{1}{6^p}[\frac{324}{5}b_5^{(p)}162b_6^{(p)}]\hfill \\ A_6^{(p)}=\frac{1}{6^p}\frac{324}{5}b_6^{(p)}\hfill \end{array}$$
(193)
We note that for all p
$$\underset{m=1}{\overset{4}{}}A_m^{(p)}=1$$
(194)
In table 3 we give the values of the $`A_m^{(p)}`$โs for $`1p10`$.
| $`p`$ | $`A_1^{(p)}`$ | $`A_2^{(p)}`$ | $`A_3^{(p)}`$ | $`A_4^{(p)}`$ | $`A_5^{(p)}`$ | $`A_6^{(p)}`$ |
| --- | --- | --- | --- | --- | --- | --- |
| $`1`$ | $`1`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`2`$ | $`0`$ | $`1`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`3`$ | $`0`$ | $`0`$ | $`1`$ | $`0`$ | $`0`$ | $`0`$ |
| $`4`$ | $`0`$ | $`0`$ | $`0`$ | $`1`$ | $`0`$ | $`0`$ |
| $`5`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`1`$ | $`0`$ |
| $`6`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`1`$ |
| $`7`$ | $`\frac{5}{324}`$ | $`\frac{49}{216}`$ | $`\frac{203}{162}`$ | $`\frac{245}{172}`$ | $`\frac{175}{36}`$ | $`\frac{7}{2}`$ |
| $`8`$ | $`\frac{35}{648}`$ | $`\frac{1009}{1296}`$ | $`\frac{2695}{648}`$ | $`\frac{13811}{1296}`$ | $`\frac{245}{18}`$ | $`\frac{133}{18}`$ |
| $`9`$ | $`\frac{665}{5832}`$ | $`\frac{6307}{3888}`$ | $`\frac{98915}{11664}`$ | $`\frac{9065}{432}`$ | $`\frac{10913}{432}`$ | $`\frac{49}{4}`$ |
| $`10`$ | $`\frac{245}{1296}`$ | $`\frac{62167}{23328}`$ | $`\frac{53375}{3888}`$ | $`\frac{774575}{23328}`$ | $`\frac{4165}{108}`$ | $`\frac{7609}{432}`$ |
## Appendix B The energy matrix
The energy matrix $`ฯต^{(d)}`$, defined by Eq. (26) can be immediately calculated by means of the equation of motion (21) and the recurrence rule (23) \[see Tables 1, 2, 3\]. The matrix $`\mathrm{\Omega }^{(d)}`$ is defined as the matrix whose columns are the eigenvectors of the matrix $`ฯต^{(d)}`$. In this Appendix we report the expressions of $`ฯต^{(d)}`$ and $`\mathrm{\Omega }^{(d)}`$ for the various dimensions.
### B.1 One dimension
$$ฯต^{(1)}=\left(\begin{array}{ccc}\mu & 2V& 0\\ 0& \mu & 2V\\ 0& V& 3V\mu \end{array}\right)\mathrm{\Omega }^{(1)}=\left(\begin{array}{ccc}1& 2^2& 1\\ 0& 2& 1\\ 0& 1& 1\end{array}\right)$$
(195)
### B.2 Two dimensions
$$ฯต^{(2)}=\left(\begin{array}{ccccc}\mu & 4V& 0& 0& 0\\ 0& \mu & 4V& 0& 0\\ 0& 0& \mu & 4V& 0\\ 0& 0& 0& \mu & 4V\\ 0& \frac{3}{8}V& \frac{25}{8}V& \frac{35}{4}V& 10V\mu \end{array}\right)$$
(196)
$$\mathrm{\Omega }^{(2)}=\left(\begin{array}{ccccc}1& 4^4& 2^4& (\frac{4}{3})^4& 1\\ 0& 4^3& 2^3& (\frac{4}{3})^3& 1\\ 0& 4^2& 2^2& (\frac{4}{3})^2& 1\\ 0& 4& 2& (\frac{4}{3})& 1\\ 0& 1& 1& 1& 1\end{array}\right)$$
(197)
### B.3 Three dimensions
$$ฯต^{(3)}=\left(\begin{array}{ccccccc}\mu & 6V& 0& 0& 0& 0& 0\\ 0& \mu & 6V& 0& 0& 0& 0\\ 0& 0& \mu & 6V& 0& 0& 0\\ 0& 0& 0& \mu & 6V& 0& 0\\ 0& 0& 0& 0& \mu & 6V& 0\\ 0& 0& 0& 0& 0& \mu & 6V\\ 0& \frac{5}{54}V& \frac{49}{36}V& \frac{203}{27}V& \frac{245}{12}V& \frac{175}{6}V& 21V\mu \end{array}\right)$$
(198)
$$\mathrm{\Omega }^{(3)}=\left(\begin{array}{ccccccc}1& 6^6& 3^6& 2^6& \left(\frac{3}{2}\right)^6& \left(\frac{6}{5}\right)^6& 1\\ 0& 6^5& 3^5& 2^5& \left(\frac{3}{2}\right)^5& \left(\frac{6}{5}\right)^5& 1\\ 0& 6^4& 3^4& 2^4& \left(\frac{3}{2}\right)^4& \left(\frac{6}{5}\right)^4& 1\\ 0& 6^3& 3^3& 2^3& \left(\frac{3}{2}\right)^3& \left(\frac{6}{5}\right)^3& 1\\ 0& 6^2& 3^2& 2^2& \left(\frac{3}{2}\right)^2& \left(\frac{6}{5}\right)^2& 1\\ 0& 6& 3& 2& \left(\frac{3}{2}\right)& \left(\frac{6}{5}\right)& 1\\ 0& 1& 1& 1& 1& 1& 1\end{array}\right)$$
(199)
## Appendix C The normalization matrix
We recall the definition of the normalization matrix
$`I^{(d)}(๐ข,๐ฃ)`$ $`=`$ $`\{\psi ^{(d)}(๐ข,t),\psi ^{(d)}{}_{}{}^{}(๐ฃ,t)\}`$ (200)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{๐ค}{}}e^{\text{i}๐ค(๐_i๐_j)}I^{(d)}(๐ค)`$
It is straightforward to see that use of the hermiticity condition (43 ) leads to the fact that we have to calculate only the matrix elements $`I_{1,m}^{(d)}(๐ค)(m=1,2,\mathrm{}2d+1)`$. The calculations of these is very easy when one observes the following anticommutating rule
$$\begin{array}{c}\{c(๐ข,t)[\nu ^\alpha (i)]^p,c^{}(๐ฃ,t)\}=\delta _{\mathrm{๐ข๐ฃ}}[\nu ^\alpha (i)]^p\hfill \\ \underset{n=1}{\overset{p}{}}(1)^n\frac{1}{(2d)^{n1}}\left(\begin{array}{c}p\\ n\end{array}\right)\alpha _{\mathrm{๐ข๐ฃ}}c(๐ข,t)[\nu ^\alpha (i)]^{pn}c^{}(๐ฃ,t)\hfill \end{array}$$
(201)
By taking the expectation value of (201) we obtain in momentum space
$`I_{1,m}^{(d)}(๐ค)`$ $`=`$ $`\kappa ^{(m1)}\alpha (๐ค){\displaystyle \underset{n=1}{\overset{m1}{}}}(1)^n`$ (204)
$`\times {\displaystyle \frac{1}{(2d)^{n1}}}\left(\begin{array}{c}m1\\ n\end{array}\right)C_{1,mn}^{(d)\alpha }`$
with the definitions
$$C^{(d)\alpha }=\psi ^{(d)\alpha }(i)\psi ^{(d)}(i)\kappa ^{(p)}=[\nu ^\alpha (i)]^p$$
(205)
### C.1 One dimension
$$I^{(1)}(๐ค)=\left(\begin{array}{ccc}I_{1,1}^{(1)}& I_{1,2}^{(1)}& I_{1,3}^{(1)}\\ I_{1,2}^{(1)}& I_{1,3}^{(1)}& I_{2,3}^{(1)}\\ I_{1,3}^{(1)}& I_{2,3}^{(1)}& I_{3,3}^{(1)}\end{array}\right)$$
(206)
where
$$\begin{array}{c}I_{1,1}^{(1)}(๐ค)=1\hfill \\ I_{1,m}^{(1)}(๐ค)=\kappa ^{(m1)}\alpha (๐ค)\underset{n=1}{\overset{m1}{}}(1)^n\frac{1}{(2)^{n1}}\hfill \\ \times \left(\begin{array}{c}m1\\ n\end{array}\right)C_{1,mn}^{(d)\alpha }(m=2,3)\hfill \\ I_{m,3}^{(1)}(๐ค)=\underset{n=1}{\overset{2}{}}A_n^{(m+1)}I_{1,n+1}^{(1)}(๐ค)(m=2,3)\hfill \end{array}$$
(207)
### C.2 Two dimensions
$$I^{(2)}(๐ค)=\left(\begin{array}{ccccc}I_{1,1}^{(2)}& I_{1,2}^{(2)}& I_{1,3}^{(2)}& I_{1,4}^{(2)}& I_{1,5}^{(2)}\\ I_{1,2}^{(2)}& I_{1,3}^{(2)}& I_{1,4}^{(2)}& I_{1,5}^{(2)}& I_{2,5}^{(2)}\\ I_{1,3}^{(2)}& I_{1,4}^{(2)}& I_{1,5}^{(2)}& I_{2,5}^{(2)}& I_{3,5}^{(2)}\\ I_{1,4}^{(2)}& I_{1,5}^{(2)}& I_{2,5}^{(2)}& I_{3,5}^{(2)}& I_{4,5}^{(2)}\\ I_{1,5}^{(2)}& I_{2,5}^{(2)}& I_{3,5}^{(2)}& I_{4,5}^{(2)}& I_{5,5}^{(2)}\end{array}\right)$$
(208)
where
$$\begin{array}{c}I_{1,1}^{(2)}(๐ค)=1\hfill \\ I_{1,m}^{(2)}(๐ค)=\kappa ^{(m1)}\alpha (๐ค)\underset{n=1}{\overset{m1}{}}(1)^n\frac{1}{(4)^{n1}}\hfill \\ \times \left(\begin{array}{c}m1\\ n\end{array}\right)C_{1,mn}^{(d)\alpha }(m=2,\mathrm{}5)\hfill \\ I_{m,5}^{(2)}(๐ค)=\underset{n=1}{\overset{4}{}}A_n^{(m+3)}I_{1,n+1}^{(2)}(๐ค)(m=2,\mathrm{}5)\hfill \end{array}$$
(209)
### C.3 Three dimensions
$$I^{(3)}(๐ค)=\left(\begin{array}{ccccccc}I_{1,1}^{(3)}& I_{1,2}^{(3)}& I_{1,3}^{(3)}& I_{1,4}^{(3)}& I_{1,5}^{(3)}& I_{1,6}^{(3)}& I_{1,7}^{(3)}\\ I_{1,2}^{(3)}& I_{1,3}^{(3)}& I_{1,4}^{(3)}& I_{1,5}^{(3)}& I_{1,6}^{(3)}& I_{1,7}^{(3)}& I_{2,7}^{(3)}\\ I_{1,3}^{(3)}& I_{1,4}^{(3)}& I_{1,5}^{(3)}& I_{1,6}^{(3)}& I_{1,7}^{(3)}& I_{2,7}^{(3)}& I_{3,7}^{(3)}\\ I_{1,4}^{(3)}& I_{1,5}^{(3)}& I_{1,6}^{(3)}& I_{1,7}^{(3)}& I_{2,7}^{(3)}& I_{3,7}^{(3)}& I_{4,7}^{(3)}\\ I_{1,5}^{(3)}& I_{1,6}^{(3)}& I_{1,7}^{(3)}& I_{2,7}^{(3)}& I_{3,7}^{(3)}& I_{4,7}^{(3)}& I_{5,7}^{(3)}\\ I_{1,6}^{(3)}& I_{1,7}^{(3)}& I_{2,7}^{(3)}& I_{3,7}^{(3)}& I_{4,7}^{(3)}& I_{5,7}^{(3)}& I_{6,7}^{(3)}\\ I_{1,7}^{(3)}& I_{2,7}^{(3)}& I_{3,7}^{(3)}& I_{4,7}^{(3)}& I_{5,7}^{(3)}& I_{6,7}^{(3)}& I_{7,7}^{(3)}\end{array}\right)$$
(210)
where
$$\begin{array}{c}I_{1,1}^{(3)}(๐ค)=1\hfill \\ I_{1,m}^{(3)}(๐ค)=\kappa ^{(m1)}\alpha (๐ค)\underset{n=1}{\overset{m1}{}}(1)^n\frac{1}{(6)^{n1}}\hfill \\ \times \left(\begin{array}{c}m1\\ n\end{array}\right)C_{1,mn}^{(d)\alpha }(m=2,\mathrm{}7)\hfill \\ I_{m,7}^{(3)}(๐ค)=\underset{n=1}{\overset{6}{}}A_n^{(m+5)}I_{1,n+1}^{(3)}(๐ค)(m=2,\mathrm{}7)\hfill \end{array}$$
(211)
## Appendix D The spectral matrices
The spectral density matrices $`\sigma _{ab}^{(d,n)}(๐ค)`$ can be immediately calculated by means of the knowledge of the matrices $`\mathrm{\Omega }^{(d)}`$ and $`I^{(d)}`$ through Eq. (33).
### D.1 One dimension
$`\sigma ^{(1)}`$ $`=\mathrm{\Sigma }_1\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)\sigma ^{(2)}=\mathrm{\Sigma }_2\left(\begin{array}{ccc}1& 2^1& 2^2\\ 2^1& 2^2& 2^3\\ 2^2& 2^3& 2^4\end{array}\right)`$ (218)
$`\sigma ^{(3)}`$ $`=\mathrm{\Sigma }_3\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right)`$ (222)
where
$$\begin{array}{c}\mathrm{\Sigma }_1=I_{1,1}3I_{1,2}+2I_{1,3}\hfill \\ \mathrm{\Sigma }_2=4(I_{1,2}I_{1,3})\hfill \\ \mathrm{\Sigma }_3=I_{1,2}+2I_{1,3}\hfill \end{array}$$
(223)
### D.2 Two dimensions
$`\sigma ^{(1)}`$ $`=`$ $`\mathrm{\Sigma }_1\left(\begin{array}{ccccc}1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\end{array}\right)`$ (229)
$`\sigma ^{(2)}`$ $`=`$ $`\mathrm{\Sigma }_2\left(\begin{array}{ccccc}1& 4^1& 4^2& 4^3& 4^4\\ 4^1& 4^2& 4^3& 4^4& 4^5\\ 4^2& 4^3& 4^4& 4^5& 4^6\\ 4^3& 4^4& 4^5& 4^6& 4^7\\ 4^4& 4^5& 4^6& 4^7& 4^8\end{array}\right)`$ (235)
$$\sigma ^{(3)}=\mathrm{\Sigma }_3\left(\begin{array}{ccccc}1& 2^1& 2^2& 2^3& 2^4\\ 2^1& 2^2& 2^3& 2^4& 2^5\\ 2^2& 2^3& 2^4& 2^5& 2^6\\ 2^3& 2^4& 2^5& 2^6& 2^7\\ 2^4& 2^5& 2^6& 2^7& 2^8\end{array}\right)$$
(236)
$`\sigma ^{(4)}`$ $`=`$ $`\mathrm{\Sigma }_4\left(\begin{array}{ccccc}1& \left(\frac{4}{3}\right)^1& \left(\frac{4}{3}\right)^2& \left(\frac{4}{3}\right)^3& \left(\frac{4}{3}\right)^4\\ \left(\frac{4}{3}\right)^1& \left(\frac{4}{3}\right)^2& \left(\frac{4}{3}\right)^3& \left(\frac{4}{3}\right)^4& \left(\frac{4}{3}\right)^5\\ \left(\frac{4}{3}\right)^2& \left(\frac{4}{3}\right)^3& \left(\frac{4}{3}\right)^4& \left(\frac{4}{3}\right)^5& \left(\frac{4}{3}\right)^6\\ \left(\frac{4}{3}\right)^3& \left(\frac{4}{3}\right)^4& \left(\frac{4}{3}\right)^5& \left(\frac{4}{3}\right)^6& \left(\frac{4}{3}\right)^7\\ \left(\frac{4}{3}\right)^4& \left(\frac{4}{3}\right)^5& \left(\frac{4}{3}\right)^6& \left(\frac{4}{3}\right)^7& \left(\frac{4}{3}\right)^8\end{array}\right)`$ (242)
$`\sigma ^{(5)}`$ $`=`$ $`\mathrm{\Sigma }_5\left(\begin{array}{ccccc}1& 1& 1& 1& 1\\ 1& 1& 1& 1& 1\\ 1& 1& 1& 1& 1\\ 1& 1& 1& 1& 1\\ 1& 1& 1& 1& 1\end{array}\right)`$ (248)
with
$$\begin{array}{c}\mathrm{\Sigma }_1=I_{1,1}\frac{1}{3}(25I_{1,2}70I_{1,3}+80I_{1,4}32I_{1,5})\hfill \\ \mathrm{\Sigma }_2=\frac{16}{3}(3I_{1,2}13I_{1,3}+18I_{1,4}8I_{1,5})\hfill \\ \mathrm{\Sigma }_3=4(3I_{1,2}19I_{1,3}+32I_{1,4}16I_{1,5})\hfill \\ \mathrm{\Sigma }_4=\frac{16}{3}(I_{1,2}7I_{1,3}+14I_{1,4}8I_{1,5})\hfill \\ \mathrm{\Sigma }_5=\frac{1}{3}(3I_{1,2}22I_{1,3}+48I_{1,4}32I_{1,5})\hfill \end{array}$$
(249)
### D.3 Three dimensions
$`\sigma ^{(1)}`$ $`=`$ $`\mathrm{\Sigma }_1\left(\begin{array}{cccc}1& 0& \mathrm{}& 0\\ 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0\end{array}\right)`$ (254)
$`\sigma ^{(2)}`$ $`=`$ $`\mathrm{\Sigma }_2\left(\begin{array}{cccc}1& 6^1& \mathrm{}& 6^6\\ 6^1& 6^2& \mathrm{}& 6^7\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 6^6& 6^7& \mathrm{}& 6^{12}\end{array}\right)`$ (259)
$`\sigma ^{(3)}`$ $`=`$ $`\mathrm{\Sigma }_3\left(\begin{array}{cccc}1& 3^1& \mathrm{}& 3^6\\ 3^1& 3^2& \mathrm{}& 3^7\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 3^6& 3^7& \mathrm{}& 3^{12}\end{array}\right)`$ (264)
$`\sigma ^{(4)}`$ $`=`$ $`\mathrm{\Sigma }_4\left(\begin{array}{cccc}1& 2^1& \mathrm{}& 2^6\\ 2^1& 2^2& \mathrm{}& 2^7\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 2^6& 2^7& \mathrm{}& 2^{12}\end{array}\right)`$ (269)
$$\sigma ^{(5)}=\mathrm{\Sigma }_5\left(\begin{array}{cccc}1& \left(\frac{3}{2}\right)^1& \mathrm{}& \left(\frac{3}{2}\right)^6\\ \left(\frac{3}{2}\right)^1& \left(\frac{3}{2}\right)^2& \mathrm{}& \left(\frac{3}{2}\right)^7\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \left(\frac{3}{2}\right)^6& \left(\frac{3}{2}\right)^7& \mathrm{}& \left(\frac{3}{2}\right)^{12}\end{array}\right)$$
(270)
$`\sigma ^{(6)}`$ $`=`$ $`\mathrm{\Sigma }_6\left(\begin{array}{cccc}1& \left(\frac{6}{5}\right)^1& \mathrm{}& \left(\frac{6}{5}\right)^6\\ \left(\frac{6}{5}\right)^1& \left(\frac{6}{5}\right)^2& \mathrm{}& \left(\frac{6}{5}\right)^7\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \left(\frac{6}{5}\right)^6& \left(\frac{6}{5}\right)^7& \mathrm{}& \left(\frac{6}{5}\right)^{12}\end{array}\right)`$ (275)
$`\sigma ^{(7)}`$ $`=`$ $`\mathrm{\Sigma }_7\left(\begin{array}{cccc}1& 1& \mathrm{}& 1\\ 1& 1& \mathrm{}& 1\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& 1& \mathrm{}& 1\end{array}\right)`$ (280)
here
$$\mathrm{\Sigma }_p=\underset{m=1}{\overset{7}{}}B_m^{(p)}I_{1,m}^{(3)}$$
(281)
The coefficients $`B_m^{(p)}`$are given in Table 4.
| $`p`$ | $`B_1^{(p)}`$ | $`B_2^{(p)}`$ | $`B_3^{(p)}`$ | $`B_4^{(p)}`$ | $`B_5^{(p)}`$ | $`B_6^{(p)}`$ | $`B_7^{(p)}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| $`1`$ | $`1`$ | $`\frac{147}{10}`$ | $`\frac{406}{5}`$ | $`\frac{441}{2}`$ | $`315`$ | $`\frac{1134}{5}`$ | $`\frac{324}{5}`$ |
| $`2`$ | $`0`$ | $`36`$ | $`\frac{1566}{5}`$ | $`1044`$ | $`1674`$ | $`1296`$ | $`\frac{1944}{5}`$ |
| $`3`$ | $`0`$ | $`45`$ | $`\frac{1053}{2}`$ | $`\frac{4149}{2}`$ | $`3699`$ | $`3078`$ | $`972`$ |
| $`4`$ | $`0`$ | $`40`$ | $`508`$ | $`2232`$ | $`4356`$ | $`3888`$ | $`1296`$ |
| $`5`$ | $`0`$ | $`\frac{45}{2}`$ | $`297`$ | $`\frac{2763}{2}`$ | $`2889`$ | $`2754`$ | $`972`$ |
| $`6`$ | $`0`$ | $`\frac{36}{5}`$ | $`\frac{486}{5}`$ | $`468`$ | $`1026`$ | $`\frac{5184}{5}`$ | $`\frac{1944}{5}`$ |
| $`7`$ | $`0`$ | $`1`$ | $`\frac{137}{10}`$ | $`\frac{135}{2}`$ | $`153`$ | $`162`$ | $`\frac{324}{5}`$ |
## Appendix E Relations between the Ising and spinless model
In this Appendix, we want to recall the main results of the Ising model and establish the relations between the two models. For an infinite chain the simplest method is the use of the transfer matrix method. The details of calculations are well known and can be found in many textbooks. For example, we refer the reader to Refs.Baxter (1982); Goldenfel (1992); Lavis and Bell (1999). The magnetization per site is given by
$$S(i)=\frac{\mathrm{sinh}(\beta h)}{\sqrt{\mathrm{sinh}^2(\beta h)+e^{4\beta J}}}$$
(282)
The two-point correlation function $`S(i)S(i+j)`$ has the expression
$$S(i)S(i+j)=S(i)^2+(1S(i)^2)p^j$$
(283)
where
$$p=\frac{\gamma ^{(2)}}{\gamma ^{(1)}}$$
(284)
$`\gamma ^{(1)\text{ }}`$and $`\gamma ^{(2)}`$ are the eigenvalues of the transfer matrix
$$\begin{array}{c}\gamma ^{(1)}=e^{\beta J}[\mathrm{cosh}(\beta h)+\sqrt{\mathrm{sinh}^2(\beta h)+e^{4\beta J}}]\hfill \\ \gamma ^{(2)}=e^{\beta J}\left[\mathrm{cosh}(\beta h)\sqrt{\mathrm{sinh}^2(\beta h)+e^{4\beta J}}\right]\hfill \end{array}$$
(285)
The three-point correlation function is given byMarsh (1966)
$$\begin{array}{c}S(i)S(i+j)S(i+j+r)=S(i)^3\hfill \\ +S(i)[1S(i)^2](p^j+p^rp^{j+r})\hfill \end{array}$$
(286)
The relations between the Ising and fermionic models are
$$\begin{array}{c}h=\frac{1}{2}(\mu Vd)\hfill \\ J=\frac{1}{4}Vd\hfill \end{array}$$
(287)
$$\nu (i)=\frac{1}{2}[1+S(i)]$$
(288)
$$\lambda ^{(1)}=\frac{1}{4}[1+2S(i)+S(i)S(i+a)]$$
(289)
$$\kappa ^{(2)}=\frac{1}{8}[3+4S(i)+S(i)S(i+2a)]$$
(290)
$`\lambda ^{(2)}`$ $`={\displaystyle \frac{1}{16}}[3+7S(i)+4S(i)S(i+a)`$
$`+S(i)S(i+2a)+S(i)S(i+a)S(i+2a)]`$ (291)
By recalling (116) and by means of (288), the magnetization in the fermionic model has the expression
$$S(i)=(1T_2)\sqrt{\frac{T_1}{T_12T_1T_2+2T_2^2}}$$
(292)
By oserving that $`T_1`$ and $`T_2`$ can be expressed as
$$\begin{array}{c}T_1=1\mathrm{tanh}\left(\frac{\beta \mu }{2}\right)=\frac{2e^{2\beta h}}{e^{2\beta h}+e^{4\beta J}}\hfill \\ T_2=1\mathrm{tanh}\left(\frac{\beta (\mu V)}{2}\right)=1\mathrm{tanh}(\beta h)=\frac{2}{e^{2\beta h}+1}\hfill \end{array}$$
(293)
it is straigtforward to see that (292) is the same as (282). In the fermionic model, by means of (142), we have
$$\frac{\nu (i)\nu (i+j)\nu ^2}{\nu \nu ^2}=\frac{S(i)S(i+j)S(i)^2}{1S(i)^2}=p^j$$
(294)
where the parameter $`p`$ is expressed in terms of $`\nu `$ and $`\lambda ^{(1)}`$ by means of (138). By using (VI) and (116), and by recalling (115) and (293), it is easy to see that the expression of $`p`$, given by (138) is exactly equal to the expression (284). Then, the two-point correlation function of the fermionic model exactly agree with the expression (283) of the Ising model. The parameters $`\kappa ^{(2)}`$, $`\lambda ^{(1)}`$ and $`\lambda ^{(2)}`$ can be calculated in the fermionic model by putting (116) into (112)-( VI), and in the Ising model by means of (289)-(291). After lengthy, but straightforward, calculations, using the relations (115) and (293), it is possible to show that there is an exact agreement. |
warning/0506/astro-ph0506336.html | ar5iv | text | # Hot Stars in Old Stellar Populations: A Continuing Need for Intermediate AgesThe data presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation.
## 1 Introduction
Stellar population analysis offers a powerful, if difficult to interpret, method of understanding the formation histories of nearby early-type galaxies (see Rose, 1985; Gonzรกlez, 1993; Trager et al., 2000b; Caldwell, Rose & Concannon, 2003; Mehlert et al., 2003, for just a few examples). This analysis relies primarily on the comparison of a hydrogen Balmer absorption-line strength to a metal absorption-line strength (or a combination of metal lines) to break the ageโmetallicity degeneracy (Worthey, 1994), as the Balmer lines are (non-linearly) sensitive to the temperature of the main-sequence turnoff and the metal lines are sensitive to the temperature of the red giant branch. One can therefore determine accurate ages for the old stellar populations found in early-type galaxies. However, other hot star populations such as blue horizontal branch stars or blue straggler stars can significantly increase the observed Balmer-line strengths of old stellar populations (see, e.g., Burstein et al., 1984; Rose, 1985; Rose & Tripicco, 1986; Rose, 1994; de Freitas Pacheco & Barbuy, 1995; Maraston & Thomas, 2000; Lee, Yoon & Lee, 2000; Trager et al., 2000a).
In this paper we explore the effect of a specific kind of hot star population, that is, old, metal-poor populations containing blue horizontal branch stars, on the inferred stellar population ages and compositions of early-type galaxies in the Coma Cluster. These galaxies appear have significant intermediate-aged populations due to their enhanced Balmer lines. We use blue indexes first described by Rose (1985, 1994) to determine the level of contamination of the galaxy spectra by blue horizontal branch (BHB) stars. We then subtract model spectra representing populations containing these stars from the observed spectra and determine ages, metallicities, and enhancement ratios from the residual spectra. Finally these stellar population parameters are compared with those determined from the observed spectra to quantify the effect of a contaminating population of old, metal-poor stars on the spectra of early-type galaxies.
Throughout this paper, we refer to populations with metallicities $`[\mathrm{Z}/\mathrm{H}]1.5`$ as โmetal-poorโ (and thus possessing BHB stars) and populations with ages $`1t10`$ Gyr as โintermediate agedโ.
## 2 Data
The line strengths discussed in this paper are derived from multi-slit spectra of twelve early-type galaxies in the Coma Cluster, centred on the cD galaxy NGC 4874, taken with the Low-Resolution Imaging Spectrograph (LRIS: Oke et al., 1995) on the Keck II 10-m Telescope. Details relevant to the current study are summarised here; for a complete description of the acquisition, reduction, and calibration of these spectra and the extraction of Lick/IDS absorption-line strengths we refer interested readers to Trager, Faber & Dressler (in preparation; hereafter TFD05).
### 2.1 Observations
Spectra were obtained in three consecutive 30-minute exposures on 7 April 1997 UT with the red side of LRIS, with seeing $`\mathrm{FWHM}0.8`$ arcsec, through clouds. A slit width of 1 arcsec was used in conjunction with the 600 line $`\mathrm{mm}^1`$ grating blazed at 5000 ร
, giving a resolution of 4.4 ร
FWHM ($`\sigma =1.9`$ ร
) and a wavelength coverage of typically 3500โ6000 ร
, depending on slit placement. Spectra of Lick/IDS standard G and K giant stars and F9โG0 dwarfs (Worthey et al., 1994) were observed on the same and subsequent nights through the LRIS 1 arcsec long slit using the same grating to be used for calibration to the Lick/IDS system (see Sec. 2.3 below).
Individual two-dimensional spectra of each galaxy were extracted from the multi-slit images after standard calibrations (overscan correction, bias removal, dark correction, and flat field correction), mapping of the geometric distortions, wavelength calibration, and sky subtraction<sup>1</sup><sup>1</sup>1For NGC 4874, which filled its slitlet, and for D128 and NGC 4872, whose spectra were contaminated by that of NGC 4874, sky subtraction was performed first using the โskyโ information at the edge of their slitlets and then corrected by comparing this sky spectrum to the average sky from all other slitlets. The excesses in these slits were added back into the final extracted spectra. were performed following the methodology of Kelson (2003). Both individual one-dimensional spectra and variance-weighted, combined spectra were then extracted from the two-dimensional spectra. In order to simulate an equivalent circular aperture to match with other line strength work in the Coma Cluster, the extracted spectra were weighted by distance from the object centre. For the present study the spectra were extracted with an equivalent circular diameter aperture of 2.7 arcsec, matching the Lick/IDS galaxy aperture (Trager et al., 1998) and the fibre diameter of the large sample of line strengths of Coma galaxies of Moore et al. (2002). Finally, the spectra were flux-calibrated using observations of spectrophotometric standard stars.
### 2.2 Line strengths on the Rose system
Rose (1985, 1994) has developed an absorption-line strength system, based on a combination of line-depth ratios and equivalent widths, which provides powerful tools for decomposing the hot-to-cool and giant-to-dwarf star ratios in composite populations. The Rose indexes of interest here are Ca ii (= Ca ii H+H$`ฯต`$/Ca ii K) and Hn/Fe ($`=\mathrm{H}\theta /\lambda 3859+\mathrm{H}\delta /\lambda 4045+\mathrm{H}\gamma /\lambda 4325`$, Caldwell & Rose 1998), which together indicate the presence of hot stars in old, metal-rich stellar populations (Rose, 1985, 1994; Caldwell et al., 2003).
The LRIS spectra were first smoothed to a common velocity dispersion of $`230\mathrm{km}\mathrm{s}^1`$ as suggested by Caldwell et al. (2003), except for NGC 4874, which was left unsmoothed at its intrinsic velocity dispersion of $`271\mathrm{km}\mathrm{s}^1`$. Indexes were then computed by finding the minimum intensity of each absorption line using a cubic spline interpolation of a small region around the centre of each line and then dividing the appropriate combination of lines to determine the index value. For example, the $`\mathrm{H}\delta /\lambda 4045`$ index is determined by finding the minimum intensity of the $`\mathrm{H}\delta `$ absorption line and dividing it by the minimum intensity of the Fe I $`\lambda 4045`$ absorption line. Indexes were determined independently from each of the three exposures after smoothing and the mean and sample standard deviations were used as the final index value and error, respectively. Within the errors of the two Rose indexes of interest (Hn/Fe and Ca ii), the indexes determined from the variance-weighted combined spectra and from the mean of the individual exposures are identical. However, the uncertainties determined from the standard deviations of the indexes determined from the individual exposures appear to be more reliable than those determined from estimation of the photon noise in the combined spectra (cf. Caldwell et al., 2003). Rose line strengths and errors for the current sample, measured in a synthesised circular aperture of 2.7-arcsec diameter, are presented in Table 1.
No correction for emission fill-in of the Hn/Fe index has been attempted. The fluxes of the H$`\alpha `$ emission in each galaxy required to make the correction using the recipe of Appendix A of Caldwell et al. (2003) are unknown, although we expect that the corrections will be small: the typical correction to Hn/Fe in Caldwell et al. (2003) is 0.012, which is a correction of less than 1 Gyr for an 8 Gyr-old population. The precise value of Hn/Fe is however not used in the analysis that follows.
Furthermore, we cannot calibrate our index strengths on to a Rose โsystemโ, as the Jones (1996) stellar library used by Caldwell et al. (2003) is smoothed to a velocity dispersion of 103 $`\mathrm{km}\mathrm{s}^1`$ (the intrinsic resolution of our stellar spectra is roughly 150 $`\mathrm{km}\mathrm{s}^1`$), and no published index strengths exist for the galaxies in the present study. However, the location of the Coma galaxies in the Hn/FeโCa ii diagram is coincident with the distribution of field and Virgo early-type galaxies in Caldwell et al. (2003), giving confidence that the moderate velocity smoothing is sufficient to calibrate the Coma galaxies on to a Rose-like โsystemโ (Fig. 1a).
### 2.3 Line strengths on the Lick/IDS system
The Lick/IDS absorption-line strength system has been developed by Faber, Burstein, and their collaborators (e.g. Burstein et al., 1984; Worthey et al., 1994; Trager et al., 1998) to determine the stellar content of early-type galaxies (as seen in the models of, e.g., Worthey, 1994; Thomas, Maraston & Bender, 2003). For the purpose of breaking the ageโmetallicity degeneracy inherent in colours and line strengths of old stellar populations (see, e.g., OโConnell, 1980), the $`\mathrm{H}\beta `$, $`\mathrm{Mg}b`$, Fe5270 and Fe5335 indexes are among the best-understood and best-calibrated (e.g., Trager et al., 2000a); we will use these four indexes to determine stellar population parameters in this study.
In order to measure Lick/IDS line strengths and then to place them on the Lick/IDS system, a series of spectral manipulations and calibrations are required. For the current observations, the complete series of steps performed is discussed in detail in TFD05; here we briefly review the process.
The first step in determining the line strengths of a galaxy is to measure its systemic velocity and velocity dispersion. These are needed to place the index bandpasses on the spectrum and to calibrate galaxy line strengths on to the Lick/IDS *stellar* system used by most stellar population models, including Worthey (1994). Both of these quantities are measured following the direct-fitting algorithm of Kelson et al. (2000) over the rest-frame wavelength range 4200โ5100 ร
.
Next, the spectrum is smoothed to the Lick/IDS resolution (Worthey & Ottaviani, 1997) using a variable-width Gaussian filter. The Lick/IDS index bandpasses are then placed on the spectrum and indexes measured in either ร
or magnitudes, depending roughly on the width of the central bandpasses (Trager et al., 1998). Errors are determined from the object spectrum and its associated variance spectrum.
For galaxies, a correction for fill-in of $`\mathrm{H}\beta `$ absorption by emission is required (Gonzรกlez, 1993). This is determined from the strength of the \[O iii\]$`\lambda 5007`$ line as determined from the residual spectrum after subtracting the best-fitting spectrum from the models of Vazdekis (1999) from the observed spectrum<sup>2</sup><sup>2</sup>2We use the Vazdekis (1999) models for the emission correction because we achieve lower $`\chi ^2`$ values in the spectral fitting (i.e., better fits) when using these models instead of the Worthey (1994) models. This is due to the denser coverage in age of the Vazdekis (1999) models with respect to the Worthey (1994) models.. The correction to $`\mathrm{H}\beta `$ is then determined using the \[O iii\]$`\lambda 5007`$$`\mathrm{H}\beta `$ correction suggested by Trager et al. (2000a): $`\mathrm{\Delta }\mathrm{H}\beta =0.6\times `$\[O iii\], where the \[O iii\] index is defined by Gonzรกlez (1993) and is positive for emission within the central bandpass. This correction never exceeds $`0.09`$ ร
for the present galaxies, or about 4% in $`\mathrm{H}\beta `$ strength.
Finally, two corrections are required to bring the galaxy indexes on to the Lick/IDS stellar system: small offsets resulting from the fact that the Lick/IDS system is not based on flux-calibrated spectra, and a velocity dispersion correction to account for the velocity broadening of the galaxies (Gonzรกlez, 1993; Trager et al., 1998). The former correction is performed using observations of the Lick/IDS stars taken in the run; the latter is performed using the polynomial corrections given in Trager et al. (1998). The final Lick/IDS index strengths of early-type galaxies in the Coma Cluster for the four lines of interest and for \[O iii\] are given in Table 1.
## 3 Analysis
Our goal is to determine the fractional contribution of hot stars, *assumed to arise from the blue horizontal branch stars of an old, metal-poor population,* to the light at 4000 ร
of early-type galaxies in the Coma Cluster and then to correct their age- and metallicity-sensitive line strengths for this contamination. We then compute the stellar population parameters age $`t_{\mathrm{SSP}}`$, metallicity $`[\mathrm{Z}/\mathrm{H}]_{\mathrm{SSP}}`$, and enhancement ratio $`[\mathrm{E}/\mathrm{Fe}]_{\mathrm{SSP}}`$ for each galaxy as observed and after correction for hot-star contamination. We can therefore determine the effect of these stars on the inferred stellar populations of early-type galaxies. Here SSP refers to the equivalent *single stellar population*, that is, a population of stars formed at the same time $`t_{\mathrm{SSP}}`$ with the same chemical composition $`[\mathrm{Z}/\mathrm{H}]_{\mathrm{SSP}}`$ and $`[\mathrm{E}/\mathrm{Fe}]_{\mathrm{SSP}}`$, with the same line strengths as the galaxy.
### 3.1 Models
We use the models of Worthey (1994, hereafter W94) to analyse the stellar populations of early-type galaxies, extended to cover non-solar abundance ratios using the Tripicco & Bell (1995) response functions as described in Trager et al. (2000a). In particularly, we use the โvanillaโ W94 models, with extended horizontal branches at low metallicities, described in Leitherer et al. (1996). While these models provide line strengths on the Lick/IDS system, they do not normally provide Rose index values nor spectra (as in, e.g., Vazdekis, 1999; Bruzual & Charlot, 2003). Leonardi & Worthey (2000, see also ) extended the W94 models to produce spectra using the empirical stellar spectra library of Jones (1996) and theoretical stellar spectra generated with the SYNTHE and ATLAS programs (R. L. Kurucz 1995, priv. comm.). We use these model spectra in our analysis of the Hn/FeโCa ii plane (Fig. 1) and in our correction of the observed early-type galaxy spectra for contamination by hot-star light.
To determine stellar population parameters, we fit observed (or corrected) line strengths to model line strengths using the method described in TFD05. Briefly, we determine stellar population parameters using a nonlinear least-squares code based on the Levenberg-Marquardt algorithm. Stellar population models are interpolated on the fly to produce model indexes which are compared to the observed (or corrected) indexes $`\mathrm{H}\beta `$, $`\mathrm{Mg}b`$, Fe5270 and Fe5335. Uncertainties are determined by taking the dispersion of stellar population parameters from 500 Monte Carlo trials using the errors of the observed line strengths, assuming that these errors are normally distributed.
### 3.2 Method
Our method is based on the suggestion of Rose (1985, 1994) and Caldwell et al. (2003) to use the Hn/FeโCa ii plane to determine the presence and amount of hot-star light at 4000 ร
in early-type galaxies. The W94 models predict an asymptotic Ca ii strength of 1.24 for *metal-rich* populations with ages $`t_{\mathrm{SSP}}2`$ Gyr (Fig. 1)<sup>3</sup><sup>3</sup>3The Vazdekis (1999) models have a slightly stronger asymptotic Ca ii strength of 1.25, which implies slightly higher hot-star fractions by typically about 1% of the light at 4000 ร
. However, we do not use the Vazdekis (1999) models for the analysis of the impact of the hot-star population on the stellar population parameters because these models (1) appear not to have the same flux calibration in the โblueโ and โredโ models and (2) have not been modified to account for enhancement-ratio variations.. A hot-star spectrum is first smoothed to instrumental resolution of LRIS and to the velocity dispersion of the galaxy in question and then subtracted from the observed spectrum in increments of $`f_{\mathrm{hot}}=0.005`$ (where $`f_{\mathrm{hot}}`$ is the fraction of light coming from the old, metal-poor population at 4000 ร
) until the measured Ca ii strength in the residual spectrum reaches the asymptotic old, metal-rich value. A demonstration of this process is given in Figure 2 for NGC 4867.
We choose three baseline hot-star models, whose line strengths are given in Table 2. The first is a 17 Gyr-old single stellar population model with $`[\mathrm{Z}/\mathrm{H}]=1.5`$ dex and therefore an extended horizontal branch running from red to blue; the assumed $`[\mathrm{E}/\mathrm{Fe}]=0`$ dex, that is to say, having the same $`[\mathrm{E}/\mathrm{Fe}]`$ as the calibrating stars (Thomas et al., 2003). As described in Trager et al. (2000a), the W94 model ages are older by 10โ25 per cent than models based on isochrones from the Padova group (Bertelli et al., 1994); this 17 Gyr-old model is equivalent in its line strengths to a 15 Gyr old model from, say, Thomas et al. (2003). The second is a younger, metal-poor single-stellar population with an age of 12 Gyr, $`[\mathrm{Z}/\mathrm{H}]=1.5`$ dex, and $`[\mathrm{E}/\mathrm{Fe}]=0`$ dex, also with an extended horizontal branch. The third is the observed spectrum of the globular cluster NGC 6254 (M10) from the compilation of Schiavon et al. (2005)<sup>4</sup><sup>4</sup>4http://www.astro.virginia.edu/ rps7v/GCs/intro.html, which has $`[\mathrm{Fe}/\mathrm{H}]=1.51`$ and an extended, blue horizontal branch (see, e.g., Rosenberg et al., 2000). The indexes given in Table 2 are derived from the spectrum of Schiavon et al. (2005) and corrected using the LRIS-based flux corrections.
We note here that Schiavon et al. (2004) present an alternative approach using the ratio $`\mathrm{H}\delta _\mathrm{F}/\mathrm{H}\beta `$ as a function of Fe4383. We have not attempted this method here, as it is dependent on knowing the response of $`\mathrm{H}\delta _\mathrm{F}`$ to $`\alpha `$-enhancements, which have only just become available and are currently in the process of being incorporated into stellar population models (Thomas et al., 2004; Korn et al., 2005; Lee & Worthey, 2005).
## 4 Results and Discussion
Figure 1b shows the results of the correction for hot-star light on the Hn/Fe and Ca ii strengths for three of the Coma galaxies. Table 3 gives the required fractions to correct the spectra for hot-star light for all twelve galaxies: these range from 3 per cent of the total light within the 2.7 arcsec aperture at 4000 ร
for NGC 4864 to 22 per cent for NGC 4867 when the hot-star light comes from the 17 Gyr old model, with a mean value of 7.9 per cent and an rms scatter of 2.9 per cent (here we use the bi-weight mean and scatter described in Beers, Flynn, & Gebhardt, 1990). When the 12 Gyr model is used, these fractions range from 3โ19 per cent with a mean of 7.0 per cent and an rms scatter of 2.6 per cent, and when NGC 6254 is used, the range is 5โ30 per cent with a mean of 10.2 per cent and an rms scatter of 3.7 per cent.
The inferred masses of the metal-poor components are also listed in Table 3; the bi-weight mean mass fractions of the metal-poor component are $`8.3\pm 6.1`$ per cent for the 17 Gyr old population and $`5.4\pm 3.9`$ per cent for the 12 Gyr old population. Note that the mass-to-light ratio of NGC 6254 is unknown and therefore the hot-star mass fractions are also unknown for this case. These average mass fractions are consistent with the conclusion of Worthey et al. (1996) that most or all elliptical galaxies have a metal-poor fraction $`5`$ per cent, with a few exceptions depending on the age of the metal-poor population; GMP 3565, NGC 4867, and NGC 4873 are the most extreme cases. We have not fully explored alternative explanations in this short paper, but GMP 3565 is the most metal-poor of the galaxies, and thus may have a stronger metal-poor tail if it has the same abundance distribution as other galaxies shifted to lower mean abundance. The other two galaxies are candidates to violate the 5 per cent rule, but explanations such as UV-upturn populations, multiple-age populations, and blue stragglers have yet to be explored.
After the appropriate fraction of hot-star light has been subtracted from the observed spectrum, Lick/IDS line strengths are recomputed (including velocity dispersion, emission, and systemic corrections). The results are shown in Figure 3 as vectors pointing from the observed line strengths to the corrected line strengths. As is clear from this figure, the change in the line strengths point along vectors of *nearly constant age* with an increase in metallicity $`[\mathrm{Z}/\mathrm{H}]_{\mathrm{SSP}}`$ and a slight increase in $`[\mathrm{E}/\mathrm{Fe}]_{\mathrm{SSP}}`$ being the major effects. This is borne out by the inferred stellar population parameters in Table 3 and the fractional changes in Table 4: only NGC 4867 changes its age by more than $`1\sigma `$. However, NGC 4867 is extrapolated far off of the W94 grids after subtraction in all three cases, and so the inferred ages and metallicities are suspect.
We therefore find that *the presence of even moderate amounts of light from hot stars does not significantly alter the inferred ages of early-type galaxies in the presence of intermediate-aged populations,* provided that the hot-star light comes from an old, metal-poor population containing blue horizontal branch stars. We however do not claim that the ages of early-type galaxies are *unaffected* by hot stars. In fact, we reproduce in Fig. 3 the result of Maraston & Thomas (2000) that the *oldest* stellar populations ($`t>10`$ Gyr) look younger by 2โ5 Gyr when ancient, metal-poor populations are superimposed on ancient, metal-rich populations. Rather, we suggest that when an intermediate-aged population is present, i.e., when $`\mathrm{H}\beta 1.5`$ ร
, the influence of hot stars on the inferred age is nearly negligible (see also Thomas et al., 2005). The difference between the effect of hot-star light on intermediate-aged and old populations is mostly due to the change in slope of the W94 grids at old ages and high metallicities in the $`\mathrm{H}\beta `$$`[\mathrm{MgFe}]`$ diagram and to the high $`\mathrm{H}\beta `$ strengths of the observed galaxies due, presumably, to intermediate-aged stellar populations.
The increased inferred metallicity but nearly constant age in the corrected population can be understood in the context of the Appendix of Trager et al. (2000b), which demonstrated that stellar populations add as vectors in the $`\mathrm{H}\beta `$โmetal-line spaces (such as $`\mathrm{H}\beta `$$`[\mathrm{MgFe}]`$). In effect, the โtwoโ populations in the Coma early-type galaxiesโthe metal-rich, intermediate-aged population with high $`\mathrm{H}\beta `$ and high $`[\mathrm{MgFe}]`$, and the metal-poor, old population with high $`\mathrm{H}\beta `$ (from the hot blue horizontal branch stars) and low $`[\mathrm{MgFe}]`$โadd to produce observed populations with high $`\mathrm{H}\beta `$ and moderate $`[\mathrm{MgFe}]`$ strengths. The strengthening of the $`\mathrm{H}\beta `$ line from the hot stars is compensated by a dilution of the continuum around the metal lines (Fig. 2), which weakens their measured equivalent widths. These weakened metal lines combined with the stronger $`\mathrm{H}\beta `$ line serve to preserve the inferred age of the galaxy while lowering the observed metallicity.
## 5 Conclusions
We have examined the effect of hot (horizontal branch) stars on the inferred stellar population parameters of early-type galaxies using observations of twelve $`\mathrm{H}\beta `$-strong early-type galaxies in the Coma Cluster and spectra drawn from stellar population models as well as an observed spectrum of a metal-poor globular cluster with a purely blue horizontal branch. If the hot-star light comes from ancient, metal-poor populations (Rose, 1985, 1994; Lee et al., 2000; Maraston & Thomas, 2000; Caldwell et al., 2003) typically contributing $`10`$ per cent of the light at 4000 ร
(as detected in the Ca ii index), the ages of these galaxies are not significantly affected by the correction for this hot-star light. For the oldest, most metal-rich galaxies, this correction can be significant, as shown by Maraston & Thomas (2000) and Fig. 3. But this correction is insignificant for the intermediate-aged populations found in these early-type Coma galaxies, and likely also for the field and group galaxies studied by Trager et al. (2000b), many of which have similarly high $`\mathrm{H}\beta `$ line strengths. We suggest therefore that the claim that old, metal-poor stars can โexplain awayโ the strong $`\mathrm{H}\beta `$ lines in these early-type galaxies (e.g., Maraston & Thomas, 2000) is overstated.
The presence of blue straggler stars in these galaxies is still a possibility to explain the enhanced Balmer-line strengths of early-type galaxies. However, as discussed by Rose (1985) and Trager et al. (2000a), populations of blue straggler stars are subject to the same constraints as other hot-star populations, as blue straggler stars typically have spectral types around mid-A. This is the same colour as the BHB stars that dominate the Balmer-line strengths of the old, metal-poor populations we have considered here. Using the same arguments we have already presented, therefore, we suggest that blue straggler stars are unlikely to affect the inferred ages of $`\mathrm{H}\beta `$-strong galaxies. We conclude that *intermediate-aged populations are still required* to explain the strong $`\mathrm{H}\beta `$ lines in early-type galaxies.
## Acknowledgments
The authors wish to recognise and acknowledge the very significant cultural role and reverence that the summit of Mauna Kea has always had within the indigenous Hawaiian community. We are most fortunate to have had the opportunity to conduct observations from this mountain.
It is a pleasure to thank D. Kelson, C. Maraston, J. Rose, and D. Thomas for many stimulating discussions, R. Schiavon for making his spectrum of NGC 6254 available in advance of publication, L. MacArthur for a careful reading of an early version of the manuscript, and an anonymous referee for suggestions that improved the clarity of the presentation. Support for this work was provided by NASA through Hubble Fellowship grant HF-01125.01-99A to SCT awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5-26555; by a Carnegie Starr Fellowship to SCT; by NSF grants AST-0307487 and AST-0346347 to GW; by NSF grants AST-9529098 and AST-0071198 to SMF; and by NASA contract NAS5-1661 to the WF/PC-I IDT. |
warning/0506/cond-mat0506410.html | ar5iv | text | # Electronic structure and X-ray magnetic circular dichroism of CrO2
## I Introduction
Chromium dioxide is the only stochiometric binary oxide that is a ferromagnetic metal and also is the simplest and best studied half metal Coey . In recent years, magneto electronic devices such as spin valve field sensors and magnetic random access memories have emerged, where both charge and spin of electrons are exploited using spin-polarized currents and spin-dependent conduction Prinz ; Daughton . The performance of such magneto-electronic devices depends critically on the substantial spin polarization of the ferromagnetic components. CrO<sub>2</sub> attracts specific interest being the only compound reported so far experimentally to possess 100% spin polarization. The half-metallic band structure of CrO<sub>2</sub> was confirmed by several experimental techniques such as Andreev reflection, super conducting tunelling, photoemission, point-contact magnetoresistance, x-ray absorption, resonant scattering and Raman spectroscopyJi ; Parker ; Schmitt ; Dedkov ; Coey1 ; Coey2 ; Huang1 ; Gupta ; Kurmaev ; Iliev .
Half-metallic ferromagnetic materials appear as potential candidates and a lot of work is under progress to synthesis magnetic oxides such as CrO<sub>2</sub>, Fe<sub>3</sub>O<sub>4</sub> and Sr<sub>2</sub>FeMoO<sub>6</sub> which are found to posses high Curie temperature (T<sub>c</sub>). Recently Sr<sub>2</sub>CrReO<sub>6</sub> a Cr based double perovskite is found to have a largest T<sub>c</sub> of about 635 K which is the highest among the double perovskite familykato2002 . These compounds serve as perspective materials in the field of spintronicsHaghiri . Numerous theoretical works are available for CrO<sub>2</sub> explaining the electronic structure, bonding, optical and magneto-optical properties and magneto-crystalline anisotropy(MCA) etcSchwarz ; Sorantin ; Uspenskii ; Lewis ; Korotin ; Mazin ; Guo ; Oppeneer ; komelj ; Toropova . Because of the uniaxial crystal structure, CrO<sub>2</sub> is expected to have a large magnetic anisotropy which makes it the favoured material for magneto-optical recording. Recently it has been shown that the low-temperature experimental data were reproduced well by LSDA itself without taking into account the Hubbard $`U`$ correction confirming that the ordered phase of CrO<sub>2</sub> is weakly correlatedToropova .
The circular dichroism-type spectroscopy became a powerful tool in the study of the electronic structure of magnetic materialsBagus ; Ebert . It has been recently demonstrated by Weller et alWeller that x-ray magnetic circular dichroism (XMCD) is also a suitable technique to probe MCA at an atomic scale, via the determination of anisotropy of the orbital magnetic moment on a specific shell and site. The x-ray absorption spectroscopy (XAS) using polarized radiation probes element specific magnetic properties of compounds via the XMCD by applying the sum rules to the experimental spectraThole ; Carra ; Laan . The application of these sum rules to itinerant systems, in particular to low symmetry systems, is debated since the sum rules are derived from atomic theoryCarra ; Chen ; Wu . Recently angle resolved XMCD technique has been applied by Georing et alGeoring ; Georing1 ; Gold for a wide range of temperature to investigate the anisotropies of orbital moments $`l_z`$ and magnetic dipole term $`t_z`$ of Cr atom in epitaxial CrO<sub>2</sub> films on TiO<sub>2</sub> substrate. From their XMCD studies they found large anisotropies of $`l_z`$ and $`t_z`$, where the latter was derived from the semi emprical van der Laanโs method of moment analysisLaan2 .
Theoretical understanding of XMCD for magnetic material is not an easy task, and several ab initio calculations have attempted to compute XMCD of transition metals and rare-earth compoundsWu ; Ebert1 ; Alouani ; Brouder ; Guo2 ; Ankudinov ; Galanakis1 ; Galanakis2 . The $`L_2`$ and $`L_3`$ edges involving electronic excitations of the $`2p`$-core electrons towards $`d`$-conduction states have attracted much attention due to the dependence of the XMCD spectra on the exchange-splitting and the spin-orbit coupling of both initial core and final conduction states. Brouder and co-workersBrouder , GuoGuo2 , and Ankudinov and RehrAnkudinov used multiple-scattering theory to study XMCD but their method, although successful, has been applied to systems with few atoms per unit cell, as their formalism is computationally involved. Finally, atomic calculations, using crystal-field symmetry, were widely applied to fit the experimental $`M_4,_5`$ edges of the rare earths and actinide compounds and the $`L_2,_3`$ edges of early transition metals. Because of the large number of parameters to fit, it is difficult to apply this formalism to delocalized $`3d`$ statesvan . Though there are lot of theoretical studies available explaining the electronic and magnetic properties of CrO<sub>2</sub> theoretical studies on the XMCD spectra were not available. Parallel to our study, Baadji and coworkersbaadji are investigating the effect of Hubbard interaction on the magnetic properties and XMCD specta of CrO<sub>2</sub> using the linear Augmented plane wave method.
In the present study, efforts were taken to study theoretically the L<sub>2,3</sub> edge of Cr and K-edge of oxygen of CrO<sub>2</sub> using relativistic full-potential linear muffin-tin orbital method (FP-LMTO) method. In Section II we briefly present the theoretical method used to calculate the electronic structure, orbital magnetic moment and XMCD spectra. The density of states and magnetic moments are discussed in Section III. In Section IV we calculate the XAS and XMCD spectra of CrO<sub>2</sub> and compare with the experiment. The conclusion of the results are presented in Section V.
## II Theoretical Details
In the present work electronic structure calculations were performed using the all-electron full-potential linear muffin-tin orbital (FP-LMTO) methodwills including the spin-orbit interaction directly in the Hamiltonian. The exchange correlation potential is parametrized using the generalized gradient approximation (GGA)pbe . In this method, space is divided into non-overlapping muffin-tin spheres surrounding the atoms, and an interstitial region. Most importantly, this method assumes no shape approximation of the potential, wave functions, or charge density. The spherical-harmonic expansion of the potential was performed up to $`l_{max}=6`$, and we used a double basis so that each orbital is described using two different kinetic energies in the interstitial region. The basis set consisted of the Cr ($`4s`$ $`4p`$ $`3d`$), O ($`2s`$ $`2p`$) LMTOs. We performed our calculations using the experimentally determined structure and atomic positions, i.e., the rutile structure with space group symmetry $`P4_2/mmm`$ and with cell parameters $`a=b=4.419`$ ร
, $`c=2.912`$ ร
Porta . The radii of the muffin-tin spheres used for Cr and O were 2.0 and 1.5 Bohr units respectively. To find the easy magnetization axis we calculated total energy only for as well as magnetization axis. The integration in reciprocal space was performed using the tetrahedron methodjepsen72 and 1100 irreducible k points in the Brillouin zone (BZ), for the magnetization axis, whereas for magnetization axis 1200 k points were used using the same BZ grid as for the previous quantization axis. To avoid numerical errors one has to use the same k grid for both spin quantization axis and let the symmetry of the crystal in presence of the spin-orbit decide for the number of of irreducible k points. The theoretical XAS and XMCD spectra was calculated using the method described elsewhereAlouani . This method was found to be successful in reproducing the experimental XAS and XMCD spectra of serval transition metal compoundsAlouani ; Galanakis1 ; Galanakis2 .
## III Density of states and magnetic moments
The electronic structure of CrO<sub>2</sub> has been extensively reported in the literatureSchwarz ; Sorantin ; Uspenskii ; Lewis ; Korotin ; Mazin ; Guo ; Oppeneer ; Toropova so in this section we briefly discuss the density of states and compare with the ealier works.
Figure 1 shows the density of states (spin up and down) of Cr-$`3d`$ and O-$`2p`$ states. For the majority spin, the Fermi level lies near a local minimum of the Cr 3d-$`t_{2g}`$ band with the DOS at the Fermi level. For the minority spin, the Fermi level falls in a gap of 1.34 eV which is in agreement with the earlier resultsSchwarz ; Lewis ; Korotin ; Mazin ; Guo . The calculated exchange splitting between the majority and the minority spin main peaks of the Cr 3d-$`t_{2g}`$ band is found to be 2.3 eV. A similar splitting of 2.5 eV has been found in FPLAPW calculationsMazin , while a splitting of 2.3 eV has been found in a recent FPLMTO calulationsGuo . All the exchange splitting calculated within the LSDA or GGA are too small when compared with the measured large splitting of about 5 eV between the main peaks in the occupied and the empty Cr-$`3d`$ DOS, from experimental photoemission studiesTsujioka . It is interesting to note that half-metallic gap persists even with spin-orbit coupling. In general the spin-orbit coupling evantually destroys the half-metallic nature in many Huesler alloys, dilute magnetic semiconductors, zinc blende type transition metal pnictides and double perovskitesMavropoulos ; apl .
The calculated spin and orbital magnetic moments along the axis of magnetisation are found to be 1.99 $`\mu _B`$ and -0.045 $`\mu _B`$/atom for Cr and -0.08 $`\mu _B`$ and -0.0017 $`\mu _B`$/atom for O respectively. The calculated spin magnetic moments of Cr and O are slightly higher when compared to the recent FP-LMTO calculations of Jeng and GuoGuo within the LSDA. In the present calculation we have used GGA to the exchange correlation which will tend to increase the magnitude of the spin moment. When comparing the orbital magnetic moments of Cr and O they compares well with the recent FLAPW calulationskomelj . From the above results it can be seen that the the orbital magnetic moment of Cr is almost quenched and is also antiparallel to the spin moment, which is consistent with Hundโs rule for 3$`d`$ shells which are less than half-filled. A similar quenched orbital magnetic moment for Cr has been recently reported in the double perovskite Sr<sub>2</sub>CrReO<sub>6</sub>apl . The orbital magnetic moment of the O atom is parallel to the spin moment because the O-2$`p`$ shell is more than half-filled. The calculated spin moment of oxygen is antiparallel to the spin moment of Cr and hence the Cr and O are coupled antiferromagnetically. The calculated orbital magnetic moments of Cr and O are found to be -0.045 and -0.0017 $`\mu _B`$/atom respectively which are slightly lower when compared with the experimental XMCD measurments of Huang et al in which they obtained a value of -0.06$`\pm `$0.02 $`\mu _B`$/Cr and -0.003$`\pm `$0.001 $`\mu _B`$/O respectively in the axis of magnetizationHuang2 . While Georing et alGeoring in their XMCD measurments on CrO<sub>2</sub> films found a relatively small contribution to the orbital magnetic moment of Cr in the axis of magnetization which is found to be -0.02 $`\mu _B`$. When comparing the true spin moment of Cr atom Georing et alGeoring obtained a value of 1.2 $`\mu _B`$. In order to explain the magnetic moment of 2 $`\mu _B`$ per unit cell which was obtained from the superconducting quantum intereference device (SQUID) measurments the above authors assumed a very large spin moment of about 0.4 $`\mu _B`$ per O which they try to explain interms of hybridization between the Cr and oxygen. Recently the same authors applied spin correction factor to their spectra from which they were able to get a spin magnetic moment of 2.4 $`\mu _B`$ which is slightly higher when compared to the present theoretical and other experimental worksGeoring2 . For the oxygen atom the spin moments are once again too low when compared to the value estimated by Georing et alGeoring . However it was recently shown by Komelj et alkomelj that calculations with Hubbard $`U`$ = 3 eV yield a spin moment of 0.1 $`\mu _B`$ which is a factor of 4 times smaller than the value obtained by Georing et alGeoring . In order to find the easy axis of magnetization we have done self consistent calculations only along the and magnetization axis. Our calculations find axis to be the easy axis of magnetization which is in agreement with the experimentsSpinu ; Yang ; Li ; Miao ; Gold .
## IV X-Ray absorbtion and magnetic circular dichroism
In this section we calculate, analyse and compare the XMCD spectra with experiment. At the core level edge XMCD is not only element specific but also orbital specific. For 3$`d`$ transition metals, the electronic states can be probed by the K, L<sub>2,3</sub>, M$`{}_{2}{}^{},_{3}^{}`$ x-ray absorption and emission spectra. The dichroism at the L<sub>2</sub> and L<sub>3</sub> edges is influenced by the spin-orbit coupling of the initial 2$`p`$ core states. The large SO splitting of the core levels gives rise to a very pronounced dichroism in comparison with the dichroism at the K-edge. In figure 2, we present the XAS and XMCD spectra for Cr atom in the axis of magnetization. We convoluted our theoretical spectra using a Lorentzian and Gaussian width of 0.5 eV. The Gaussian represents the experimental resolution while Lorentzian corresponds to the width of the core hole.
The calculated total absorbtion spectra corresponding of the Cr L<sub>2,3</sub> edge was shown in the upper part(figure 2). We have scaled our spectra in a way that the experimental and theoretical L<sub>3</sub> peaks in the absorption spectra have the same intensity. The energy difference between the L<sub>3</sub> and L<sub>2</sub> peaks is given by the spin-orbit splitting of $`p_{1/2}`$ and $`p_{3/2}`$ core states which is 8.65 eV. The calculated value agrees well with the experimental value of 8.2-8.6 eVGeoring2 . As far the absorption spectra is concerned the present theory represents well the experimental features but underestimates the intensity of the L<sub>2</sub> peak.
The intensity of the L<sub>3</sub> edge in the x-ray absorption spectra when compared to the L<sub>2</sub> peak is a bit high which is in agreement with the experiment. The calculated XMCD spectra along the axis of magnetization is shown in the lower part. The present caculations reproduce well the shape of the experimental spectra.
Using the XMCD sum rules we calculate the spin and orbital magnetic moment of Cr-3$`d`$ state which are found to be 1.82 and -0.044 $`\mu _B`$/atom respectively and these values of spin and orbital magnetic moments and are in agreement with the values obtained from the self-consistent calculations. A deviation of about 10% in the magnetic moments when comparing the values from sum rules and direct calculation may result from the spectral overlapGeoring3 . When comparing the spin and orbital magnetic moments obtained using the sum rules with the theoretical LSDA values of Huang et alHuang2 the spin moment of Cr is slight lower when compare to the theoretical value of 1.9$`\pm `$0.1 $`\mu _B`$ whereas the orbital magnetic is found to be -0.06 $`\pm `$0.02 $`\mu _B`$ which is in good agreement with the present calculations. However the experimental studies of Huang et alHuang2 could not provide quantitative information on the spin moment of Cr, because they cannot uniquely define which part of spectra belongs to the L<sub>3</sub> or L<sub>2</sub> edge. When comparing our magnetic moments obtianed from sum rules with that of the XMCD measurments of GeoringGeoring the spin magnetic moment are too low (1.2 $`\mu _B`$), whereas the spin magnetic moment(2.4 $`\mu _B`$) obtianed by the same authors after spin correction applied to the spectraGeoring2 is slightly higher when compared to the present calculations. The orbital magnetic moments are also slightly higher (-0.044 $`\mu _B`$) when compared to the experimental results of Georing in which they obtain a value of -0.02 $`\mu _B`$.
In figure 3, we have shown the oxygen K edge XAS and XMCD spectra along the axis of magnetization.
We used a similar broadening as that of the Cr L<sub>2,3</sub> for the oxygen K edge. The oxygen K edge spectrum mainly comes from excitation of the 1$`s`$ state and is less intense when compared to the L<sub>2,3</sub> edge of Cr. The exchange splitting of the initial 1$`s`$ core state is extrmely small value, therefore, only the exchange and spin-orbit splitting of the final 2$`p`$ state is responsible for the observed dichroism at the K-edge. However the present calculation reproduce fairly well the experimentally observed O-K edge and XMCD spectra.
In figure 4, the Cr L<sub>2,3</sub> XAS and XMCD spectra of CrO<sub>2</sub> along the axis of magnetization is also shown.
The core level splitting of 8.6 eV is found in this axis which is in agreement with experimental value of 8.2-8.6Georing2 . The theoretical XAS spectra agrees quite well with the experiments except the intensity of the L<sub>2</sub> peak is lower when compared to the experiment. The calculated XMCD spectra is shown in the lower part reproduce well the shape of the experimentally observed features. The calculated XMCD intensity is relatively less in the magnetization axis when compared to the axis of magnetization. For both the axes, the low intensity of the calculated L<sub>2</sub> edge makes the theoretical XMCD integrated L<sub>3</sub>/L<sub>2</sub> branching ratio to be much larger than the experimental one. This discrepancy may be due to the negligence of 3$`d`$ core-hole interaction. Using the XMCD sum rules the spin and orbital magnetic moments of Cr atom along the magnetization axis is found to be 1.844 $`\mu _B`$ and -0.004 $`\mu _B`$, respectively. The caluculated spin magnetic moment agrees well with the earlier calculationsGuo ; komelj . The orbital magnetic moment obtained from the sum rules in the axis is much lower when compares to the direct calculations. A similar discrepancy was observed in NiMnSb in which the orbital magnetic moments of Mn obtained from sum rules are much lower when compared to the direct calculationGalanakis1 .However the theoretical orbital magnetic moments are much lower when compared to the experimental values of -0.09 $`\mu _B`$Georing .
The oxygen orbital magnetic moment is almost quenched in the axis, which are similar to the earlier theoretical worksGuo ; komelj resulting in a weak intensity of the oxygen K edge in contrast to the experiments, which show the same angular dependency for the Cr-3$`d`$ orbital projections and the O-K edge XMCD.
## V Conclusions
In the present work we have carried out a detailed theoretical study on the magnetic properties of CrO<sub>2</sub>. The calculated half-metallic band structure of CrO<sub>2</sub> is in agreement with earlier studies. Our calculation confirms the c-axis to be the easy axis of magnetisation. The calculated XAS and XMCD spectra of Cr L<sub>2,3</sub> and oxygen K edge compares fairly well with the experimental results. The L<sub>3</sub>/L<sub>2</sub> branching ratio is higher when compared to the experiments which could be improved by including the 3$`d`$ core-hole interaction. In addition the spin and orbital magnetic moments obtianed from the sum rules are compared with the direct calculations.
###### Acknowledgements.
The authors acknowledge Dr. E. Goering for sharing experimental details. M.A. thanks O.K. Andersen for an invitation to the MPI during the completion of this work. |
warning/0506/astro-ph0506213.html | ar5iv | text | # The Formation Histories of Galaxy Clusters
## 1 Introduction
Due to their immense size, galaxy clusters can be easily identified with their dark matter halos, and their clustering and number counts can be reliably predicted by dark matter simulations. However in order to find or โweighโ galaxy clusters observationally, assumptions of dynamical equilibrium and less understood astrophysics usually come into play<sup>1</sup><sup>1</sup>1Weak gravitational lensing measures mass without equilibrium assumptions, but only in projection (see e.g. Refs. lens ).. Since halos form via the accretion and mergers of smaller units, and for clusters this process is occurring to the present time, such assumptions need to be specified precisely. For example, how well a galaxy cluster is described by equilibrium properties depends upon whether it has recently undergone a major merger (where โmajorโ means disrupting equilibrium) and upon the specific methods used for its detection and mass measurement. Thus even if one is only interested in galaxy clusters as โtest particlesโ or peaks in the density field, the history of cluster growth is important in order to make contact with observations. Cluster growth histories are also important for understanding other cluster properties. For example, several observational phenomena (see below) are associated with clusters which have recently undergone mergers. A related question is which cluster properties have the least sensitivity to cluster assembly histories.
In this paper we study the assembly history and degree of virialization of high-mass halos in large N-body simulations. A cluster assembly history can be characterized either by events with specific occurrence times, such as mergers or large mass changes, or by properties of its entire history, e.g. a parameterization of the mass as a function of time. We calculate several such quantities for a statistically significant sample of halos. We consider several popular formation time definitions and cluster history parameterizations and statistics. We also calculate the fraction of โrecentlyโ merged galaxy clusters as a function of redshift, back to $`z1`$.
The definitions of both โmergerโ and โrecentlyโ depend on the cluster property of interest and we explore several choices. We also find the fraction of galaxy clusters which have had a recent large mass gain (including accretion) for several choices of interval and two final/initial mass ratios. Merger histories can be reliably extracted from N-body simulations, and as such these recently merged fractions are implicit in earlier work. However, it can be difficult to obtain specific numbers from the literature for $`\mathrm{\Lambda }`$CDM models, especially if one has a particular relaxation time in mind. In part this is because previous studies of different quantities at different times have been published over several years, often using different cosmological models. Here we compute several of these quantities for a much larger sample than used in earlier work, for $`\mathrm{\Lambda }`$CDM cosmologies, and present them in a homogeneous manner in the hope that this will be a useful reference for the community.
The outline of the paper is as follows. Section 2 is a review, including pointers to earlier work on formation times and examples of observed merger phenomena. Section 3 describes the simulations and methods. A reader interested primary in the results can skip directly to Section 4, which has comparisons and distributions of some cluster formation properties, and the fractions of clusters which have recently merged or had a large mass increase, as a function of time. Three and four-body major mergers are also studied. Finally, Section 5 presents our conclusions.
## 2 Background
The growth of structure by mergers and accretion is key to the hierarchical paradigm of structure formation, and thus mergers and mass gains have been studied intensively. Previous work on cluster formation histories includes Refs. TorBouWhi97 ; Tor98 ; Col99 ; GotKlyKra01 ; Zha ; RowThoKay04 ; TKGK ; Bus03 . While cluster assembly is a complex and ongoing process in hierarchical models, it is often useful to have some measure of when the cluster โformedโ. Refs. GotKlyKra01 ; RowThoKay04 ; TKGK , each with 10-20 clusters, considered formation time definitions including the redshift, $`z_{\mathrm{jump}}`$, of the most recent large $`\mathrm{\Delta }M`$ over a short time. Both RowThoKay04 ; TKGK also found a characteristic formation time $`z_f`$ associated with the entire cluster growth curve using the parameterization of Ref. Wec02 . This parameterization works extremely well for galaxy sized halos and correlates with other properties such as concentration. For galaxy clusters, Ref. TKGK introduced a generalization to help better match the more recently active formation histories of galaxy clusters. Ref. Zha had a sample similar to ours and found a โturning pointโ time where halos went from a quickly growing phase to a more slowly growing phase, this turning point was correlated with concentration. There are other formation times considered in the literature. For instance Ref. vdB02 found the average mass accretion history of halos generated by extended Press-Schechter. Ref. RowThoKay04 also looked at the amount of mass gain coming from large $`\mathrm{\Delta }M`$ โjumpsโ as fraction of cluster mass. We consider these properties and their distribution for our large sample of over 500 clusters for a $`\mathrm{\Lambda }`$CDM $`\sigma _8=0.8`$ model in ยง4.
Our second set of results is the frequency of recent mergers and recent large mass gains for some fixed lookback time and definition of โrecent.โ Analytic estimates of merger rates, in many cases combined with simulations, include those by Refs. analytic ; Tor98 ; GotKlyKra01 ; CohBagWhi01 (see also the extension reported in Nic02 ). The previous work closest to the merger counts considered here is by Ref. GotKlyKra01 who found the major merger distribution as a function of redshift for a system which had 11 clusters at $`z=0`$ (with a merger defined as an increase in halo mass such that $`M_f/M_i>1.33`$) and that of Ref. CohBagWhi01 which also studied the recently merged population, but for a smaller number of clusters and not in as much detail. It should be noted that merger rates and the number of clusters which have recently undergone a merger are slightly different, it is the latter we study here.
Mergers can alter X-ray temperatures, luminosities, cluster galaxy light, cluster galaxy velocities and Sunyaev Zelโdovich (SZ; Ref. SZ ) โfluxโ<sup>2</sup><sup>2</sup>2Several simulations with different heating and feedback prescriptions have found small merger induced scatter in the integrated SZ flux relation (Refs. WhiHerSpr02 ; Mot05 ; Koc05 although there are examples where a $`50\%`$ increase in total SZ flux for one sound crossing has been seen Sar05 ). Significant scatter in the mass-flux relation instead appears to be dominated by that due to line of sight projection WhiHerSpr02 .. Observational consequences of merging clusters of galaxies have been studied analytically and with numerical simulations, e.g. Refs. numer ; RicSar01 ; RowThoKay04 . In these, the time scale for major disturbances of clusters due to mergers seems only weakly dependent upon the impact parameter (the magnitude of the disturbance has a stronger dependence, e.g. see Refs RicSar01 ). We do not consider impact parameters below. The errors induced by ignoring merger effects (on X-ray luminosities and temperatures) for cosmological parameter estimates has been considered in RanSarRic02 <sup>3</sup><sup>3</sup>3They used an analytic model SomKol99 to get a merger history, in principle our calculations here could also be used for this purpose.. More generally, unaccounted for mergers can disguise cluster masses and alter survey selection functions.
Observationally, some signatures of mergers can be used to flag unrelaxed clusters. For instance substructure, the presence of multiple peaks in the cluster surface density on scales larger than the constituent galaxies, has been studied extensively. Several different observational methods have been used to quantify the amount of substructure in X-ray clusters (see the review Buo02 and references therein). The fraction of clusters exhibiting substructure has been found to be significant in several surveys. For example, Ref. JonFor99 visually found substructure in 41% of 208 Einstein IPC images, Ref. Moh95 measured the emission weighted centroid variation for 65 Einstein clusters to get a substructure fraction of 61% (see also moresub for related work) and Ref. Sch01 found, using three different indicators, substructure in $`52\pm 7\%`$ of 470 clusters in the ROSAT all-sky survey. Substructure has also been searched for in other wave-bands. For example, substructure was detected in over 80% of 25 low richness 2dFGRS clusters (Bur04 ; see their references for other studies). Following Ref. TsaBuo96 , substructure has been considered as a means of constraining cosmology, e.g. in Jel05 using 40 Chandra clusters over a range of redshifts (see also struccos for relevant simulations). A meaningful comparison of our calculated number of โrecently mergedโ clusters to the number of observed galaxy clusters with substructure depends on relaxation times. These times depend upon the substructure measurement in question and are still not well estimated. For instance, relatively small sub-clumps may take a different amount of time to disappear depending on their density, see Buo02 for further discussion.
Radio halos and radio relics have been proposed as another merger indicator (see the reviews Fer04 and Sar04b ). Clusters hosting radio halos and relics tend to show other signs of merger activity, however not all clusters with other indicators of recent mergers show these radio signals. In a complete sample GioFer02 at the NRAO VLA Sky Survey surface brightness limit, only 5% of clusters had a radio halo source and 6% had a peripheral relic source, however in clusters with X-ray luminosity larger than $`10^{45}`$ erg s<sup>-1</sup>, 35% had radio halos or relics. In the ROSAT sample considered above by Ref. Sch01 , 53 clusters of the 470 had radio halos or relics, most of these were high luminosity ($`L>4.0\times 10^{44}`$ erg s<sup>-1</sup>). It has been proposed that more radio halos and relics will be found in lower mass clusters as sensitivities improve (e.g. Ref. Sar05 ). The relaxation time for radio halos and relics is not known.
The absence of mergers (i.e. a relaxed cluster) has been suggested as a requirement for cool core clusters (e.g. Refs. cool ). The fraction of cool core clusters has been measured for various samples, e.g. Ref.Bau05 found with Chandra that 55% of 38 clusters in the ROSAT Brightest Cluster Sample (with $`z0.150.4`$) show mild cooling (associated with a cooling time less that $`10`$Gyr) and 34 % show signs of strong cooling (a cooling time less than $`2`$Gyr). There are some simulations that suggest however that mergers may help form cool galaxy cluster cores (e.g. see Mot04 ), also 5/22 of the ROSAT cooling flow clusters have significant substructure Sch01 .
There are many more phenomena associated with mergers which are even more complex (see Sarazinโs list of questions about merging clusters of galaxies Sar04 ). The examples of merger indicators reported above have had their occurrence rates measured in some of the largest cluster samples available.
There are some differences between mergers and any form of mass gain<sup>4</sup><sup>4</sup>4A modified version of the PS formalism that differentiates between (instantaneous) major mergers and accretion has been developed in Refs. mer-ac .. Clusters which have gained mass primarily via accretion might be expected to be more relaxed than those who reach the same mass via a major merger. For smaller halos those with a relatively large mass gain over a short time are more clustered than halos of the same mass but without this mass gain ScaTha03 . The subset of recently merged halos, in contrast, does not appear to be so biased Per03 . The authors of Ref. ScaTha03 conjecture that objects in denser regions tend to have more nearby material to accumulate, resulting in a bias, while the same is not true of recent mergers.
## 3 Methods
To investigate these questions we used two N-body simulations run with a TreePM code Whi02 . Each simulation evolved $`512^3`$ particles in a periodic cubical box, $`300h^1`$Mpc on a side, for a particle mass of $`1.67\times 10^{10}h^1M_{}`$. The models had $`h=0.7`$, $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ and $`\sigma _8=0.8`$, 1.0. The two normalizations span the observationally preferred range. Outputs were dumped at equal intervals<sup>5</sup><sup>5</sup>5This is the same time interval considered by Ref. RowThoKay04 and shorter than that considered by Refs. GotKlyKra01 ; TKGK , though the latter groups still had their time spacings well within the lower limits for a merging time scale. of conformal time, $`\delta \tau =100h^1`$Mpc (comoving), starting at $`z=2`$. Other parameters and details of the simulations can be found in YanWhi04 .
For each output groups are identified with a friends of friends algorithm (FoF; Dav85 ) using a linking length $`b=0.15`$ in units of the mean inter-particle spacing. These groups correspond roughly to all particles above a density of $`3/(2\pi b^3)140`$ times the background density. We did not consider other group finders, a comparison between HOP EisHut98 and FoF merger fractions for a smaller sample was considered in CohBagWhi01 . In that case there were both fewer recently merged clusters and fewer clusters altogether for the FoF case. See also Ref. ScaTha03 for a comparison of HOP and FoF mergers and accretion for smaller mass halos. They estimate what minimum mass fraction in accretion is robust for different group fractions and find that a 20% mass increase (for a shorter length of time and smaller halos than we consider) is outside the โnoiseโ region for the FoF vs. HOP group finders.
For each group and output time we compute a number of properties, including the mass, velocity dispersion and potential energy of the group. All properties are computed using only group members. Masses henceforth refer to friends of friends masses $`M_{\mathrm{fof}}`$ unless otherwise stated. To get $`M_{200}`$ or $`M_{\mathrm{vir}}`$ conversion is needed Whi01 โ for a fitting formula see HuKra03 . From the simulation we find $`M=1.2M_{200}`$ at $`z=0`$, and $`M=0.98M_{200}`$ at $`z1`$. We consider halos to be clusters if $`M10^{14}h^1M_{}`$ (i.e. 5971 or more particles); at $`z=0`$ there are 909 clusters in the $`\sigma _8=1.0`$ box and 574 in the $`\sigma _8=0.8`$ box.
For the formation time studies, most of the definitions are straightforwardly derived from the cluster histories. We track the properties of all of the progenitor halos with more than 64 particles for all of our clusters. From these histories it is straightforward to find the mass of the largest progenitor as a function of time. To find mergers, the two largest predecessors are found for each cluster at an earlier time. The final cluster is said to have had a (major) merger between the earlier time and the later time if the ratio of particle numbers going from both of the predecessors into the final cluster is above some minimum<sup>6</sup><sup>6</sup>6This was in part to try to minimize the effect of particles which were not actually bound being taken to be part of a clusterโs predecessor.. In the case where the time interval of interest does not coincide with a dump time for the data, the two closest times bracketing the time of interest are taken, and the merger fraction is linearly interpolated between these (starting from the time closest to the time of interest). One relaxation time we consider is the crossing time, in this case we start at the earlier time and interpolate between final times.
We checked our definitions of mergers in two ways. Instead of using the mass of only progenitor particles which actually go into the final halo in the mass ratio, one can use the full mass of the progenitor halos, this changes the merged fraction by less than 5% for the longer time intervals and almost always less than 10% for the short time intervals. We prefer our definition as it might help counteract the inclusion of many unbound particles in halos due to our FoF group finder. Secondly, we identify mergers by looking at the initial and final halos and comparing mass ratios of the initial halos, without tracking their intermediate behavior until they merged to the final halo. A concern was that two predecessors with an initial mass ratio of e.g. 1:5 might first become a mass ratio of 1:10, say, if the largest progenitor grew significantly in mass relative to the second largest, before merging. Looking at intermediate times, the mass ratio of predecessors is differs for the two definitions, but mostly for almost equal mass predecessors. Once the predecessor mass ratio drops to $``$ 1:3 the scatter between the two ways of defining predecessor mass ratio gets quite small. Using the more accurate definition based on intermediate time steps produces a slightly larger number of recently merged clusters, again less than 5%. That is, our approximation in some cases underestimates the number of mergers by a small amount.
We vary the mass ratios and time intervals, and consider intervals equally spaced in real time, in light crossings (scaling as $`a`$) and in terms of โcrossing timesโ (scaling as $`(G\overline{\rho })^{1/2}a^{3/2}`$). We call the lookback time the time at which the cluster is observed, and the relaxation time the interval prior to the lookback time within which the merger or mass jump (โlarge $`\mathrm{\Delta }M`$โ) has or hasnโt occurred. Intervals of 100, 300 and $`600ah^1`$Mpc light crossings correspond to approximately 1, 3 and $`6\times a/2`$Gyr. We call these $`\delta \tau =100,300,600h^1`$Mpc in the following. Although the simulations go back to $`z2`$, we only plot back to lookback times where the statistics have any constraining power.
We also consider three- and four-body merger fractions, since such systems have been observed (see, e.g. Refs. threebody ). We define three-body mergers as the case when the second and third largest predecessors each contribute to the final halo at least 20% of the mass contributed by the largest predecessor. Similarly, four-body mergers require the fourth largest predecessor to also contribute at least 20% of the mass contributed by the largest predecessor.
## 4 Results
### 4.1 Formation times and properties
A classic paradigm for cluster assembly histories is the spherical top-hat collapse model. In this model a uniform, overdense region โbreaks awayโ from the expansion of the universe, evolving as a self-contained positive curvature universe. This toy model is often used to motivate the threshold density for virialization or to estimate formation times LiddleText . The evolution of such perturbations does not closely resemble the rich structure seen in the formation of massive halos in N-body simulations, as noted by many authors referenced above, however it is interesting to see how it performs nonetheless. We show in Fig. 1 the evolution of the peak circular velocity ($`Mv_c^3`$) of the most massive progenitor for 6 clusters, chosen at random from amongst the most massive ($`M>6\times 10^{14}h^1M_{}`$) in our $`\sigma _8=0.8`$ simulation. Since the mass in the spherical top-hat model is difficult to compare to simulations (it is constant and has uniform overdensity), we use peak circular velocity. The peak of the circular velocity curve, $`v_c^2GM(<r)/r`$, is computed for each progenitor using the minimum of the cluster potential as a center for defining $`r`$. The dashed line shows the evolution predicted using the spherical top-hat collapse model for a cluster which virializes at the present. We see that the spherical infall description has a clusterโs peak circular velocity growing most quickly right before virialization, while clusters forming in a cosmological simulation have more of a steady growth over time<sup>7</sup><sup>7</sup>7For the spherical top-hat collapse model, the cluster mass is constant and the density/radius change with time. In physical coordinates the perturbation first expands (falling $`v_c`$) and then collapses (rising $`v_c`$).. We have also marked on the plot four of the formation times defined and discussed below, $`z_{\mathrm{jump}}`$, $`z_{1/2}`$, $`z_f`$ and $`z_{\mathrm{tp}}`$.
A different view of the formation of clusters is given in Fig. 2 which focuses on the same 6 clusters as Fig. 1. In the top panels we show the mass of the largest progenitor as a function of time. In the middle panel we show the ratio of the 1D velocity dispersion to the peak circular velocity โ this has been used as a proxy for the degree of virialization of the clusters KneMul99 . In the lower panel we show the relation between the total kinetic and potential energy explicitly. Both energies have been computed using only the particles in the FoF group, but no proxy has been used, unlike the middle panel.
Even focusing only on the most massive progenitor, rather than the full merger tree, the upper panels show the rich phenomenology of cluster formation. There are several obvious mass jumps interspersed between periods of relatively smooth accretion. The lower panels show that the cluster is disturbed by mergers and large mass jumps with strong departures from the expected virial relation $`2\mathrm{K}\mathrm{E}\mathrm{PE}`$. The departures last for several of our output times. The โstandardโ vacuum virial relation $`2\mathrm{K}\mathrm{E}\mathrm{PE}`$ is not satisfied by these clusters even in their quiescent phase. The ratio $`2\mathrm{K}\mathrm{E}/\mathrm{PE}`$ is typically larger than unity due to the steady accretion of material onto the cluster<sup>8</sup><sup>8</sup>8This can be thought of as a surface pressure term which modifies the virial relation, boosting the kinetic energy.ColLac96 ; KneMul99 and there is an overall decline in $`2\mathrm{K}\mathrm{E}/\mathrm{PE}`$ with time. Averaging over the entire cluster sample we find for objects of fixed mass the KE/PE ratio decreases with decreasing redshift. For objects at fixed redshift the ratio increases with increasing mass. For objects at fixed number density the ratio decreases with decreasing redshift. For fixed mass cut or number density the decline is steeper below $`z=1`$ than above $`z=1`$ and within errors the decline below $`z=1`$ has the same shape for all samples. We hypothesize that this decline in KE/PE as due to the cessation of structure formation due to $`\mathrm{\Lambda }`$ domination below $`z1`$ (plus a small shift in the mass of objects currently undergoing rapid accretion). We will see a similar drop in merger activity below $`z1`$ in the next section.
We attempted to further quantify the relaxation of clusters to the background virial relation as a function of time since last major disturbance (c.f. Ref. RowThoKay04 ), but we found the scatter was too large for us to draw robust conclusions.
Even though cluster formation is a complex process it is sometimes useful to attempt to describe it with a single โformation timeโ. Any such compression of information must be imperfect, and to some extent arbitrary. Depending on the phenomenon of interest different times may be more or less appropriate. For this reason several definitions in the literature exist for the โformation timeโ of a cluster. We calculate several of these and their distributions for the $`\sigma _8=0.8`$ sample of 574 clusters below.
The first definition we consider is when a cluster had its most recent large $`\mathrm{\Delta }M`$, i.e. $`M_f/M_i1.2`$ in an interval $`\delta \tau =100h^1`$Mpc. In this section a mass โjumpโ refers to a mass gain of at least this much. For any cluster $`z_{\mathrm{jump}}`$ is the redshift of the most recent such jump. In our sample 13 clusters had no large mass jumps at any time after $`z2`$, two of these clusters are more massive than $`3\times 10^{14}h^1M_{}`$. Ref. RowThoKay04 found that 2 out of their sample of 20 did not have any large mass jumps. Another common definition is when the cluster has reached at least half its mass, for this time there is an analytic formula NusShe99 which was tested in (and improved using) simulations in Ref. SheTor04 . We define $`z_{1/2}`$ as the earliest output time when the cluster had at least half of its mass. The relation between these two definitions is shown in the upper left hand panel of Fig. 3. There is a large amount of scatter between the two definitions but there is a striking transition at $`z=0.5`$. For clusters which have not reached half of their present day mass before $`z=0.5`$, almost all have had a large $`\mathrm{\Delta }M`$ jump after reaching half-mass, i.e. they have at least half of their mass before they have their last large mass jump. However, for clusters which attained $`z_{1/2}`$ before redshift $`z=0.5`$ (286 clusters), slightly under one third (80) had their most recent mass jump even earlier (i.e. $`z_{\mathrm{jump}}>z_{1/2}`$). An intuitive reason for this is that clusters which reached half of their mass early on have not been gaining mass quickly and so are less likely to have large mass jumps after $`z_{1/2}`$. This general shape persists if one considers some other โformation timeโ such as $`z_{1/3}`$ or $`z_{3/4}`$ (with obvious definitions), the $`z`$ value where $`z_{\mathrm{jump}}`$ starts becoming larger (that is, earlier) than $`z_{1/3}(z_{3/4})`$ is larger (smaller) than that for $`z_{1/2}`$. A larger fraction (12/33) of the clusters with mass greater than $`3\times 10^{14}h^1M_{}`$ and $`z_{1/2}>0.5`$ had $`z_{\mathrm{jump}}>z_{1/2}`$. If one is interested in formation times in order to find undisturbed clusters, for lower redshifts it appears that $`z_{\mathrm{jump}}`$ may be a more conservative estimate.
A third definition of a cluster formation time is when the potential well of the object becomes deep enough to be considered โa clusterโ. We take this, somewhat arbitrarily, to be when the object has reached $`10^{14}h^1M_{}`$, and call this $`z_{14}`$.
All of these definitions rely on special events in the mass accretion history. An alternative is to use the whole history of the most massive progenitor. While this is less information than in the entire tree, it provides a global view of the formation process. We consider three fits used in the past for simulations. For galaxy sized halos, Ref. Wec02 found that the mass accretion histories could be fit by $`M(a)=M_0e^{2a_fz}=M_0e^{2z/(1+z_f)}`$, and that $`z_f`$ correlated well with other cluster properties, i.e. concentration. Ref. TKGK found that for clusters they could obtain better fits if they generalized the function to $`\stackrel{~}{M}(a)=M_0a^pe^{2\stackrel{~}{a}_fz}=a^pM(a)`$, where $`\stackrel{~}{a}_f=1/(1+\stackrel{~}{z}_f)`$. A third form was proposed by Zha , who fit the the mass accretion history to the form $`M_i/M_{\mathrm{tp}}=f\left(\rho _{\mathrm{vir}}(z_{\mathrm{tp}})/\rho _{\mathrm{vir}}(z_i)\right)`$ where $`\rho _{\mathrm{vir}}(z)`$ is the virial density at redshift $`z`$ computed from spherical top-hat collapse and $`f`$ interpolates between two power-laws Zha . The latter authors found that halos generally transition from a period of rapid accretion to a phase of slow accretion and used the transition redshift, $`z_{\mathrm{tp}}`$ as a proxy for โformationโ time.
We fit these three parameterizations to our sample, estimating $`z_f,\stackrel{~}{z}_f,z_{tp}`$ through a least squares fit of $`\mathrm{ln}(M_i/M_0)`$ against $`z_i`$, with all points being equally weighted. The values of $`z_f,\stackrel{~}{z}_f`$ are allowed free range, however $`z_{\mathrm{tp}}`$, due to its definition, is constrained to the range $`0<z<2`$ where we have data<sup>9</sup><sup>9</sup>9Also, we fit $`M`$ not $`M_{\mathrm{vir}}`$ but this should have little effect since $`M_{\mathrm{vir}}1.08M`$ at $`z=0`$ and $`1.10M`$ at high $`z`$.. We compare this formation time $`z_f`$ with $`z_{1/2}`$ at upper right in Fig. 3; they are correlated (as also found by Ref. TKGK ). We show in Fig. 3, bottom, the comparison of this โturning pointโ time $`z_{\mathrm{tp}}`$ with $`z_{\mathrm{jump}}`$ and $`z_f`$ above. The largest correlation shown is between $`z_{\mathrm{tp}}`$ and $`z_f`$, as shown in Fig. 3.
Although these parameterizations do describe general trends in the cluster mass histories, good fits were not found for all the cluster histories. (The best fits using our least squares criterion was for the models of Ref. TKGK , which have an extra parameter. However, similarly to TKGK , we found that many of the clusters had either $`p=0`$ (278/574) or $`\stackrel{~}{a}_f=0`$ (119/574).) This made results sensitive to the $`z`$-range used and the interpretation of the fit coefficient difficult. To illustrate the difficultly, we show in Fig. 4 examples of fits for 3 clusters.
In more detail, for the fit for $`z_f`$ (introduced by Ref. Wec02 ), we investigated the dependence on the range of $`z`$ used in the fit. Ref. TKGK fit their $`z_f`$ to histories for $`0z<0.5`$ and for $`0z<10`$ and found that the results were consistent. We found extremely large variations comparing $`0z<0.5`$ and $`0z<2.07`$. Our fits may have differed in success from theirs for two reasons. Our clusters tend to be more massive (theirs were between $`5.8\times 10^{13}h^1M_{}M_{180\rho _b}2.5\times 10^{14}h^1M_{}`$, which are even lower masses in terms of our mass definition), as they note, less massive halos tend to have better fits in general. Secondly, our sample is a lot larger, allowing us to sample a broader range of behavior. Even some clusters which gained more than 90% of their mass in the range we consider have $`z_f`$ which depends sensitively on the fit range. This is a simple consequence of the fact, which can be seen above, that the functional form does not well describe the shapes of the mass accretion histories of these objects. Finally we note that for clusters with recent mass jumps the fit prefers $`z_f<0`$. As it is very plausible to say that these systems are still in the process of forming we report these values above without renormalization.
The other two parameterizations also had difficulties with some subset of the clusters. It should be noted however that our redshift range was smaller than TKGK and that they did not report robustness of the $`p0`$ fits to smaller $`z`$ ranges. Ref. Zha actually considers a sample of similar size to ours for their fit, and large scatter is evident around the parameterized fit in their results as well, they find a useful correlation between $`z_{tp}`$ and concentration nonetheless.
In addition to comparing the different formation times to each other, we also use the cluster sample to get a snapshot of the formation time distribution. In Fig. 5 we show the distributions of these different formation times for all the clusters at $`z=0`$. Clusters which had no large mass jump since $`z=2`$ are put in the $`z_{\mathrm{jump}}=2`$ bin for completeness.
There is quite a spread in these formation times, however some trends are clear. By $`z=0.2`$ over 85% of the clusters have had their last large mass jump and over 81% of them have formed according to the definition of Ref. Wec02 , while 88% of them have at least half of their mass. The different definitions have different biases though, e.g., of the redshift formation times shown, only the definition of Ref. Wec02 can go negative.
A major part of our difficulty in fitting smooth curves to the mass accretion histories is the presence of significant mass jumps. These jumps are evident in the examples and are a general feature in cluster formation. The statistics of these jumps are of interest. Fig. 6 gives the distribution of the number of times since $`z=2`$ that the mass of each cluster changes by 20%. This is the same quantity considered by Ref. RowThoKay04 for their sample of 20 clusters more massive than $`1.2\times 10^{14}h^1M_{}`$ โ our sample is more than an order of magnitude larger allowing us to more completely characterize the distribution. We find that there is little dependence on the cluster mass โ lines for clusters above 1, 2 and $`3\times 10^{14}h^1M_{}`$ are shown. The average (median) number of jumps for each of these samples is 4.0 (4), 4.3 (4), 4.6 (5), respectively, with a wide spread. The right panel of Fig. 6 shows how much of the cluster mass is gained in these events โ slightly less than half of the clusters (272/574) get at least half of their mass in jumps of 20% of larger. For the $`M>3\times 10^{14}h^1M_{}`$ subsample the fraction is larger: 41/79.
With a several large number of jumps common, it appears difficult to find good smooth and simple parameterized fits for all of the histories. However, given the correlations between the different fit โformation timesโ and between these and e. g. concentration, these fits may be very useful approximations in a large number of cases.
### 4.2 Time dependence of recently merged or large $`\mathrm{\Delta }M`$ cluster fraction
The above shows that the present day cluster population has a very wide spread of formation times and frequent major disruptions. For observational purposes, another question is: how many clusters at a given redshift (โlookback timeโ) have had a large disruption recently (within a given โrelaxation timeโ).
The fraction of clusters with a large mass change $`\mathrm{\Delta }M`$ in a given interval is shown in Fig. 7, top, as the dotted lines. This is the total mass change, via mergers or accretion, within the 1.0 and $`2.5`$Gyr time intervals. The two lines show the fraction of clusters with final/initial mass $`1.2`$ and $`1.33`$ (these ratios were chosen by Refs. RowThoKay04 and GotKlyKra01 ), with the lower line for the higher mass ratio.
We may be interested in the most disruptive mass gains, which occur when the progenitor ratios are smallest. We refer to these events as major mergers. Figure 7, top, also shows the fraction of recently merged clusters as a function of time for minimum mass ratios 1:3, 1:5 and 1:10. For a lookback time of $`7`$Gyr ($`z0.8`$) for instance and the $`2.5`$Gyr relaxation time, this means considering predecessors $`9.5`$Gyr ago (at $`z1.6`$).
The error bars are calculated using binomial statistics. If $`M`$ of the $`N`$ clusters have merged the most likely<sup>10</sup><sup>10</sup>10The mean value is $`(M+1)/(N+2)`$. fraction is $`f=M/N`$ with variance:
$$\sigma _f^2=\frac{M(NM)+1+N}{(N+2)^2(N+3)}.$$
(1)
We use symmetric error bars. One can see here that the number of clusters which have had a recent major merger increases in the past, reaching a dramatic 80% for 1:10 or smaller mergers $`7`$Gyr ago (around $`z0.83`$) for a relaxation time of $`2.5`$Gyr. For shorter relaxation times fewer mergers have occurred, as expected. For fixed relaxation times the amount of accretion relative to merging changes with lookback time.
In the lower two panels different relaxation times are considered. For some phenomena associated with mergers the relaxation times of interest depend on the time of observation. We saw in the previous section that departures from the virial relation lasted several hundred Mpc of conformal time. This is not unexpected. A typical relaxation time is likely some fraction of the halo dynamical time. Since clusters have $`\overline{\rho }10^2\rho _{\mathrm{crit}}`$ and characteristic times scale as $`\rho ^{1/2}`$, typical timescales should be $`0.1t_\mathrm{H}`$ where $`t_\mathrm{H}H^1`$ is the Hubble time. Ref. RowThoKay04 noted that X-ray disturbances in their gas simulations lasted $`\delta \tau =300h^1`$Mpc (scaling as $`a`$ rather than the crossing time $`a^{3/2}`$). This corresponds to roughly $`1.3h^1`$Gyr at the present, and shorter times in the past. Simulations reported in RicSar01 have major merger related disturbances lasting approximately one crossing time. From RosBorNor02 , we take this to be $`1`$Gyr at the present. These two relaxation times are used in Fig. 7 bottom, with that found by Ref. RowThoKay04 at left and $`a^{3/2}`$Gyr at right. For the latter case, the $`a`$ used is for the initial time, before the merger, as clusters which merge later have a larger $`a`$ and thus longer relaxation time, and thus will also still be unrelaxed at the given lookback time. The ratio between the fractions with mergers and large $`\mathrm{\Delta }M`$ does not change as much with lookback time compared to the case where the relaxation time was independent of the lookback time. This suggests that equating large $`\mathrm{\Delta }M`$ jumps with mergers is more reliable for phenomena whose relaxation times scale with $`a`$.
In Fig. 8 we compare different samples of clusters, changing cosmology ($`\sigma _81.0`$) and mass (taking $`M>3.0\times 10^{14}h^1M_{}`$). At the top are the fractions of recently merged or recent large $`\mathrm{\Delta }M`$ clusters for both $`\sigma _8=0.8`$ and $`\sigma _8=1.0`$, for the relaxation time is $`\delta \tau =600h^1`$Mpc. At present this lookback time is slightly over $`2.5`$Gyr. The error bars are smaller for $`\sigma _8=1.0`$ because there are nearly twice as many clusters in the sample. The $`\sigma _8=1.0`$ clusters have had fewer mergers and mass jumps than those for $`\sigma _8=0.8`$; clustering is less evolved for $`\sigma _8=0.8`$ so a cluster at fixed mass is more likely to be forming in our redshift range. At the bottom, the recently merged or recent large $`\mathrm{\Delta }M`$ clusters are shown for more massive clusters, $`M>3.0\times 10^{14}h^1M_{}`$, with $`\sigma _8=0.8`$, for relaxation times of $`\delta \tau =300h^1`$Mpc (left) and for $`2.5`$Gyr (right). As there are significantly fewer clusters (79 rather than 574) at $`z=0`$, the error bars are a lot larger and are only shown for the top line. Comparing with the same plots in Fig. 7, it is seen that more of the massive clusters have had recent major mergers.
Values for $`\sigma _8=1`$ are reported in Tables 1-4 for four lookback times, along with $`N_{\mathrm{clus}}`$ so that the errors can be estimated.
For all relaxation times there is an increase in the fraction of recently merged or recent large $`\mathrm{\Delta }M`$ as the lookback time increases, i.e. there were more major mergers and large mass jumps in the past. The fraction depends on the relaxation time chosen. An overview is provided in Fig. 9 where we compare the ratios of merged and large $`\mathrm{\Delta }M`$ samples near the present to that of about $`6`$Gyr ago, for several different relaxation times in our two simulations. A doubling of the recently merged or recent large $`\mathrm{\Delta }M`$ fraction between the present and $`z0.67`$ is not uncommon, and sometimes even larger increases occur. The errors in these ratios for the relaxation time of $`\delta \tau =100h^1`$Mpc are quite large. For $`\sigma _8=0.8(1.0)`$, from 44(29)% for the 1:3 mergers to 32(28)% for $`M_f/M_i1.2`$. For the other relaxation times the errors on the ratio are the largest for the 1:3 mergers (10-25%) and smallest for the mass jumps with $`M_f/M_i1.2`$ (3-12%). The errors on the other cases fall in between. More details can be obtained from the earlier plots and the Tables.
Another question of interest might be: in how many major mergers do three or four predecessors have comparable mass? We took one particular definition as described in the methods section. These three-body mergers were fairly common, the exact numbers depend upon relaxation times considered. For the $`\sigma _8=0.8`$ case, about 15% of the 1:5 mergers were three-body for $`\delta \tau =30h^1`$Mpc and slightly less than 30% for $`\delta \tau =600h^1`$Mpc on average, with little evidence of growth with time. For $`\sigma _8=1.0`$ the fraction of three-body mergers to 1:5 major mergers was closer to 10% on average and seemed to grow with time for $`\delta \tau =300h^1`$Mpc. For $`\delta \tau =600h^1`$Mpc this growth was more evident, starting around 12% and growing to around 30%. For fixed relaxation times, 2.5 (1.0) Gyr, the fraction of three-body major mergers to 1:5 major mergers showed a definite increase with lookback time. For $`\sigma _8=0.8`$ it started at 30% (10%) and grew to 80% (35%) at lookback times of $`7`$Gyr. For $`\sigma _8=1.0`$ the fraction started at 12% (5%) and grew to 55% (25%) at lookback times of $`7`$Gyr. This means that even though a 1:5 merger might not be considered โmajorโ by some definitions, it is very possible that there are two of these going on at the same time, making these mergers more energetic.
There were many fewer four-body mergers. The largest effect was for the $`2.5`$Gyr relaxation time, where four-body mergers were about one third of the three-body mergers. In the largest case, for $`\sigma _8=0.8`$ and $`2.5`$Gyr relaxation time, the total fraction of clusters which had had a recent four-body merger only reached 5% or above for lookback times greater than $`4`$Gyr.
## 5 Conclusions
Clusters of galaxies represent the current endpoint of structure formation. As the largest systems which have had time to virialize in a universe with hierarchical structure formation, they make excellent laboratories for cosmology, large-scale structure and galaxy formation. The formation of galaxy clusters via a combination of mergers and accretion of smaller objects is crucial to understanding many of their present day properties.
In this paper we investigated the formation of galaxy clusters in some detail, and the major events which define this process, extending earlier studies TorBouWhi97 ; Tor98 ; Col99 ; GotKlyKra01 ; Zha ; RowThoKay04 ; TKGK ; Bus03 ; analytic discussed in ยง2. While galaxy cluster formation properties can be reliably calculated with the present generation of large, high-resolution N-body simulations, it is difficult to find specific numbers in the literature. To remedy this we have calculated, for a sample of hundreds of clusters, the degree of virialization, formation time, and for two values of the clustering amplitude $`\sigma _8`$, the fraction of disturbed galaxy clusters for many definitions of disturbed.
We began by showing the time histories of a few clusters to illustrate the typical formation pattern: periods of smooth accretion punctuated by large increases in mass. We then turned to various characterizations of this process, applied to our statistical sample as a whole.
We first calculated the ratio of kinetic to potential energy. At early times clusters are hotter than the vacuum virial relation (2KE$`=`$PE) would predict due to continuous infall of material. With the onset of $`\mathrm{\Lambda }`$ domination and the cessation of structure growth the excess kinetic energy drops and their mass accretion rate slows.
We then turned to formation time definitions in the literature; while no single number can capture the complexity of a cluster merging tree, the concept of a formation time encodes much useful information. We compared different formation time definitions which are relevant for different types of observations. For instance a recent large mass jump may not be of interest if one only wants to know when the cluster first was detectable by SZ decrement (which presumably depends more on the depth of the potential well), but may be relevant for studies which rely on a cluster being old enough to be dynamically relaxed. The different formation times we considered are correlated, with large scatter. Formation times relying upon a smooth parameterized fit to galaxy cluster histories were most correlated with each other, but the results were sensitive to the fit methodology, due to the large amount of recent sporadic mass gain. We compared the distributions of the different formation times for our whole cluster sample and for our more massive clusters, again finding some correlation but also significant scatter.
We then turned from smooth parameterizations to the characteristic abrupt jumps in mass over time for galaxy clusters. On average, galaxy clusters have had at least 4 large mass jumps since $`z2`$, this number increases with cluster mass; about half of the clusters get at least half of their mass in these jump events.
The above measurements were for our $`\sigma _8=0.8`$ sample of 574 clusters. The second part of the paper reported the time dependence of the fraction of clusters which have had a recent mass jump or major merger, using in addition a $`\sigma _8=1.0`$ sample of 909 clusters. These fractions should be of use for estimates of the number of โrelaxedโ clusters available in surveys (given a relaxation time for the phenomena of interest) and for helping to constrain relaxation times for phenomena associated with mergers or mass jumps where the phenomenonโs occurrence fraction has been measured but not its relaxation time. In Figs. 7-9 and tables 1-4 we give the merger fractions for a number of different situations. These fractions also serve to quantify the previously seen trend of more recent mergers (fractionally) at higher redshift. Comparing $`z=0`$ with $`z=0.67`$ (about $`6`$Gyr ago), the fraction of clusters which had a recent merger or mass jump is doubled by almost any definition and for some definitions the increase is even larger.
The simulations used here were performed on the IBM-SP at NERSC. JDC thanks T. Abel, S. Allen, G. Bryan, R. Gal, J. Hennawi, R. Kneissl, A. Kravtsov and P. Ricker for helpful discussions and was supported in part by NSF-AST-0205935. MJW was supported in part by NASA and the NSF. |
warning/0506/astro-ph0506368.html | ar5iv | text | # Jet Formation in Black Hole Accretion Systems I: Theoretical Unification Model
## 1. Introduction
Jets are a common outcome of accretion, yet the observed jet properties, such as collimation and speed, are not uniform between systems. This is despite the fact that the basic physics (general relativistic magnetohydrodynamics (GRMHD)) to describe such systems is black hole mass-invariant. Thus, it is worth-while to determine the unifying, or minimum number of, pieces of physics that would explain most of the features of gamma-ray bursts (GRBs), x-ray binaries, and active galactic nuclei (AGN) (Ghisellini & Celotti, 2002; Ghisellini, 2003; Meier, 2003). To understand jet formation requires at least explaining the origin of the energy, composition, collimation, and Lorentz factor. The goal of this paper, and the companion numerical models paper (McKinney, 2005b), is to explain these for GRBs, AGN, and x-ray binaries.
Primarily we discuss two types of jets: Poynting-dominated jets typically dominated in energy flux by Poynting flux and dominated in mass by electron-positron pairs for AGN and x-ray binaries, while dominated in mass by electron-proton pairs for GRBs ; and Poynting-baryon jets with about equal Poynting flux and rest-mass flux and dominated in mass by baryons. The latter are sometimes referred to as coronal outflows due to their origin. Generically this model is similar to, e.g., Sol et al. (1989), while here the motivation is based upon the results of recent GRMHD numerical models. This two-component jet model is one key to understanding the diversity of jet observations. The Poynting-dominated jet is likely powered by the Blandford-Znajek effect, while the Poynting-baryon jet is likely powered by both Blandford-Znajek power and the release of disk gravitational binding energy (McKinney, 2005a). Collimation of the polar Poynting-dominated jet may be due transfield balance against the broader Poynting-baryon jet or by self-collimating hoop stresses.
Among all the black hole accretion systems, it appears that the least unifiable is the observed emission. While the radiative physics is not black hole mass invariant, the observed differences suggest that the environment likely plays a significant role in the emission. For example, while both blazars and GRBs exhibit non-thermal emission, long-duration GRBs are harder with higher luminosity, while blazars are softer with higher luminosity (Ghirlanda et al., 2004, 2005). Also, GRBs lead to apparently most of the energy in $`\gamma `$-rays and less than $`10\%`$ to the sub-$`\gamma `$-ray afterglow (Piran, 2005). On the contrary, blazars apparently release only $`10\%`$ in $`\gamma `$-rays and the rest is produced in the radio lobe (Ghisellini & Celotti, 2002). Despite the difficulties in understanding the emission processes in some jet systems, the jet itself is probably produced by a universal process.
The disk and jet radiative physics are keys to understanding the evolution of the jet and why different systems have different terminal Lorentz factors. Through radiative annihilation of photons in AGN and x-ray binary systems, the radiative physics may illuminate the origin of jet composition by determining the electron-positron mass-loading the Poynting-lepton jet, and so the Lorentz factor of the jet. For GRBs, the radiative annihilation of neutrinos and the effect of Fick diffusion by free neutrons from the corona into the jet (Levinson & Eichler, 2003) may give an understanding of the Lorentz factor of the jet and the origin of baryon-contamination.
The rest of this section briefly reviews the types of black hole accretion systems and discusses jets in each. At the end is an outline of the paper.
### 1.1. GRBs
Neutron stars and black holes are associated with the most violent of post-Big Bang events: supernovae and some gamma-ray bursts (GRBs) and probably some x-ray flashes (XRFs) (for a general review see Woosley 1993; Wheeler, Yi, Hรถflich, & Wang 2000). Observations of a supernova light curve (SN2003dh) in the afterglow of GRB 030329 suggest that at least some long-duration GRBs are probably associated with core-collapse events (Stanek et al., 2003; Kawabata et al., 2003; Uemura et al., 2003; Hjorth et al., 2003).
Neutrino processes and magnetic fields are both important to understand core-collapse. In unraveling the mechanism by which core-collapse supernovae explode, the implementation of accurate neutrino transport has been realized to be critical to whether a supernova is produced in simulations (Messer et al. and collaborators, 1998). This has thus far been interpreted to imply that highly accurate neutrino transport physics is required, but this could also mean additional physics, such as a magnetic field, could play a significant role. Indeed, all core-collapse events may be powered by MHD processes rather than neutrino processes (Leblanc & Wilson, 1970; Symbalisty, 1984; Woosley & Weaver, 1986; Duncan & Thompson, 1992; Khokhlov et al., 1999; Akiyama, Wheeler, Meier, & Lichtenstadt, 2003). Core-collapse involves shearing subject to the Balbus-Hawley instability as in accretion disks (Akiyama, Wheeler, Meier, & Lichtenstadt, 2003). All core-collapse explosions are significantly polarised, asymmetric, and often bi-polar indicating a strong role of rotation and a magnetic field (see, e.g., Wang & Wheeler 1996; Wheeler, Yi, Hรถflich, & Wang 2000; Wang, Howell, Hรถflich, & Wheeler 2001; Wang et al. 2002; Wang, Baade, Hรถflich, & Wheeler 2003, and references therein). Possible evidence for a magnetic dominated outflow has been found in GRB 021206 (Coburn & Boggs, 2003), marginally consistent with a magnetic outflow directly from the inner engine (Lyutikov, Pariev, & Blandford, 2003), although these observations remain controversial.
Black hole accretion is the key source of energy for many GRB models. Collapsar type models suggest that a black hole forms during the core-collapse of some relatively rapidly rotating massive stars. The typical radius of the accretion disk likely determines the duration of long-duration GRBs (Woosley, 1993; Paczynski, 1998; MacFadyen & Woosley, 1999). An accretion disk is also formed as a result of a neutron star or black hole collisions with another stellar object (Narayan et al., 1992, 2001).
GRBs are believed to be the result of an ultrarelativistic jet. Indirect observational evidence of relativistic motion is suggested by afterglow achromatic light breaks and the โcompactness problemโ suggests GRB material must be ultrarelativistic with Lorentz factor $`\mathrm{\Gamma }100`$ to emit the observed nonthermal $`\gamma `$-rays (see, e.g., Piran 2005). Direct observational evidence for relativistic motion comes from radio scintillation of the ISM (Goodman, 1997) and measurements of the afterglow emitting region from GRB030329 (Taylor et al., 2004a, b).
Typical GRB jet models invoke either a hot neutrino-driven jet or a cold Poynting flux-dominated jet, while both allow for comparable amounts of the accretion energy to power the jet (Popham et al., 1999). A neutrino-driven jet derives its energy from neutrino annihilation from gravitational energy and the jet is thermally accelerated. However, strong outflows can be magnetically driven (Bisnovatyi-Kogan & Ruzmaikin, 1976; Lovelace, 1976; Blandford, 1976). In particular, black hole rotational energy can be extracted as a Poynting outflow (Blandford & Znajek, 1977).
### 1.2. X-ray Binaries
Long after their formation, neutron stars and black holes often continue to produce outflows and jets (Mirabel & Rodrรญguez, 1999). These include x-ray binaries (for a review see Lewin et al. 1995; McClintock & Remillard 2003), neutron star as pulsars (for a review see Lorimer 2001 on ms pulsars and Thorsett & Chakrabarty 1999 on radio pulsars) and soft-gamma ray repeaters (SGRs) (Thompson & Duncan, 1995, 1996; Kouveliotou et al., 1999). In the case of x-ray binaries, the companion starโs solar-wind or Roche-lobe forms an accretion disk. Many x-ray binaries in their hard/low state (and radio-loud AGN) show a correlation between the x-ray luminosity and radio luminosity (Merloni et al., 2003), which is consistent with radio synchrotron emission from a jet and x-ray emission from a geometrically thick, optically thin, Comptonizing disk.
Some black hole x-ray binaries have jets (Mirabel et al., 1992; Fender, 2003), such as GRS 1915+105 with apparently superluminal motion ($`\mathrm{\Gamma }3`$) (Mirabel & Rodriguez, 1994; Mirabel & Rodrรญguez, 1999; Fender & Belloni, 2004), but may have $`\mathrm{\Gamma }1.5`$ (Kaiser et al., 2004). Synchrotron radiation from the jet suggests the presence of a magnetized accretion disk. Observations of a broad, shifted, and asymmetric iron line from GRS 1915+105 is possible evidence for a relativistic accretion disk (Martocchia et al., 2002), although this feature could be produced by a jet component.
The standard paradigm is that relativistic jets from x-ray binaries are probably produced by the Blandford-Znajek effect. However, Gierliลski & Done (2004) suggest that at least some black holes, such as GRS 1915+105, have slowly rotating black holes. If this is correct, then another mechanism is required to produce jets. Indeed, jets or outflows are produced from systems containing NSs, young stellar objects, supersoft x-ray white dwarfs, symbiotic white dwarfs, and even UV line-driven outflows from massive O stars. Indeed, a baryon-loaded coronal outflow with $`\mathrm{\Gamma }1.53`$ can be produced from a black hole accretion disk and not require a rapidly rotating black hole (McKinney & Gammie, 2004). Nonrelativistic outflows were found even in viscous hydrodynamic simulations (Stone et al., 1999; Igumenshchev & Abramowicz, 1999, 2000; McKinney & Gammie, 2002). Such baryon-loaded outflows or jets are sufficient to explain most known x-ray binaries without invoking rapidly rotating black holes, and thus unifies such mildly relativistic jets in neutron star and black hole x-ray binaries.
### 1.3. AGN
Active galactic nuclei (AGN) have long been believed to be powered by accretion onto supermassive black holes (Zelโdovich, 1964; Salpeter, 1964). Observations of MCG 6-30-15 show an iron line feature consistent with emission from a relativistic disk with $`v/c0.2`$ (Tanaka et al., 1995; Fabian et al., 2002), although the lack of a temporal correlation between the continuum emission and iron-line emission may suggest it is a jet-related feature (Elvis, 2000).
AGN are observed to have jets with $`\mathrm{\Gamma }10`$ (Urry & Padovani, 1995; Biretta et al., 1999), even $`\mathrm{\Gamma }30`$ (Begelman et al., 1994; Ghisellini & Celotti, 2001; Jorstad et al., 2001), while some observations imply $`\mathrm{\Gamma }200`$ (Ghisellini et al., 1993; Krawczynski et al., 2002; Konopelko et al., 2003). Some radio-quiet AGN show evidence of weak jets (Ghisellini et al., 2004), which could be explained as a coronal outflow (McKinney & Gammie, 2004) and not require a rapidly rotating black hole. Observations imply the existence of a two-component jet structure with a Poynting jet core and a dissipative surrounding component (Ghisellini & Madau, 1996; Ghisellini et al., 2005). The energy structure of the jet and wind are important in understanding the feedback effect that controls size of the black hole and may determine the $`M\sigma `$ relation (Springel et al., 2004; Di Matteo et al., 2005).
### 1.4. Outline of Paper
ยง 2 summarizes the proposed unified model to explain jet formation in all black hole accretion systems.
ยง 3 discusses why ideal MHD must break down in magnetospheres and why the Goldreich-Julian charge density is never reached. A preliminary model is derived that describes the pair-loading and baryon-loading of the Poynting-dominated jet.
ยง 4 determines the Lorentz factor of Poynting-dominated jets. The GRB jet Lorentz factor is shown to be based upon electron-positron pair and baryon loading. We show that GRBs likely have electron-proton jets with Lorentz factor at large distances of $`100\mathrm{\Gamma }_{\mathrm{}}10^3`$.
Based upon pair creation rates for AGN and x-ray binaries, we show that relativity low radiatively efficient AGN, such as M87, have electron-positron jets with $`2\mathrm{\Gamma }_{\mathrm{}}10`$. Radiatively efficient systems, such as microquasar GRS1915+105, likely do not have Poynting-dominated lepton jets but rather the observed jets are a relativistic coronal outflow from the inner-disk.
ยง 5 discusses Poynting-baryon jets and how they can explain various observational features of jets in AGN and x-ray binaries.
ยง 6 summarizes the key results and fits from GRMHD numerical models (McKinney, 2005b) used in this paper.
ยง 7 discusses the results and their possible implications.
ยง 8 summarizes the key points.
Appendix A discusses breakdown of ideal MHD by electron-positron pair creation by radiative annihilation and electron-proton pair creation by ambipolar and Fick diffusion. See also the discussion in McKinney (2005b). Appendix B gives a succinct summary of conserved flow quantities in GRMHD used in section 4. Appendix C gives a derivation for the lab frame stationary GRMHD forces along and perpendicular to the flow (field) line in the lab frame. This elucidates the origin of acceleration and collimation. Appendix D gives the formulae for Comptonization and pair annihilation used in section 4.
## 2. GRMHD Pair Injection Model of Jet Formation
The jet energy, composition, collimation, and Lorentz factor are likely determined in a similar way for all black hole accretion systems. The particle acceleration mechanism and particle composition of the jet remained unexplained in McKinney & Gammie (2004). However, if field lines tie the black hole to large distances, then the source of matter is likely pair creation since the amount of matter that diffuses across field lines is much smaller (Phinney, 1983; Levinson & Eichler, 1993; Punsly, 2001). Thus, the Poynting-dominated jet composition is electron-positron pair dominated in AGN and x-ray binaries.
However, in GRB systems, free neutrons lead to baryon contamination due to Fick diffusion across the field lines and subsequent rapid collisionally-induced avalanche decay to an electron-proton plasma (Levinson & Eichler, 2003). The pair annihilation rates are much faster than the dynamical time, and due to the temperature decrease, the electron-positron pair rest-mass exponentially drops beyond the fireball formation near the black hole. Thus, the GRB jet composition is likely dominated by electron-proton pairs.
GRMHD numerical models confirmed that accretion of a thick disk with height ($`H`$) to radius ($`R`$) ratio of $`H/R0.1`$ with a homogeneous poloidal orientation self-consistently creates large scale fields that tie the black hole to large distances (McKinney & Gammie, 2004; Hirose et al., 2004). Accretion of an irregular field loads the jet with baryons and lowers the speed of the jet. However, the existence of a mostly uniform field threading the disk arises naturally during core-collapse supernovae and NS-BH collision debris disks. In AGN and solar-wind capture x-ray binary systems, the accreted field is probably uniform (Narayan et al., 2003; Uzdensky & Spruit, 2005). Roche-lobe overflow x-ray binaries, however, might accrete a quite irregular field geometry. The field geometry that arrives at the black hole, after travelling from the source of material (molecular torus, star(s), etc.) to the black hole horizon, likely depends sensitively on the reconnection physics.
The reason why each system has some observed Lorentz factor has not been well-understood. One key idea of this paper is that the terminal Lorentz factor is determined by the toroidal magnetic energy per unit pair mass density energy near the location where pairs can escape to infinity (beyond the so-called โstagnation surfaceโ). Put another way, the Lorentz factor is determined by the energy flux per unit rest mass flux for the rest-mass flux in pairs beyond the stagnation surface. For GRBs, neutron diffusion is crucial to explain (and limit) the Lorentz factor. For AGN and x-ray binaries, since a negligible number of baryons cross the field lines, pair-loading is crucial to determine the Lorentz factor of the Poynting-dominated jet since this determines the rest-mass flux or density.
Figure 1 shows the basic picture for GRB systems, while figure 2 shows the basic picture for AGN and x-ray binary systems. An accreting, spinning black hole creates a magnetically dominated funnel region around the polar axis. The rotating black hole drives a Poynting flux into the funnel region, where the Poynting flux is associated with the coiling of poloidal magnetic field lines into toroidal magnetic field lines. The accretion disk emits neutrinos in a GRB model ($`\gamma `$-ray and many soft photons for AGN and x-ray binaries) that annihilate and pair-load the funnel region within some โinjection region.โ For GRB systems, neutrons Fick-diffuse across the field lines and collisionally decay into an electron-proton plasma.
Many pairs (any type) are swallowed by the black hole, but some escape if beyond some โstagnation surface,โ where the time-averaged poloidal velocity is zero and positive beyond. Pairs beyond the stagnation surface are then accelerated by the Poynting flux in a self-consistently generated collimated outflow. In the electromagnetic (EM) jet, the acceleration process corresponds to a gradual uncoiling of the magnetic field and a release of the stored magnetic energy that originated from the spin energy of the black hole.
One key result of this paper is that the release of magnetic energy need not be gradual once the toroidal field dominates the poloidal field, in which case pinch (and perhaps kink) instabilities can occur and lead to a nonlinear coupling (e.g. a shock) that converts Poynting flux into enthalpy flux (Eichler, 1993; Begelman, 1998). In the proposed GRB model, this conversion reaches equipartition and the jet becomes a โmagnetic fireball,โ where the toroidal field instabilities drive large variations in the jet Lorentz factor and jet luminosity.
In AGN systems, nonthermal synchrotron from shock-accelerated electrons and some thermal synchrotron emission releases the shock energy until the synchrotron cooling times are longer than the jet propagation time. For AGN, jet acceleration is negligible beyond the extended shock zone, as suggested for blazars beyond the โblazar zoneโ (Sikora et al., 2005). In x-ray binary systems, the shock is not as hot and also unlike in the AGN (at least those like M87) case the jet can be optically thick. Thus these x-ray binary systems self-absorbed synchrotron emit if they survive Compton drag.
For all these systems, at large radii patches of energy flux and variations in the Lorentz factor develop due to toroidal instabilities. These patches in the jet could drive internal shocks and at large radii they drive external shocks with the surrounding medium. The EM jet is also surrounded by a mildly relativistic matter coronal outflow/jet/wind, which is a material extension of the corona surrounding the disk. This Poynting-baryon, coronal outflow collimates the outer edge of the Poynting-dominated jet, which otherwise internally collimates by hoop stresses. The luminosity of the Poynting-baryon jet is determined, like the Poynting-dominated jet, by the mass accretion rate, disk thickness, and black hole spin.
This model is studied analytically in this paper, while in a companion paper we study this model numerically using axisymmetric, nonradiative, GRMHD simulations to study the self-consistent process of jet formation from black hole accretion systems (McKinney, 2005b). Those simulations extend the work of McKinney & Gammie (2004) by including pair creation (and an effective neutron diffusion for GRB-type systems) to self-consistently treat the creation of jet matter, investigating a larger dynamic range in radius, and presenting a more detailed analysis of the Poynting-dominated jet structure.
Unless explicitly stated, the units in this paper have $`GM=c=1`$, which sets the scale of length ($`r_gGM/c^2`$) and time ($`t_gGM/c^3`$). The mass scale is determined by setting the (model-dependent) observed (or inferred for GRB-type systems) mass accretion rate ($`\dot{M}`$\[$`gs^1`$\]) equal to the accretion rate through the black hole horizon as measured in a simulation. So the mass is scaled by the mass accretion rate at the horizon, such that $`\rho _{0,disk}\dot{M}[r=r_H]t_g/r_g^3`$ and the mass scale is then just $`m\rho _{0,disk}r_g^3=\dot{M}[r=r_H]t_g`$. Unless explicitly stated, the magnetic field strength is given in Heaviside-Lorentz units, where the Gaussian unit value is obtained by multiplying the Heaviside-Lorentz value by $`\sqrt{4\pi }`$.
The value of $`\rho _{0,disk}`$ can be determined for different systems. For example, a collapsar model with $`\dot{M}=0.1\mathrm{M}_{}s^1`$ and $`M3\mathrm{M}_{}`$, then $`\rho _{0,disk}3.4\times 10^{10}\mathrm{g}\mathrm{cm}^3`$. M87 has a mass accretion rate of $`\dot{M}10^2\mathrm{M}_{}\mathrm{yr}^1`$ and a black hole mass of $`M3\times 10^9\mathrm{M}_{}`$ (Ho, 1999; Reynolds et al., 1996) giving $`\rho _{0,disk}10^{16}\mathrm{g}\mathrm{cm}^3`$. GRS 1915+105 has a mass accretion rate of $`\dot{M}7\times 10^7\mathrm{M}_{}\mathrm{yr}^1`$ (Mirabel & Rodriguez, 1994; Mirabel & Rodrรญguez, 1999; Fender & Belloni, 2004) with a mass of $`M14\mathrm{M}_{}`$ (Greiner et al., 2001), but see Kaiser et al. (2004). This gives $`\rho _{0,disk}3\times 10^4\mathrm{g}\mathrm{cm}^3`$. This disk density scales many of the results of the paper.
## 3. Breakdown of ideal-MHD
Pair creation is critical to understand the physics of the highly magnetized, evacuated funnel region that is associated with a Poynting-dominated jet. Pair creation is often invoked in order to use the force-free electrodynamics or ideal MHD approximation in a black hole magnetosphere (see, e.g., Blandford & Znajek 1977). However, in MHD where rest-mass is treated explicitly, pair creation is not simply a passive mechanism to short out spark gaps, which is the mechanism invoked to allow the use of the force-free approximation.
Pair creation (and neutron Fick diffusion for GRB-type systems) determines the matter flow in the magnetosphere, and thus the matter-loading of any Poynting jet that emerges (Phinney, 1983; Punsly, 1991; Levinson, 2005). As shown below, these sources of mass loading self-consistently determine the Lorentz factor of the Poynting-dominated jet and allows one to understand why black hole accretion systems, while following the mass-invariant GRMHD equations of motion, show a variety of jet Lorentz factors.
For GRBs, the radiative physics and neutron diffusion is shown to determine the Lorentz factor of the Poynting-dominated jet. For AGN and x-ray binaries, the radiative physics is shown to determine the Lorentz factor of the Poynting-dominated jet by determining its energy and mass-loading.
The ideal MHD approximation (or force-free approximation in magnetically dominated regions) has been shown to be a reasonably valid theoretical framework to describe most of the nonradiative dynamically important accretion physics around a black hole in GRBs, AGN, and black hole x-ray binary systems (Phinney, 1983; McKinney, 2004). This approximation is the foundation of most studies of jets and winds. The ideal MHD approximation is a good approximation to describe these flow properties except 1) in current sheets, which is not treated explicitly in this paper ; 2) where pair creation contributes a nonnegligible amount of rest-mass, internal energy, or momentum density; 3) if the Goldreich-Julian (GJ) charge density is larger than the number density of charge carriers ; and 4) if the rest-mass flux due to ambipolar and Fick diffusion is negligible.
The first goal is to show that radiative annihilation into pairs establishes a density of pairs much larger than the Goldreich-Julian density. The Goldreich-Julian charge density is never reached because pair creation is completely dominated by neutrino annihilation in GRB-type systems and photon annihilation in AGN and x-ray binary systems.
Notice that the breakdown of ideal MHD is required in order to extract black hole spin energy from a stationary, axisymmetric system. Wald (1974) showed that a rotating black hole induces a parallel electric current in the surrounding magnetosphere such that the plasma becomes nondegenerate (i.e. $`E^iB_i0`$). Bekenstein & Oron (1978) argued that if the ideal MHD approximation were valid, that no energy could be extracted from a black hole. This is because since $`u^r<0`$ at the horizon, and the radial energy flux can be written as $`T_t^r=E\rho u^r`$ (where $`E`$ is conserved along each flow line ; see appendix B), then to extract net energy ($`T_t^r>0`$) from the black hole requires $`E<0`$. However, in the ideal MHD approximation $`E>0`$ at $`r\mathrm{}`$, and by conservation of $`E`$ along each flow line, then $`E>0`$ on the horizon as well. However, based upon arguments by Goldreich & Julian (1969), Blandford & Znajek (1977) argued that as the magnetosphere is evacuated to the Goldreich-Julian charge density, the parallel electric current separates the charges. The Goldreich-Julian rest-mass density for a species of electrons is
$$\rho _{GJ}m_e\frac{\mathrm{\Omega }_HB}{2\pi cq},$$
(1)
where $`B`$ is the magnetic field strength and $`q`$ is the electron charge. Once the parallel electric current is sufficiently large, electrons are accelerated across the potential gap and photons can be emitted by curvature radiation or inverse Compton scattering. These high energy photons either self-interact or are involved in a magnetic bremsstrahlung interaction, ultimately leading to electron-positron pairs. These pairs would continuously short the induced potential difference. However, this picture does not establish how the resulting pair plasma flow behaves.
Why have ideal GRMHD numerical models demonstrated the Blandford-Znajek effect (Koide, Shibata, Kudoh, & Meier, 2002; McKinney & Gammie, 2004; De Villiers et al., 2005a; Komissarov, 2005) ? These ideal GRMHD numerical models implicitly break the ideal MHD approximation in the required way to allow the extraction of energy from the black hole. For all the initial conditions and field geometries explored by McKinney & Gammie (2004) using the โidealโ GRMHD numerical model of an accreting black hole, they always find that a highly magnetized polar region forms and any material in this magnetosphere is either rapidly driven into the black hole or driven out in a wind or jet. They find that strong field lines tie the black hole horizon to large radii. Thus, necessarily these ideal MHD models break the ideal MHD approximation at a stagnation point where the poloidal velocity $`u^p=0`$. Necessarily matter is created (at least) in this location since matter inside this surface goes into the black hole and matter beyond it goes away from the black hole. This aspect is similar to the charge-starved magnetosphere models where there is a spark-gap (Ruderman & Sutherland, 1975) where particles are generated (for a review see Levinson 2005). Once the magnetosphere reaches an axisymmetric, quasi-stationary state, then the departure from the ideal-MHD condition can be measured as deviations from conservation of the conserved flow quantities given in equations B3 to equations B12.
Notice that for a realistic accretion disk the BZ power is different than the typically used estimates (McKinney, 2005a). For $`j0.5`$, they find that the efficiency in terms of the mass accretion rate is
$$\eta _{EM,tot}=\frac{P_{tot}}{\dot{M}c^2}15\%\left(\frac{\mathrm{\Omega }_H}{\mathrm{\Omega }_H[j=1]}\right)^4,$$
(2)
and
$$\eta _{EM,jet}=\frac{P_{jet}}{\dot{M}c^2}7\%\left(\frac{\mathrm{\Omega }_H}{\mathrm{\Omega }_H[j=1]}\right)^5,$$
(3)
where $`r_g=GM/c^2`$, $`\mathrm{\Omega }_H=jc/(2r_H)`$ is the rotation frequency of the hole, $`r_H=r_g(1+\sqrt{1j^2})`$ is the radius of the horizon for angular momentum $`J=jGM^2/c`$, and $`j=a/M`$ is the dimensionless Kerr parameter, where $`1j1`$. However, net electromagnetic energy is not extracted for $`j0.5`$ (including retrograde accretion) when an accretion disk is present (McKinney & Gammie, 2004). This high efficiency is a result of the near equipartition of the magnetic field strength ($`(B^r)^2`$) in the polar region at the horizon and the rest-mass density in the disk at the horizon. If the black hole has $`j0.9`$, then $`1\%`$ of the accreted rest-mass energy is emitted back as Poynting flux in the form of a jet and $`3\%`$ is emitted back in total (so obviously $`2\%`$ goes into the disk and corona โ about equally it turns out).
### 3.1. GRB Pair Creation Model
In GRB models, such as the collapsar model, neutrino/anti-neutrino annihilation provides a source of electron-positron pairs at a much larger density than the Goldreich-Julian density and so the magnetosphere is not charge starved. The cross-field magnetic diffusion for charged species is negligible in such systems. However, free neutrons diffuse across the field lines and load the jet with an electron-proton plasma (Levinson & Eichler, 2003), and this effect is considered in the next section.
For the collapsar model, the jet has $`B3\times 10^{15}\mathrm{Gauss}`$ (McKinney, 2004) and $`j0.9`$, which gives
$$\rho _{GJ}10^9gcm^3(\mathrm{Collapsars}).$$
(4)
One can compare this to the density of pairs produced by neutrino annihilation for the GRB collapsar model. One can use the results in table 3 and figure 9b of Popham et al. (1999) and the results in table 1 in MacFadyen & Woosley (1999), which are fairly well fit to power laws, such that for models with $`\dot{M}0.1\mathrm{M}_{}s^1`$, and $`M=3\mathrm{M}_{}`$
$`\eta _{\nu \overline{\nu },ann}`$ $``$ $`{\displaystyle \frac{L_{\nu \overline{\nu },ann}}{\dot{M}c^2}}`$ (5)
$``$ $`1\%\left({\displaystyle \frac{\alpha }{0.1}}\right)\left({\displaystyle \frac{j}{0.9}}\right)^7\left({\displaystyle \frac{\dot{M}}{0.1\mathrm{M}_{}s^1}}\right)^{3.8},`$
where this fit is based on an average between the conservative and optimistic models of MacFadyen & Woosley (1999). This assumes an average neutrino energy of $`10`$MeV from the disk. This says that for the collapsar model with $`j=0.9`$ that about $`1\%`$ of the rest-mass accreted is given back as positron-electron pairs due to neutrino annihilation, which is similar to the Poynting flux from the black hole that goes into the jet region as from equation 3.
Published results of neutrino annihilation rates as a function of position (Popham et al., 1999) can be used to obtain a preliminary model of pair creation and incorporated into a GRMHD model. The details of this preliminary model end up not affecting the results, and a more self-consistent model is left for future work. The results primarily depend on the overall annihilation luminosity and the basic radial dependence of the energy injected as pairs.
Figure 9 and table 3 in Popham et al. (1999) and table 1 of MacFadyen & Woosley (1999) can be used to obtain approximate radial and height dependent fits of the energy density rate of depositing pairs into the jet region. Their figure 9 shows that the height and radial dependence of the pair annihilation luminosity per unit distance. These follow approximate power laws or exponential laws for $`j0.2`$. A reasonable fit is that
$$P[R]e^{\frac{Rr_g}{2.4r_g}}$$
(6)
and
$$Q[z]e^{\frac{zr_g}{3.5r_g}}$$
(7)
for the luminosity per unit distance. The coefficient is determined by the total annihilation luminosity ($`L_{\nu \overline{\nu },ann}`$). As in Popham et al. (1999), photon null geodesic transport in curved spacetime is neglected such that
$$L_{\nu \overline{\nu },ann}2\pi A_{r_g}^{\mathrm{}}_{\theta =0}^{arctan[H/R]}P[R]Q[z]r^2\mathrm{sin}\theta d\theta dr.$$
(8)
Figure 6 of Popham et al. (1999) can be used to obtain the disk thickness to radius ratio
$$H/R0.1\left(\frac{r}{2r_g}\right)^{2/3}$$
(9)
for $`\dot{M}=0.1\mathrm{M}_{}s^1`$. This allows one to determine that $`AL_{\nu \overline{\nu },ann}/191`$. Thus the energy generation rate can be written as
$$\frac{\dot{e}_{\nu \overline{\nu },ann}}{\dot{\rho }_{0,disk}c^2}\frac{\eta _{\nu \overline{\nu }}}{N_A}P[R]Q[z](\mathrm{Collapsars}),$$
(10)
where $`N_A191`$ and we have defined $`\dot{\rho }_{0,disk}\dot{M}/r_g^3`$. However, the above $`H/R`$ assumes the jet fills around the disk. Rather, there is likely a thick corona between the disk and jet (McKinney & Gammie, 2004). Motivated by those simulations and the simulations discussed in this paper, the jet region is presumed to exist within
$$\theta _j1.0\left(\frac{r}{3r_g}\right)^{1/3}$$
(11)
for $`r100r_g`$. In that case $`N_A70`$. Thus, a significant fraction of the pairs are absorbed by the corona. However, the corona could also contribute significantly to neutrino production (Ramirez-Ruiz & Socrates, 2005). Thus, $`70N_A191`$ are reasonable limits.
Some fraction of the total energy deposited goes into pair rest-mass, pair internal energy, radiation, and pair momentum. Here $`f_\rho `$ denotes the fraction turned into lab-frame mass, $`f_h`$ the fraction turned into lab-frame internal energy and radiation, and $`f_m`$ the fraction turned into momentum energy. Thus $`1=f_\rho +f_h+f_m`$. The pair rest-mass density creation rate is defined as
$$\frac{\dot{\rho }_{e^{}e^+}}{\dot{\rho }_{0,disk}}=f_\rho \frac{\dot{e}_{\nu \overline{\nu },ann}}{\dot{\rho }_{0,disk}},$$
(12)
where $`\rho _{e^{}e^+}=\rho _{0,e^{}e^+}u^t`$.
One can obtain a rough density measure in the injection region by assuming the characteristic time scale for moving the pairs once they have formed is the light crossing time at the stagnation surface $`t_{stag}t_g(r_{stag}/r_g)`$. Then $`\rho _{0,e^{}e^+}\dot{\rho }_{0,e^{}e^+}t_{stag}`$. With $`\rho _{0,disk}\dot{\rho }_{0,disk}t_{stag}`$,
$$\frac{\rho _{e^{}e^+}}{\rho _{0,disk}}f_\rho \left(\frac{\dot{e}_{\nu \overline{\nu },ann}}{\dot{\rho }_{0,disk}c^2}\right)\left(\frac{r_{stag}}{r_g}\right).$$
(13)
Notice that many pairs fall into the black hole, so do not contribute to the jet rest-mass or energy. Only those pairs beyond the stagnation surface survive the gravity of the black hole, such that the total annihilation luminosity into the jet is
$$L_{esc}=2\pi _{r=r_{stag}}^{\mathrm{}}_{\theta =0}^{\theta _j}\dot{e}_{\nu \overline{\nu },ann}r^2\mathrm{sin}\theta drd\theta ,$$
(14)
which is a similar integral as performed before. However, notice that particles injected with $`r<r_{stag}`$, by definition, fall into the black hole since they are inside the stagnation surface where $`u^p<0`$. This is unlike the BZ-power, which in steady state is well-defined and conserved through the stagnation surface (Levinson, 2005). For $`r_{stag}=r_g`$, all the injected mass reaches infinity by definition. For $`r_{stag}=10r_g`$, only $`11\%`$ of the mass injected reach infinity. Because any mass injected lost to the black hole is of no consequence to the acceleration at large distances, then the true efficiency of pairs that load the jet is
$$\eta _{esc}=\frac{L_{esc}}{\dot{M}c^2}$$
(15)
rather than $`\eta _{\nu \overline{\nu },ann}`$ for all $`r_{stag}`$. One can show that $`\eta _{esc}=\eta _{\nu \overline{\nu },ann}`$ for $`r_{stag}=r_g`$, but is reduced to $`\eta _{esc}0.11\eta _{\nu \overline{\nu },ann}`$ for $`r_{stag}=10r_g`$ due to the loss of pairs into the black hole.
The pairs annihilate after formed by neutrino annihilation. Equation D7 gives the pair annihilation rate. For GRB models, such as the collapsar model, the pair annihilation timescale is $`t_{pa}10^{16}\mathrm{s}GM/c^310^5\mathrm{s}`$ and similarly all along the jet. Thus, the pairs annihilate and form a thermalized fireball. A fraction $`f_\rho +f_h`$ of the energy injected is turned into a electron-positron pair-radiation fireball. The typical angle of scattering neutrinos gives $`f_mf_\rho +f_h`$ (Popham et al., 1999). Thus, for the fireball formation region $`f_\rho +f_hf_m0.5`$ within factors of a few, and this is independent of the energy of the neutrinos or the efficiency of annihilation.
All of the mass energy thermalizes into the fireball with a temperature given by equation A.2. The formation fireball pair mass plus pair internal energy density plus radiation internal energy is
$$\frac{q_{0,tot}u^tu_t+p_{e^{}e^+}+p_\gamma }{\rho _{0,disk}c^2}\left(\frac{\dot{e}_{\nu \overline{\nu },ann}}{\dot{\rho }_{0,disk}c^2}\right)\left(\frac{r_{stag}}{r_g}\right),$$
(16)
where $`q_{0,tot}(\rho _{0,e^{}e^+}c^2+u_{0,e^{}e^+}+u_{0,\gamma })`$. This equation connects the energy injection process in terms of the GRMHD equations of motion for a given energy-at-infinity injection rate. See also appendix A.
Using the discussion here and the equations in appendix A, one can show that for $`\alpha =0.1`$, $`\dot{M}=0.1\mathrm{M}_{}s^1`$, and $`j=0.9`$, the fireball temperature is $`T10^{10}\mathrm{K}1\mathrm{M}\mathrm{e}\mathrm{V}`$ in the injection region. Thus pair rest-mass energy is nearly in equipartition with the pair internal energy and radiation. In particular, $`f_\rho f_h/8.5`$. Since $`f_\rho +f_h0.5`$, then $`f_\rho 0.05`$, $`f_h0.45`$, and $`f_m0.5`$.
Beyond the initial fireball formation, the Boltzmann factor gives that once the temperature drops below $`T6\times 10^9`$K the number of pairs decreases exponentially with temperature. However, the fireball is optically thin only at much larger radii of $`r10^810^{10}r_g`$. So until that radius, the radiation provides an inertial drag on the remaining pair plasma and the gas is radiation dominated.
For example, for $`j=0.9`$, $`\dot{M}=0.1\mathrm{M}_{}s^1`$, $`r_{stag}=6r_g`$, then the initial fireball rest-mass is
$$\frac{\rho _{e^{}e^+}}{\rho _{0,disk}}10^6(r_{stag}=6r_g).$$
(17)
If instead the model were an $`\alpha =0.01`$ model, then the efficiency would be about $`10`$ times less and the density ratio would be about $`10`$ times less at
$$\frac{\rho _{e^{}e^+}}{\rho _{0,disk}}10^7(\alpha =0.01).$$
(18)
In any case this is roughly
$$\rho _{e^{}e^+}10^3gcm^3,$$
(19)
which is about 12 orders of magnitude larger than the GJ density given in equation 4, and so the black hole is far from starved of charges.
### 3.2. GRB Baryon Contamination
Neutron diffusion across field lines leads to baryon contamination of the (otherwise) electron-positron-radiation jet. The neutrons Fick-diffuse (Levinson & Eichler, 2003) or diffusion due to ambipolar diffusion (see appendix A) across the field lines and baryon-load the jet. The neutrons undergo a fast collisional avalanche into protons and electrons that are then carried along with the electromagnetic and Compton-thick outflow. As shown in appendix A, the mass injection rate of neutrons (and so proton+electrons) is
$$\dot{M}_{inj,Fick}7\times 10^5\dot{M}_{acc}$$
(20)
where the density of electron-proton plasma is
$$\rho _{pe^{}}3\times 10^7\left(\frac{r}{r_g}\right)^{4/3}\rho _{0,disk}.$$
(21)
Notice that this mass injection rate and density are comparable to the injection-region rest-mass density in electron-positron pairs for $`\alpha =0.01`$. As mentioned above, father out in the jet the electron-positron pairs annihilate and contribute only an additional $`10\%`$ to the internal energy. Thus the total internal energy is sufficient to describe the gas without including pair annihilation and the rest-mass in baryons is sufficient to describe the gas rest-mass for models with $`\alpha =0.01`$ and for all models for $`r10r_g`$ where the pair-density is exponentially smaller than the baryon density in the jet. Thus, the injection of โpairsโ described in the previous section can also approximately account for the Fick diffusion of neutrons. This fact is exploited to simulate the collapsar model in McKinney (2005b).
### 3.3. AGN and X-ray Binaries Pair Creation Model
In AGN and x-ray binaries, the pairs are produced by scattering of $`1\mathrm{M}\mathrm{e}\mathrm{V}`$ $`\gamma `$-rays with other photons (for a review see Phinney 1983 and chpt. 6, 9, and 10 in Punsly 2001). These $`\gamma `$-rays could be produced by, for example, Comptonization of disk photons through a gas of relativistic electrons, non-thermal particle acceleration in shocks (see, e.g., Nishikawa et al. 2003), or reconnection events. For example, the (non-radiative) simulations of McKinney & Gammie (2004) show an extended corona that could be source of Comptonization. They also find an edge between the corona and funnel that contains frequent shocks with sound Mach number $`M100`$. Also, they found reconnection is quite vigorous in the plunging region at $`r36r_g`$, leading in part to the hot coronal outflow.
These sites of Comptonization, shocks, and reconnection are likely sources of the requisite $`\gamma `$-rays. In place of a detailed model of these processes, it is assumed that some fraction ($`f_\gamma `$) of the true bolometric luminosity is in the form of these $`\gamma `$-rays that do collide with softer photons in the funnel region. For a bolometric luminosity $`L_{bol}\eta _{eff}\dot{M}c^2`$, then the annihilation efficiency is
$$\eta _{\gamma \gamma ,ann}f_\gamma \frac{L_{bol}}{\dot{M}c^2}=f_\gamma \eta _{eff},$$
(22)
where $`\eta _{eff}`$ is the total radiative efficiency, which could include emission from both the disk and jet near the base. Notice that $`\eta _{eff}`$ depends on the black hole spin, among other things. However, the value of $`\eta _{eff}`$ is obtained from (model-dependent) values for the mass accretion rate and bolometric luminosity.
Extrapolating from gamma-ray observations of black hole x-ray binary systems suggests that in either the quiescent or outburst phase, $`f_\gamma 1\%`$ of the true bolometric luminosity is in the form of $`>1`$MeV photons (see, e.g., Ling & Wheaton 2005). These are likely produced quite close to the black hole. For an injection region with a half opening angle for the jet of $`\theta _j57^{}`$ (McKinney & Gammie, 2004), about $`(2\theta _j)^2/(4\pi )1/3`$ of these photons enter the jet region. Thus, it is assumed that a large fraction of these $`>1`$MeV photons give up their energy into producing pairs in the funnel region with some fraction of the energy going into rest-mass ($`f_\rho `$). These pairs do not annihilate so form a collisionless plasma (see appendix A). See also Punsly (2001) for why $`f_\gamma 1\%`$ is reasonable, based on assuming the infall rate is equal to the pair creation rate. In general $`f_\gamma `$ depends on the state of the accretion flow, and a self-consistent determination is left for future work.
As in the collapsar case this gives us a density rate or a typical density. In this case the stagnation surface is close to the black hole since the emission is likely always optically thin, thus $`r_{stag}=3r_g`$. The author knows of no calculations that give a radial dependence for the annihilation energy generation rate. A reasonable choice is the same radial dependence as for neutrino annihilation, since while the efficiency of neutrino scattering is much lower, the radial structure is determined by a similar calculation. The typical thick disk ADAF model assumes the electrons and protons are weakly coupled, which gives a disk thickness mostly independent of radius such that $`H/R0.85`$ and is weakly dependent on $`\alpha `$ (Narayan & Yi, 1995). For the AGN and x-ray binaries, jets are presumed to occur in the presence of a thick (ADAF-like) disk close to the black hole. While observations show that the disk at $`r6r_g`$ undergoes state transitions, steady jets are only observed in the low-hard state when the disk is likely geometrically thick. For this $`H/R`$, equation 8 gives $`A=L_{\gamma \gamma ,ann}/53`$ or $`N_A=53`$. Now equations 12 and 13 also apply for photon annihilation but with $`N_A=53`$.
For AGN accretion disks, Phinney (1983) already showed that $`\gamma `$ rays from the accretion disk corona interact with x-rays to produce electron-positron pairs in sufficient density above the Goldreich-Julian density. For example, M87 has a nuclear bolometric luminosity of $`L_\gamma 2\times 10^{42}\mathrm{erg}\mathrm{s}^1`$ and a mass accretion rate of $`\dot{M}10^2\mathrm{M}_{}\mathrm{yr}^1`$ (Ho, 1999; Reynolds et al., 1996) giving $`\eta _{eff}4\times 10^3`$. Unlike neutrinos, $`\gamma `$-rays are efficient at creating pairs and most of the $`\gamma `$-rays are at around $`1`$MeV so the fraction of the energy put into rest-mass is $`f_\rho 1`$. If $`f_\gamma 1\%`$, then
$$\rho _{e^{}e^+}10^{23}gcm^3(f_\gamma 1\%).$$
(23)
The field strength in M87 is $`B0.1\mathrm{to}50\mathrm{G}\mathrm{a}\mathrm{u}\mathrm{s}\mathrm{s}`$ (McKinney, 2004). For $`j=0.9`$ this gives that
$$\rho _{GJ}10^{32}gcm^3(\mathrm{M87}).$$
(24)
This is about 9 orders of magnitude lower than the pair creation established density, so the black hole is not charge starved. For M87,
$$\frac{\rho _{e^{}e^+}}{\rho _{0,disk}}10^7(\mathrm{M87}).$$
(25)
For black hole x-ray binaries a similar calculation is performed. For example, GRS 1915+105 has a mass accretion rate of $`\dot{M}7\times 10^7\mathrm{M}_{}\mathrm{yr}^1`$ and a bolometric luminosity $`L10^{40}\mathrm{erg}\mathrm{s}^1`$ (Mirabel & Rodriguez, 1994; Mirabel & Rodrรญguez, 1999; Fender & Belloni, 2004) with a mass of $`M14\mathrm{M}_{}`$ (Greiner et al., 2001), but see Kaiser et al. (2004). This gives $`\eta _{eff}0.26`$. If $`f_\gamma 1\%`$ and $`f_\rho 1`$, then
$$\rho _{e^{}e^+}10^9gcm^3(f_\gamma 1\%).$$
(26)
The field strength is $`B10^6\mathrm{Gauss}`$ if the disk is in the thick ADAF-like state (McKinney, 2004). For $`j=0.9`$ this gives that
$$\rho _{GJ}10^{19}gcm^3(\mathrm{GRS1915}+105).$$
(27)
This is about 9-10 orders of magnitude larger than the GJ density, so the black hole is not charge starved. For GRS 1915+105,
$$\frac{\rho _{e^{}e^+}}{\rho _{0,disk}}10^5(\mathrm{GRS1915}+105).$$
(28)
Notice that GRS 1915+105, and many x-ray binaries, are more radiatively efficient than most AGN. This means x-ray binaries have jets loaded with more pair-mass density per unit disk density. This will impact the presence and speed of any Poynting-lepton jet, as described in the next section.
Finally, clearly temporal variations in the disk structure and mass accretion rate directly affect the actual pair creation rate in the jet region. A self-consistent treatment of this (time-dependent) radiative physics is left for future work.
## 4. Jet Lorentz Factor
The Lorentz factor of the jet can be measured either as the current time-dependent value, or, using information about the GRMHD system of equations, one can estimate the Lorentz factor at large radii from fluid quantities at small radii. The Lorentz factor as measured by a static observer at infinity is
$$\mathrm{\Gamma }u^{\widehat{t}}=u^t\sqrt{g_{tt}}$$
(29)
in Boyer-Lindquist coordinates, where no static observers exist inside the ergosphere. This is as opposed to $`Wu^t\sqrt{1/g^{tt}}`$, which is the Lorentz factor as measured by the normal observer as used by most numerical relativists.
For a GRMHD model, to determine the Lorentz factor at $`r\mathrm{}`$, notice that equations B11 and B12 specify that $`E`$ and $`L`$ are conserved along each flow line. For flows with a magnetic field that is radially asymptotically small compared to near the black hole, $`B_\varphi 0`$ as $`r\mathrm{}`$. Also, for a nonradiative fluid, the internal energy is lost to kinetic energy, and so $`h1`$ as $`r\mathrm{}`$. Thus $`\mathrm{\Gamma }_{\mathrm{}}=u_t[r=\mathrm{}]`$. Now, since $`\mathrm{\Phi }`$ and $`\mathrm{\Omega }_F`$ are conserved along a flow line, then trivially
$`\mathrm{\Gamma }_{\mathrm{}}`$ $`=`$ $`E=hu_t+\mathrm{\Phi }\mathrm{\Omega }_FB_\varphi `$ (30)
$`u_{\mathrm{}}^{\widehat{\varphi }}`$ $`=`$ $`L=hu_\varphi +\mathrm{\Phi }B_\varphi ,`$ (31)
where $`h=(\rho _0+u_g+p)/\rho _0`$ is the specific enthalpy, $`\mathrm{\Phi }`$ is the conserved magnetic flux per unit rest-mass flux, $`\mathrm{\Omega }_F`$ is the conserved field rotation frequency, and $`B_\varphi `$ is the covariant toroidal magnetic field. $`E`$ and $`L`$ simply represent the conserved energy and angular momentum flux per unit rest-mass flux. Notice the matter and electromagnetic pieces are separable.
Since $`E`$ is the hydrodynamic plus magnetic energy flux per unit rest-mass flux, $`B_\varphi 0`$ simply represents the conversion of Poynting flux into rest-mass flux and $`h1`$ represents conversion of thermal energy into rest-mass flux. These are what accelerate the flow. Alternatively stated for the magnetic term, from equation C14, the fluid is accelerated by the magnetic toroidal gradients associated with the Poynting outflow.
A rough estimate of the Lorentz factor at infinity $`\mathrm{\Gamma }_{\mathrm{}}`$ can be estimated by assuming all the enthalpy flux or all the Poynting flux is lost to rest-mass flux that reach infinity. Then from equation B7, one can break up the matter and electromagnetic numerator and average the numerator and denominator separately to obtain
$$\mathrm{\Gamma }_{\mathrm{}}=\mathrm{\Gamma }_{\mathrm{}}^{(MA)}+\mathrm{\Gamma }_{\mathrm{}}^{(EM)}\left(\eta _{esc}+\eta _{EM}\right)\left(\frac{\dot{M}_{disk}}{\dot{M}_{esc}}\right),$$
(32)
where here โescโ refers to those pairs that escape the black hole gravity. In steady-state, only those pairs beyond the stagnation surface escape to feed the jet.
Before estimating the value of the Lorentz factor for various systems, the toroidal field is shown to be the source of the acceleration in ideal MHD. This also allows a probe of the structure of the jet, which is not possible in the above averages. In ideal-MHD, $`E`$ and $`L`$ in equations B11 and B12 give the Lorentz factor at infinity and the angular momentum per particle at infinity.
The enthalpy is just
$$h=(\rho +u+p)/\rho =1+\frac{4u}{3\rho }=1+\frac{4(1f_\rho )}{3f_\rho }$$
(33)
for a relativistic gas of electron-positron pairs that have $`f_\rho `$ of energy into rest-mass and the rest into thermal or momentum energy. Now the magnetic term is
$$\mathrm{\Phi }\mathrm{\Omega }_FB_\varphi =\frac{B^\varphi B_\varphi }{\rho _0u^t(v^\varphi /\mathrm{\Omega }_F1)}.$$
(34)
Notice that $`E`$ diverges for $`v^\varphi =\mathrm{\Omega }_F`$ (or $`u^r=u^\theta =0`$), where the ideal MHD approximation breaks down. Thus, $`\mathrm{\Gamma }_{\mathrm{}}`$ is also determined by the non-ideal MHD physics of particle creation in that region.
For an extended particle creation region, $`\mathrm{\Gamma }_{\mathrm{}}`$ depends on the mass, enthalpy, and momentum injected as a function of radius and $`\theta `$ in the jet region (Levinson, 2005; Punsly, 2001). For a narrow ($`\delta rr_g`$), yet distributed, particle creation region, then plausibly ideal MHD is completely reestablished at a slightly larger radius where $`v^\varphi \mathrm{\Omega }_F\mathrm{\Omega }_H/2`$ (McKinney & Gammie, 2004). In this case $`\mathrm{\Phi }\rho \mathrm{\Omega }_Fu^t/B^\varphi `$ (or $`u^r/B^r\mathrm{\Omega }_Fu^t/B^\varphi `$). Thus, the magnetic piece is
$$\mathrm{\Phi }\mathrm{\Omega }_FB_\varphi \frac{B^\varphi B_\varphi }{\rho },$$
(35)
where $`\rho =\rho _0u^t`$ is the lab-frame density. Written in Boyer-Lindquist coordinates in an orthonormal basis, then the magnetic piece is
$$\mathrm{\Gamma }_{\mathrm{}}^{(EM)}\frac{(B^{\widehat{\varphi }})^2}{\rho },$$
(36)
where $`B^{\widehat{\varphi }}=\sqrt{g_{\varphi \varphi }}B^\varphi `$. Thus, in Boyer-Lindquist coordinates
$$\mathrm{\Gamma }_{\mathrm{}}1+\frac{4(1f_\rho )}{3f_\rho }+\frac{(B^{\widehat{\varphi }})^2}{\rho c^2}$$
(37)
for $`r>r_H`$. Thus the fluid energy at infinity is due to conversion of thermal energy and toroidal field energy into kinetic energy. The actual value of $`\mathrm{\Gamma }_{\mathrm{}}`$ depends on how narrow is the injection region and the location of the stagnation surface. The relevant magnetic field is the toroidal magnetic field strength beyond the injection region where ideal MHD is mostly reestablished. Notice that $`\sigma =b^2/\rho _0`$ is often used to parameterize the Lorentz factor for a magnetically dominated flow, while perhaps equation 37 is more useful.
The polar field on the black hole horizon is approximately monopolar even for rapidly rotating holes with $`j0.95`$ (McKinney & Gammie, 2004). The monopole field solution can then be used, when properly normalized, to give the functional dependence of $`B^{\widehat{\varphi }}`$ near the black hole. Then one may use equation 37 to obtain the approximate terminal Lorentz factor. First, the black hole emits a Poynting flux given by equation 33 in McKinney & Gammie (2004) with the use of their section 2.3.2 for the magnetic field. In an asymptotic expansion, which is good to factors of 2 for any $`j`$ even at $`r=3r_g`$, then
$$B^{\widehat{\varphi }}C\frac{j}{8r/r_g}\mathrm{sin}\theta .$$
(38)
Modifications at higher black hole spin are within factors of $`1.5`$ in magnitude and there is a factor of $`1.5`$ enhancement near the polar axis compared to lower spin. In order to use this simple BZ monopole, the jet regionโs โCโ coefficient can be obtained from GRMHD numerical models (McKinney & Gammie, 2004). For a $`H/R0.20.4`$ disk model, they find that the normalization constant โ$`C`$โ in the BZ monopole solution is
$$C0.7\sqrt{\rho _{0,disk}c}$$
(39)
(see equations 47-49 in McKinney & Gammie 2004). This clearly states that the toroidal field in the polar region is nearly in equipartition with the rest-mass density in the disk. One can plug this into equation 37, but this equation would have limited usefulness since the injection region is broad and the density is difficult to estimate. However, notice that $`\mathrm{\Gamma }_{\mathrm{}}`$ has angular structure if the magnetic term dominates since then $`\mathrm{\Gamma }_{\mathrm{}}\mathrm{sin}^2\theta `$. This is important to jet structure described in McKinney (2005b). Otherwise the result has the same qualitative features as equation 46.
### 4.1. Lorentz Factors in Collapsar Systems
This section shows that, without invoking super-efficient neutrino emission mechanisms, only the Blandford-Znajek driven process can drive the flow to the necessary minimal Lorentz factor to avoid the compactness problem and unequivocally generate a GRB. Some GRBs require up to $`\mathrm{\Gamma }500`$ (Lithwick & Sari, 2001), so any invoked mechanism must be able to explain this. Some observation/models suggest some bursts have even $`\mathrm{\Gamma }1000`$ (Soderberg & Ramirez-Ruiz, 2003).
Equation 32 can be written as
$$\mathrm{\Gamma }_{\mathrm{}}=\left(\eta _{esc}+\eta _{EM,jet}\right)\left(\frac{\dot{M}}{\dot{M}_{esc}}\right).$$
(40)
Equation 15 gives $`\eta _{esc}`$, which accounts for pair capture by the black hole. Here $`\eta _{EM}`$ is given in equation 3 and as generally noted before, $`\eta _{esc}=\eta _{\nu \overline{\nu },ann}`$ for $`r_{stag}=r_g`$, but is reduced to $`\eta _{esc}0.11\eta _{\nu \overline{\nu },ann}`$ for $`r_{stag}=10r_g`$ due to the loss of pairs into the black hole. Based upon GRMHD numerical models studied in McKinney (2005b), a likely value is $`r_{stag}=5r_g`$, for which $`\eta _{esc}0.5\eta _{\nu \overline{\nu },ann}`$.
GRB-type systems are different than AGN and x-ray binary systems, because neutron diffusion baryon-loads the jet. The indirect injection of protons and electrons is shown to dominate the rest-mass in the jet because the electron-positron pairs annihilate to negligible rest-mass. The Fick diffusion mass injection rate is given by equation A10 and is
$$\dot{M}_{inj,Fick}7\times 10^5\dot{M}_{acc}$$
(41)
(Levinson & Eichler, 2003), and so $`\dot{M}_{esc}=\dot{M}_{inj,Fick}`$.
Plugging in the efficiencies for neutrino annihilation (equation 5) and Poynting-dominated jet efficiency (equation 3) into equation 40, one finds for $`\dot{M}0.1\mathrm{M}_{}s^1`$ for $`j0.1`$ that
$$\mathrm{\Gamma }_{\mathrm{}}140\left(g_{\nu \overline{\nu },ann,esc}+g_{EM,jet}\right)$$
(42)
where
$$g_{\nu \overline{\nu },ann,esc}\left(\frac{\eta _{esc}}{\eta _{\nu \overline{\nu },ann}}\right)\left(\frac{\alpha }{0.1}\right)\left(\frac{j}{0.9}\right)^7\left(\frac{\dot{M}}{0.1\mathrm{M}_{}s^1}\right)^{3.8}$$
(43)
where for $`\dot{M}0.1\mathrm{M}_{}s^1`$, the power of $`3.8`$ sharply levels out to $`0`$ (Di Matteo et al., 2002). Also,
$$g_{EM,jet}7\left(\frac{j}{1+\sqrt{1j^2})}\right)^5.$$
(44)
These $`g`$โs are just the normalized the efficiencies. It turns out that for $`j0.9`$ that the efficiencies are similar for the shown normalization of parameters. Thus, one might expect that they contribute equally to the energy of the jet.
If $`r_{stag}=r_g`$ and $`\alpha =0.1`$, then for the collapsar model with $`\dot{M}=0.1\mathrm{M}_{}s^1`$ and $`j=0.9`$, then neutrino annihilation and BZ power are equal and typically give $`\mathrm{\Gamma }100140`$, which is sufficient to avoid the compactness problem for typical bursts.
However, $`r_{stag}5r_g`$ is more likely. The neutrino-driven term then gives $`\mathrm{\Gamma }_{\mathrm{}}70`$ and is dominated by the BZ-driven term that stays at $`\mathrm{\Gamma }_{\mathrm{}}100`$. This is a result of the loss of pairs into the black hole. Notice that the results of (Popham et al., 1999; Di Matteo et al., 2002) and others do not include Kerr geometry to trace null geodesics. While the efficiency for annihilation should increase due to gravitational focusing, more pairs are also absorbed by the black hole. Thus, it is unlikely that gravitational focusing helps the neutrino-driven mechanism.
Also, $`\alpha =0.1`$ probably considerably overestimates the viscous dissipation rate of a true MHD flow. Despite typically $`\alpha =0.1`$ being used by many authors studying viscous models of disks, Stone & Pringle (2001) showed that $`\alpha =0.01`$ is more representative of MHD disks near the black hole. This implies that the energy generation rate in the disk and the generation rate of neutrinos is lower by about an order of magnitude. For normal neutrinos this gives $`\mathrm{\Gamma }_{\mathrm{}}14`$, insufficient to explain GRBs.
Another problem for the neutrino annihilation mechanism is that, with the inclusion of optically thick neutrino transport, the efficiency of neutrino emission is another few times lower for the collapsar model (Di Matteo et al., 2002). This gives $`\mathrm{\Gamma }_{\mathrm{}}5`$, clearly a serious problem for the neutrino annihilation model of GRBs.
The annihilation efficiency approximately scales with the average neutrino energy (Popham et al., 1999). One must invoke super-efficient neutrinos with an average neutrino energy of $`210`$MeV (Ramirez-Ruiz & Socrates, 2005) in order to obtain a neutrino annihilation power comparable to the BZ power. However, this is near the peak neutrino energy estimated to come from a hot corona (Ramirez-Ruiz & Socrates, 2005), and the corona is not expected to be the dominant source of neutrinos, so the average neutrino energy should be smaller.
Notice that choosing $`j=1`$ only increases the neutrino efficiency, and so $`\mathrm{\Gamma }_{\mathrm{}}`$, by a factor of $`2`$. Such a black hole spin is only achievable when $`H/R0.01`$ (Gammie, Shapiro, & McKinney, 2004), which is not representative of any GRB model (Kohri et al., 2005).
So, without evoking super-efficient neutrino emission mechanisms and unreasonably large neutrino annihilation efficiencies based upon only optically thin emission, one must turn to the Poynting flux to drive the jet. For $`j=0.9`$ one has $`\mathrm{\Gamma }_{\mathrm{}}100`$, sufficient to avoid the compactness problem is many GRBs and only invokes an obtainable black hole spin. Unlike the neutrino annihilation case, the BZ efficiency has been computed self-consistently from GRMHD numerical models of GRB-type disks (McKinney, 2005b).
Just as the neutrino-case has a proposed super-efficient mechanism, there are super-efficient magnetic models that would increase the terminal Lorentz factor. Eventually after accreting a large amount of magnetic flux, the magnetic pressure dominates the ram pressure of the accretion flow and suspends the flow (Igumenshchev et al., 2003; Narayan et al., 2003). GRMHD numerical models are suggested to not have simulated for long enough to see this effect. In this case a larger amount of magnetic flux threads the black hole. In this magnetically arrested disk (MAD) model, the field strength is comparable to the rest-mass density in this case rather than only a fraction of it. Then, the BZ efficiency is $`100\%`$ for $`j=1`$ and an estimate of the jet efficiency is about $`50\%`$, which is $`10`$ times larger than previously. Thus potentially $`\mathrm{\Gamma }1000`$ is achievable with $`j=0.9`$. Only super-efficient neutrinos with average energy $`2100`$MeV can obtain such a jet power.
At the moment there is an insufficient study of the neutrino annihilation rates for disks that have large optically thick regions, so $`\mathrm{\Gamma }_{\mathrm{}}`$ is not directly estimated for BH-NS or NS-NS collisions. However, at higher mass accretion rates the neutrino emission efficiency levels off rather than increasing (Di Matteo et al., 2002), suggesting these systems should also be dominated by Poynting-flux energy like for collapsars.
In summary, the neutrino-driven mechanism is probably dominated by the BZ power. Corrections due to optically thick neutrino-transport, a realistic choice for $`\alpha `$ based upon MHD models, and the loss of pairs into the black hole are the primary reducing factors compared to previous expectations. However, the BZ-effect generates sufficiently energetic GRBs while only invoking an obtainable black hole spin.
Also, while a hydrodynamic jet mixes and can destroy jet structure, an electromagnetic jet can internally evolve and could have a distribution of Lorentz factors in the jet. This is what is found in McKinney (2005b). Thus, while the electromagnetic jet may on average have $`\mathrm{\Gamma }100`$, the core of the jet may have $`\mathrm{\Gamma }1000`$. Thus, without invoking super-efficient mechanisms, only an electromagnetic BZ-driven jet can explain all observed/inferred GRB Lorentz factors.
Based upon the density scaling from the simulations in McKinney (2005b) that are summarized in section 6.2, and based upon appendix D, the fireball is optically thick to Compton scattering until $`r10^810^{10}r_g`$. Simulations discussed in McKinney (2005b) show that Poynting flux is continuously converted to heat in shocks and so all the energy flux leads to acceleration. Thermal acceleration occurs until the fireball is optically thin. This acceleration reaches $`\mathrm{\Gamma }1001000`$ before the internal shocks generate the nonthermal emission.
### 4.2. Lorentz Factors in AGN and X-ray Binary Systems
This section shows that the disk+jet radiative physics is crucial to determine the Lorenz factor of jets. It is shown that radiatively inefficient AGN, such as M87, should have jets with $`2\mathrm{\Gamma }10`$, while radiatively efficient systems, such as GRS 1915+105, may have jets with $`\mathrm{\Gamma }2`$, but they may be Compton dragged to nonrelativistic velocities.
For AGN and x-ray binaries equation 32 can be written as
$$\mathrm{\Gamma }_{\mathrm{}}=\frac{1}{f_\rho }\left(\eta _{esc}+\eta _{EM}\right)\left(\frac{\dot{M}c^2}{L_{esc}}\right),$$
(45)
for a fraction $`\dot{M}_{esc}=f_\rho L_{esc}`$ of mass that escapes the gravity of the hole. Equation 15 shows that the mass-energy loading of the jet is reduced, compared to the total injected mass-energy, due to loss of pairs into the black hole. Since $`\eta _{esc}L_{esc}/(\dot{M}c^2)`$, then
$$\mathrm{\Gamma }_{\mathrm{}}\frac{1}{f_\rho }\left(1+\frac{\eta _{EM}}{\eta _{esc}}\right)$$
(46)
Here $`\eta _{EM}`$ is given in equation 3. As noted before, $`\eta _{esc}=\eta _{\gamma \gamma ,ann}`$ for $`r_{stag}=r_g`$, but is reduced to $`\eta _{esc}0.11\eta _{\gamma \gamma ,ann}`$ for $`r_{stag}=10r_g`$ due to the loss of pairs into the black hole.
For these systems $`r_{stag}3r_g`$, since the disk is likely optically thin and emits harder radiation closer to the black hole. This gives $`\eta _{pairs}/\eta _{ann}0.83`$ for the ADAF model of $`H/R`$. Here the fraction of energy in rest-mass is $`f_\rho 1`$ due to the efficiency of photon-photon annihilation for $`>1`$MeV photons off the plentiful softer photons. As discussed before, observations suggest that the fraction of bolometric luminosity in $`\gamma `$-rays is $`f_\gamma 0.01`$. Since the radiative efficiency should also depend on black hole spin, it is assumed a typical system has $`j0.9`$ since that is a plausible equilibrium spin (Gammie, Shapiro, & McKinney, 2004). A detailed model of the radiative efficiency as a function of black hole spin is left for future work.
Equation 46 then gives that
$$\mathrm{\Gamma }_{\mathrm{}}1+\left(\eta _{eff}\right)^1.$$
(47)
For M87, $`\eta _{eff}=0.004`$ so $`\mathrm{\Gamma }_{\mathrm{}}=250`$, while for GRS1915+105 $`\eta _{eff}=0.26`$ so $`\mathrm{\Gamma }_{\mathrm{}}=5`$.
As cited in the introduction, the observed apparent Lorentz factor of AGN jets is $`\mathrm{\Gamma }_{\mathrm{}}30`$, while inferred Lorentz factor of some AGN jets is $`\mathrm{\Gamma }_{\mathrm{}}200`$, in basic agreement with the above estimates. However, in particular for M87 this is rather large. Also, the GRS1915+105 estimate is a bit large.
The key difference between the collapsar event and AGN or x-ray binaries is that in the collapsar case the photon luminosity (including Compton upscattered by $`\mathrm{\Gamma }^2`$) is negligible compared to the jet luminosity, but see Ghisellini et al. (2000); Lazzati et al. (2004). In M87 the bolometric luminosity $`L_{bol}`$ is almost that of the jet luminosity $`P_{jet}`$. In GRS 1915+105 the bolometric luminosity is greater than the jet luminosity.
The below discussion is a preliminary check on how radiative processes affect the above results. The results of simulations from McKinney (2005b) are invoked in order to obtain the density and magnetic structure of the Poynting-lepton jet as summarized in section 6.2. The simulations also show that at $`r10^210^4r_g`$, any remaining Poynting flux is shock-converted into enthalpy flux until they are in equipartition. A self-consistent simulation with synchrotron emission would then show the continuous loss of Poynting flux until the synchrotron cooling timescale is longer than the jet propagation timescale. This still suggests the jet magnetic field is finally in equipartition, but that much of the energy is lost and so cannot acceleration the jet.
The jet, rather than just disk, radiative physics is necessary in order to explain why $`\mathrm{\Gamma }_{\mathrm{}}=250`$ is not achieved in M87 and $`\mathrm{\Gamma }_{\mathrm{}}=5`$ is not achieved in GRS 1915+105. For M87, much of the Poynting energy that leads to $`\mathrm{\Gamma }_{\mathrm{}}250`$ is converted to heat by shocks, which is then lost to nonthermal and some thermal synchrotron emission. Thus the Lorentz factor of the Poynting-lepton jet achieved by $`r10^210^3r_g`$ is likely the maximum obtainable. Numerical models of M87 discussed in McKinney (2005b) show that $`2\mathrm{\Gamma }10`$. The Poynting-lepton jet in GRS 1915+105 is likely destroyed by Compton drag or at best $`\mathrm{\Gamma }_{\mathrm{}}2`$. Other black hole accretion systems with different mass accretion rates and radiative efficiencies have to be independently checked. The below discussion of the disk+jet radiative physics should be considered preliminary since a full radiative transport is necessary to obtain a completely self-consistent solution.
#### 4.2.1 Jet Destruction by Bulk Comptonization
The previous section showed that jets from AGN and x-ray binaries survive loading by pair production from $`\gamma `$-ray photons, but the produced jet may not survive Compton drag (bulk Comptonization) by the relatively soft photons emitted by the disk.
The simulation-based results are used to determine the optical depth as given in equations D3 and D5 to compute the perpendicular and parallel optical depths to Comptonization. For M87 $`\tau _{}6\times 10^6`$ and $`\tau _{}6\times 10^6`$ at $`r5r_g`$ (stagnation surface where jet starts) and $`\tau _{}4\times 10^6`$ at $`r=120r_g`$. For GRS 1915+105, $`\tau _{}11`$ and $`\tau _{}4`$ (stagnation surface) and $`\tau _{}3`$ at $`r=120r_g`$. Thus M87 is not Compton dragged by the disk photons, while GRS 1915+105 is likely strongly Compton dragged by photons that originate near the base of the jet and travel up through the jet or across the jet, or by synchrotron self-Compton drag.
A Compton-dragged jet has a limited Lorentz factor that reaches an equilibrium between decelerating and accelerating radiative processes. Here is it assumed that most of the disk emission is at the base of the jet. Then the relevant scenario for GRS 1915+105 is the one where all disk seed photons that enter the jet are scattered. An isotropic disk luminosity $`L_{bol}`$ shining on a conical jet with half-opening angle $`\theta _j`$ dumps a luminosity of $`L_{seed}\theta _j^2L_{bol}/4`$ into the jet. The photons effectively mass-load the jet and an equilibrium Lorentz factor is reached, where
$$\mathrm{\Gamma }_{\mathrm{}}\left(\frac{P_{jet}}{2L_{seed}}\right)^{1/3}\left(\frac{2\eta _{EM,jet}}{\eta _{eff}\theta _j^2}\right)^{1/3}$$
(48)
for a cold beam of electrons (see, e.g., Broderick 2004). The thermal Lorentz factor is comparable to the bulk Lorentz factor, so thermal corrections are not significant. For GRS 1915+105 this gives a nonrelativistic velocity ($`\mathrm{\Gamma }1`$) for the jet if most of the emission enters the base. Only if most of the emission enters far ($`r10^2r_g`$) from the base is up to $`\mathrm{\Gamma }2`$ possible. Thus is unlikely, so the Poynting-lepton jet that forms in radiatively efficient systems, such as GRS 1915+105, are Compton dragged to nonrelativistic velocities. Clearly the Lorentz factor is sensitive to the disk thickness, the emission from the disk, and the structure of the jet. Thus these estimates should be treated as preliminary. A self-consistent radiative transfer calculation is left for future work.
In summary, a radiatively efficient system loads the jet with more pairs from the larger number of $`\gamma `$-rays. This sets the maximum possible Lorentz factor to be smaller than for radiatively inefficient systems. For systems with relatively high density jets, such as X-ray binaries, the larger radiative efficiency also leads to an optically thick jet that can be Compton dragged.
#### 4.2.2 Jet Destruction by Pair Annihilation
Equation D7 gives the pair annihilation rate. For AGN, such as M87, the pair annihilation timescale is $`t_{pa}10^{11}\mathrm{s}GM/c^310^4\mathrm{s}`$ and for a jet propagation time $`t_{jet}r/c`$, a lower limit is $`t_{pa}/t_{jet}10^7`$ all along the jet. Thus, most pairs do not annihilate. See appendix A on how this affects the fluid approximation. See also Ghisellini et al. (1992).
For X-ray binary GRS1915+105, $`t_{pa}2\times 10^4\mathrm{s}GM/c^37\times 10^5\mathrm{s}`$ and $`t_{pa}/t_{jet}2`$ all along the jet. Thus, some nonnegligible fraction of the pairs annihilate. This also contributes to the destruction of the Poynting-lepton jet in X-ray binary systems since much of this radiative energy is lost at $`r150r_g`$ where the jet is optically thin along the jet and $`r300r_g`$ where the jet is optically thin perpendicular to the jet.
#### 4.2.3 Heat Loss by Synchrotron Emission
For M87 it was estimated that $`\mathrm{\Gamma }_{\mathrm{}}250`$, which is inconsistent with observations. However, much of the Poynting energy is converted to internal energy in shocks induced by toroidal field instabilities (McKinney, 2005b). Thus, the synchrotron cooling time might be sufficiently fast to release this internal energy that would otherwise accelerate the flow through thermal acceleration. For numerical models described in McKinney (2005b) that correspond to M87, an equipartition โmagnetic fireballโ forms between $`10^2r_g`$ and $`10^3r_g`$. If this energy could be released, then the shocks would again resume and all the Poynting and thermal energy would be lost. For an equipartition magnetic fireball half the energy is thermal, so the jet internal energy is $`u_{e^{}e^+}125\rho _{0,e^{}e^+}c^2`$ and so the thermal Lorentz factor is $`\mathrm{\Gamma }_e125`$. The synchrotron cooling time in the lab frame is
$$t_{syn}\frac{\mathrm{\Gamma }_{bulk}\mathrm{\Gamma }_em_ec^2}{P_{syn}}6\pi \frac{\mathrm{\Gamma }_{bulk}m_ec^2}{\mathrm{\Gamma }_e\sigma _TcB^2}.$$
(49)
This gives that
$$t_{syn}10^4s\left(\frac{r}{r_g}\right)^{1.4}(r<390r_g,\mathrm{M87})$$
(50)
and
$$t_{syn}1s\left(\frac{r}{r_g}\right)^3(r>390r_g,\mathrm{M87}).$$
(51)
For a typical lab frame jet propagation time of $`t_{jet}r/c`$,
$$\frac{t_{syn}}{t_{jet}}\left(\frac{r}{r_g}\right)^{0.4}(r<390r_g,\mathrm{M87})$$
(52)
and
$$\frac{t_{syn}}{t_{jet}}10^4\left(\frac{r}{r_g}\right)^2(r>390r_g,\mathrm{M87}).$$
(53)
Hence, one would not expect synchrotron cooling to take much of the internal energy away.
However, the โmagnetic fireballโ forms by shock heating and electrons are dramatically accelerated in such relativistic collisionless shocks. The shocks generate a power law (nonthermal) distribution of electrons, where much of the energy is carried by high-energy electrons (Begelman et al., 1984; Blandford & Eichler, 1987; Achterberg et al., 2001; Fender & Maccarone, 2003; Keshet & Waxman, 2005; Fender et al., 2005). Typically the distribution is $`N(E)E^{2.22}`$. For a pair plasma the maximum energy is limited by synchrotron losses (Achterberg et al., 2001), and the resulting synchrotron emission has photon energies of $`E25\mathrm{\Gamma }_e\mathrm{MeV}`$. This gives $`\gamma `$-ray and up to possibly TeV emission beamed along the jet, as in blazars. For example, Mrk 421 shows 15 minute variability, which for a mass of $`1.9\times 10^8\mathrm{M}_{}`$ would suggest an emission size on the order of the horizon size (Punch et al., 1992; Gaidos et al., 1996). However, relativistic time effects with $`\mathrm{\Gamma }_{bulk}10`$ place these emissions at $`r10^2r_g`$, which coincides with the shock-heated transfast region discussed in (McKinney, 2005b).
Thus, a significant portion of the shock-heated internal energy should be emitted by shock accelerated electrons and lost through the optically thin jet. In the shocks, inverse Compton also contributes to emission of high-energy photons and the loss of internal energy. Shock-induced population inversions may generate cyclotron masers at shock sites and lead to large brightness temperatures (Begelman et al., 2005).
Thus, it is expected that much of the jet is cold with $`\mathrm{\Gamma }_{bulk}510`$ left over from pre-shock magnetic acceleration. As described in the simulations of McKinney (2005b), patches of slightly faster or slower bulk $`\mathrm{\Gamma }`$ are present by $`r10^2r_g`$. In M87-based models, these range from $`2\mathrm{\Gamma }10`$.
The synchrotron emission angular frequency is
$$\omega _c\frac{3\mathrm{\Gamma }_e^2qB\mathrm{sin}\alpha }{2m_ec},$$
(54)
where $`\mathrm{sin}\alpha 1`$. This gives a characteristic synchrotron frequency of
$$\nu _c3\times 10^{12}\left(\frac{r}{r_g}\right)^{0.7}\mathrm{Hz}(r<390r_g,\mathrm{M87})$$
(55)
and
$$\nu _c3\times 10^{14}\left(\frac{r}{r_g}\right)^{1.5}\mathrm{Hz}(r<390r_g,\mathrm{M87}).$$
(56)
For $`r10^2r_g`$ where the fireball begins to form, this gives $`\nu _c100`$GHz (radio). By $`r10^3r_g`$, $`\nu _c10`$GHz (radio). The emission frequency depends on the mass accretion rate (and so $`\rho _{0,disk}`$) for any particular AGN. As discussed in McKinney (2005b), $`r10^2r_g`$ is also where the flow goes superfast (supersonic). Thus this is consistent with the idea that the radio-bright static knots at the base of the jet in, for example, Cen A is due to shocks in a transfast (transonic) transition (Hardcastle, 2005).
In summary, the jet in M87 likely emits most of the internal energy, generated in shocks in the transonic transition, as nonthermal synchrotron with some thermal synchrotron, such that the jet beyond $`10^310^4r_g`$ is relatively cold with $`2\mathrm{\Gamma }_{\mathrm{}}10`$.
Notice that in x-ray binaries, for example GRS 1915+105, have a jet with $`u/\rho _{0,e^{}e^+}c^22.5`$ and so thermal $`\mathrm{\Gamma }_e2.5`$. This gives that
$$t_{syn}10^7\mathrm{s}\left(\frac{r}{r_g}\right)^{1.4}(r<390r_g,\mathrm{GRS})$$
(57)
and
$$t_{syn}10^{11}\mathrm{s}\left(\frac{r}{r_g}\right)^3(r>390r_g,\mathrm{GRS}).$$
(58)
For a typical lab frame jet propagation time of $`t_{jet}r/c`$,
$$\frac{t_{syn}}{t_{jet}}10^3\left(\frac{r}{r_g}\right)^{0.4}(r<390r_g,\mathrm{GRS})$$
(59)
$$\frac{t_{syn}}{t_{jet}}10^7\left(\frac{r}{r_g}\right)^2(r>390r_g,\mathrm{GRS}),$$
(60)
where GRS denotes GRS1915+105. Thus thermal synchrotron is sufficiently fast to cool the jet. Since the jet is optically thick, as estimated above, then synchrotron self-absorption will dominate the emission process, which is what is observed (Foster et al., 1996; Fender & Belloni, 2004). The thermal synchrotron emission has
$$\nu _c2\times 10^{15}\left(\frac{r}{r_g}\right)^{0.7}\mathrm{Hz}(r<390r_g,\mathrm{GRS})$$
(61)
and
$$\nu _c3\times 10^{17}\left(\frac{r}{r_g}\right)^{1.5}\text{H}z(r>390r_g,\mathrm{GRS}).$$
(62)
Near the base this gives $`0.01`$keV emission (EUV). These soft synchrotron photons will be Compton upscattered (synchrotron self-Compton) by the $`\mathrm{\Gamma }5`$ jet to x-rays and contribute to the destruction of the Poynting-lepton jet. Like in AGN, nonthermal synchrotron likely takes away much of the shock-generated internal energy and this may account for some unidentified EGRET sources.
It is beyond the scope of the present study to establish whether nonthermal synchrotron, synchrotron self-Compton, or external Comptonization accounts for most of the high-energy luminosity.
#### 4.2.4 Other issues
Another possible way of contaminating the Poynting-dominated jet is by accreting a complicated field geometry and so baryon-loading the polar region. This turns the jet into a mixed lepton-baryon Poynting jet. Core-collapse presents the black hole with a field geometry that has an overall single poloidal sign. Compared to GRBs, AGN and x-ray binary black holes are more likely to accrete nontrivial field geometries leading to baryon contamination of the jet. Especially in Roche-lobe formed disks in x-ray binaries, it is likely that the accreted field geometry is quite tangled, so the likelihood of a relativistic Poynting jet is further reduced.
No black hole x-ray binary has been observed to have an ultrarelativistic jet (V4641 Sgr is still not confirmed, but see Chaty et al. 2003), despite the GRMHD physics in such systems being identical and the Lorentz factor is otherwise independent of the mass of the compact object. However, due to their relatively high radiative efficiency compared to AGN, x-ray binaries produce more $`\gamma `$-ray flux that increases the pair loading for a given magnetic field strength near the black hole. Also, the relatively high radiative efficiency means any Poynting-lepton jet is severely Compton dragged since the jet is optically thick. However, it is possible that there exists a large population of low radiative efficiency galactic black hole accretion systems. These radiatively inefficient systems would produce a large amount of Poynting flux per unit rest-mass flux which would be shock-converted by toroidal field instabilities into nonthermal synchrotron emission and could appear as โmicroblazars.โ However, thus-far observed x-ray binaries should not be as intrinsically luminous per rest-mass accretion rate since the Poynting flux per rest-mass flux available to shock-heating is two orders of magnitude smaller than available for AGN. This is due to the relatively high pair-loading in typical x-ray binaries. Low-luminosity x-ray binaries would behave more like blazars, and so low luminosity x-ray binary microblazars may account for some of the unidentified EGRET sources.
## 5. Relativistic Poynting-Baryon Jets
This section discusses how mildly relativistic Poynting-baryon jets can explain many jet observations. The origin of these jets is the inner-radial accretion disk. The origin of the mass is unstable convective outflows and magnetic buoyancy, and the mass fraction released is typically a few percent of the mass accretion rate (McKinney & Gammie, 2004; De Villiers et al., 2005a). The ratio of Poynting to baryon flux depends mostly on the spin of the black hole (Punsly & Coroniti, 1990a, b; McKinney & Gammie, 2004). Since the Poynting flux from a rapidly rotating black hole that is absorbed by the corona is also a few percent, the Poynting-baryon jet is heavily baryon-loaded. The heavy baryon-loading limits Poynting-baryon jets to only mildly relativistic velocities. The most relativistic, collimated, and least baryon-loaded portion of the Poynting-baryon jet is at the magnetic wall bounded by the Poynting-dominated jet.
### 5.1. Matter Jets and Outflows in AGN
Most AGN should have Poynting-baryon jets. This Poynting-baryon jet may often lead to erroneous conclusions about the nature of the jet in AGN systems.
For example, Junor et al. (1999); Biretta et al. (1999, 2002) suggest that M87 slowly collimates from about $`60^{}`$ near the black hole to $`10^{}`$ at large distances. However, two of their assumptions are likely too restrictive. First, they assumed the jet is always conical, which is apparent from figure 1 in Junor et al. (1999). If the jet is not conical this can overestimate the opening angle close to the core (i.e. perhaps $`35^{}`$ is reasonable all the way into the core). Second, their beam size is relatively large so that factors of $`2`$ error in the collimation angle are likely. Finally, and most importantly for this paper, they assumed that there is only one jet component. This likely leads to a poor interpretation of the observations. If there is a highly collimated relativistic Poynting-lepton jet surrounded by a weakly collimated Poynting-baryon jet, then this would also fit their observations.
Alternatively, if the accretion disk in M87 is a very thin SS-type disk with $`H/R0.00048`$ (McKinney, 2004), then their conclusion that there is slow collimation is plausible. However, thin disks may be much less efficient at producing jets (Livio, Ogilvie, & Pringle, 1999; Ghosh & Abramowicz, 1997) and may not be able to produce collimated jets (Okamoto, 1999, 2000). A form of the idea that winds collimate jets has also been proposed by Tsinganos & Bogovalov (2005) and applied to M87, but they consider a model where the wind slowly collimates the jet in order to fit observations. Here we suggest that the observations have been misinterpreted due to the presence of two components: a well-collimated relativistic cold Poynting-lepton jet and a mildly relativistic coronal outflow. We suggest the broader emission component is due to the coronal outflow.
Notice that more recent maps of the M87 jet-formation region show no โjet formationโ structure (Krichbaum et al., 2004). Thus, the structures seen previously may be transient features, such as associated with turbulent accretion disk producing a dynamic coronal outflow.
Measurements of the apparent jet speed in M87 reveal typically $`\mathrm{\Gamma }1.8`$ near the core while $`\mathrm{\Gamma }6`$ at larger radii. However, some core regions are associated with $`\mathrm{\Gamma }6`$ that rapidly fade (Biretta et al., 1999). This is consistent with a two-component outflow where the cold fast moving core of the jet is only observed if it interacts with the surrounding medium (or stars), the slower coronal outflow, or it undergoes internal shocks.
For relatively thin disks or slowly rotating black holes, Poynting-baryon jets could appear as โaborted jetsโ (Ghisellini et al., 2004).
The classical AGN unification models (Urry & Padovani, 1995) invoke a dominant role for the molecular torus and broad-line emitting clouds, while the broad coronal outflow may significantly contribute to modifications and in understanding the origin of the clouds (Elvis, 2000; Elvis et al., 2004).
Other erroneous conclusions could be drawn regarding the jet composition. Entrainment, which could occur at large distances when the ideal MHD approximation breaks down, causes difficulties in isolating the โproperโ jet componentโs composition. Worse is the fact that there should be two separate relativistic jet components, making it difficult to draw clear conclusions regarding the composition (Guilbert et al., 1983; Celotti & Fabian, 1993; Levinson & Blandford, 1996; Sikora & Madejski, 2000). However, is has been recently suggested that only electron-positron jets could explain FRII sources (Kino & Takahara, 2004).
It is also often assumed that if the jet is highly collimated that it is also highly relativistic near the black hole, which would suggest Comptonization of disk photons should produce clear spectral features (Sikora & Madejski, 2000). However, the jet may rather slowly accelerate and quickly collimate, which is universally what GRMHD numerical models find.
### 5.2. Jets in X-ray Binaries
The results of the previous section suggest that the term โmicroquasarโ does not accurately reflect the jet formation process. If radiatively efficient systems have no Poynting-lepton jet, then what produces their jets? Mildly relativistic jets from black hole microquasars may be produced by the inner-radial disk rather directly by the black hole. The above results suggest that GRS1915+105 may not have a Poynting-lepton jet during its quiescent accretion phase in the low-hard state. All black hole accretion systems with a thick disk have a mildly relativistic $`1\mathrm{\Gamma }3`$ coronal outflow due to convective instabilities and magnetic buoyancy (McKinney & Gammie, 2004). This component is sufficiently relativistic to explain the jets from black hole (and most neutron star) x-ray binary systems. This mechanism only requires a thick disk and not necessarily a spinning black hole, where other unification models suggest that the black hole spin is necessary (Meier, 2001).
It has been suggested that the transient, more relativistic, jet produced in GRS 1915+105 is the result the formation of a thin disk as the ADAF collapses as the mass accretion rate increases (Fender et al., 2004b). No particular model of the transient jet has been suggested.
Here we give a proposal for the disk-jet coupling in black hole x-ray binary systems, such as GRS 1915+105. In the prolonged hard x-ray state the disk is ADAF-like and the system produces a Poynting synchrotron self-absorbed jet with $`\mathrm{\Gamma }2`$, which may be partially or completely Compton dragged to nonrelativistic speeds. However, in the thick state, a Poynting-baryon jet is produced with $`\mathrm{\Gamma }1.5`$. During the soft state, the disk is SS-like (Shakura & Sunyaev, 1973) and the black hole polar field is relatively weak and the system generates an uncollimated (more radial) weak optically thin Poynting outflow. During this phase there is also a weak nonrelativistic uncollimated Poynting-baryon outflow.
During the transition between hard and soft x-ray states the production of pairs decreases significantly in the funnel, but the black hole polar magnetic field has yet to decay. During this transition, an optically thin Poynting-lepton jet with $`\mathrm{\Gamma }23`$ is produced that is collimated by the remaining inner-radial ADAF-like structure or the Poynting-baryon wind that was produced prior. The Lorentz factor produced in the transition depends on the details of the disk structure, and so $`\gamma `$-ray emission, during the transition. Once the black hole field has decayed, the fast transient jet shuts down.
Alternatively, during the transition to the high-soft state a transient jet can emerge as the corona is suddenly exposed to more Poynting flux from the black hole. This last bit of coronal material can be launched off as a faster transient baryon-loaded jet. The dynamics of the state transition is left for future work. This overall picture is in basic agreement with Fender & Belloni (2004), with the additional physics of pair creation dominating the Poynting-lepton jet formation process.
It is interesting that the results of Gierliลski & Done (2004) suggest that for at least some black hole x-ray binaries that have jets, the black hole is likely not rapidly rotating (i.e. perhaps $`j0.5`$). For such black holes, there is negligible Poynting flux in the form of a Poynting-dominated jet (McKinney & Gammie, 2004). Thus, our conclusion that black hole x-ray binary jets are driven by coronal outflows is consistent with the results of Gierliลski & Done (2004). However, even if black hole x-ray binaries were rapidly rotating they might not produce Poynting-dominated jets.
SS443 is plausibly an $`M20\mathrm{M}_{}`$ black hole system that has a jet with $`v0.3c`$ (Lopez et al., 2003). Such a low jet velocity can be explained by a Poynting-baryon jet. To explain the opening angle of $`1^{}`$ the disk should be very thick near the black hole, while the pulsed jet features can be explained as an instability due to the overly thick disk self-interacting at the poles near the black hole.
## 6. Summary of Companion Paper Numerical Results
This section summarizes the results of McKinney (2005b) using a GRMHD code HARM (Gammie et al., 2003a) with an advanced inversion method (Noble et al., 2005). Kerr-Schild coordinates were used in order to avoid numerical artifacts associated with causal interactions between the inner-radial boundary and the rest of the flow. Viscous models have found this issue to be critical to avoid spurious fluctuations in the jet (McKinney & Gammie, 2002), such as might be associated with codes using Boyer-Lindquist coordinates.
### 6.1. Jet propagation
As described in detail in McKinney (2005b) and as shown in figure 3 and 4, the Poynting-dominated jet forms as the differential rotation of the disk and the frame-dragging of the black hole induce a significant toroidal field that launches material away from the black hole by the same force described in equation C14.
A coronal outflow is also generated between the disk and Poynting-dominated jet. In this model the coronal outflow has $`\mathrm{\Gamma }_{\mathrm{}}1.5`$. The coronal-funnel boundary contains shocks with a sonic Mach number of $`M_s100`$. The inner-radial interface between the disk and corona is a site of vigorous reconnection due to the magnetic buoyancy and convective instabilities present there. These two parts of the corona are about $`100`$ times hotter than the bulk of the disk. Thus these coronal components are a likely sites for Comptonization and nonthermal particle acceleration.
Figure 3 and figure 4 show the final log of density and magnetic field projected on the Cartesian z vs. x plane. For the purposes of properly visualizing the accretion flow and jet, we follow MacFadyen & Woosley (1999) and show both the negative and positive $`x`$-region by duplicating the axisymmetric result across the vertical axis. Color represents $`\mathrm{log}(\rho _0/\rho _{0,disk})`$ with dark red highest and dark blue lowest. The final state has a density maximum of $`\rho _02\rho _{0,disk}`$ and a minimum of $`\rho _010^{13}\rho _{0,disk}`$ at large radii. Grid zones are not smoothed to show grid structure. Outer radial zones are large, but outer $`\theta `$ zones are below the resolution of the figure.
Clearly the jet has pummelled its way through the surrounding medium, which corresponds to the stellar envelope in the collapsar model. By the end of the simulation, the field has been self-consistently launched in to the funnel region and has a regular geometry there. In the disk and at the surface of the disk the field is curved on the scale of the disk scale height. Within $`r10^2r_g`$ the funnel field is ordered and stable due to the poloidal field dominance. However, beyond $`r10^2r_g`$ the poloidal field is relatively weak compared to the toroidal field and the field lines bend and oscillate erratically due to pinch instabilities. The radial scale of the oscillations is $`10^2r_g`$ (but up to $`10^3r_g`$ and as small as $`10r_g`$), where $`r10r_g`$ is the radius where poloidal and toroidal field strengths are equal. By the end of the simulation, the jet has only fully evolved to a state independent of the initial conditions at $`r5\times 10^3r_g`$, beyond which the jet features are a result of the tail-end of the initial launch of the field. The head of the jet has passed beyond the outer boundary of $`r=10^4r_g`$. Notice that the magnetic field near the black hole is in an X-configuration. This is due to the BZ-effect having power $`P_{jet}\mathrm{sin}^2\theta `$, which vanishes at the polar axis. The X-configuration is also related to the fact that the disk+corona is collimating the Poynting-dominated jet. The field is mostly monopolar near the black hole, and such field geometries decollimate for rapidly rotating black holes in force-free electrodynamics (Krasnopolsky et al., 2005).
### 6.2. Summary of Fits
A summary of the fits along a fiducial field line is given. Near the black hole the half-opening angle of the full Poynting-dominated jet is $`\theta _j1.0`$, while by $`r120r_g`$, $`\theta _j0.1`$. This can be roughly fit by
$$\theta _j\left(\frac{r}{r_g}\right)^{0.4}(\mathrm{inner})$$
(63)
for $`r<120r_g`$ and $`\theta _j0.14`$ beyond. The core of the jet follows a slightly stronger collimation with
$$\theta _jr^{2/5}$$
(64)
up to $`r<120r_g`$ and $`\theta _j0.09`$ beyond. Also, roughly for M87 and the collapsar model, the core of the jet has
$$\mathrm{\Gamma }_{bulk}\left(\frac{r}{5r_g}\right)^{0.44}(\mathrm{inner})$$
(65)
for $`5<r10^3r_g`$ and constant beyond for the M87 model if including synchrotron radiation, while the collapsar model should continue accelerating and the power law will truncate when most of the internal and Poynting energy is lost to kinetic energy and the jet becomes optically thin at about $`r10^9r_g`$ or internal shocks take the energy away. If the acceleration is purely thermal without any magnetic effect, then $`\mathrm{\Gamma }r`$ (Mรฉszรกros & Rees, 1997). However, it is not clear how the equipartition magnetic field affects the acceleration. Roughly for GRS 1915+105 the core of the jet has
$$\mathrm{\Gamma }_{bulk}\left(\frac{r}{5r_g}\right)^{0.14}(\mathrm{inner})$$
(66)
for $`5<r10^3r_g`$ and constant beyond, with no account for Compton drag or pair annihilation. Also, for any jet system the base of the jet has $`\rho _0r^{0.9}(\mathrm{inner})`$ for $`r120r_g`$ and $`\rho _0r^{2.2}(\mathrm{outer})`$ beyond. For the collapsar and M87 models
$$\frac{\rho _0}{\rho _{0,disk}}1.5\times 10^9\left(\frac{r}{120r_g}\right)^{0.9}(\mathrm{inner})$$
(67)
and
$$\frac{\rho _0}{\rho _{0,disk}}1.5\times 10^9\left(\frac{r}{120r_g}\right)^{2.2}(\mathrm{outer}),$$
(68)
while for GRS 1915+105 the inner-radial coefficient is $`10^5`$ and outer is $`6\times 10^3`$. For the collapsar model, the inner radial internal energy density is moderately fit by
$$\frac{u}{\rho _{0,disk}c^2}=4.5\times 10^9\left(\frac{r}{120r_g}\right)^{1.8}(\mathrm{inner}).$$
(69)
The outer radial internal energy density is moderately fit by
$$\frac{u}{\rho _{0,disk}c^2}=4.5\times 10^9\left(\frac{r}{120r_g}\right)^{1.3}(\mathrm{outer}).$$
(70)
The transition radius is $`r120r_g`$. For M87 the internal energy is near the rest-mass density times $`c^2`$ until $`r120r_g`$ when the dependence is as for the collapsar case. For GRS1915+105 the internal energy is near the rest-mass density times $`c^2`$ until $`r120r_g`$ and then rises to about $`2.5`$ times the rest-mass density times $`c^2`$. The inner radial toroidal lab field is well fit by
$$\frac{B^{\widehat{\varphi }}}{\sqrt{\rho _{0,disk}c^2}}[\mathrm{Gauss}]=0.0023\left(\frac{r}{390r_g}\right)^{0.7}(\mathrm{inner})$$
(71)
for $`5<r<390r_g`$. The outer radial toroidal lab field is well fit by
$$\frac{B^{\widehat{\varphi }}}{\sqrt{\rho _{0,disk}c^2}}[\mathrm{Gauss}]=0.0023\left(\frac{r}{390r_g}\right)^{1.5}(\mathrm{outer})$$
(72)
for $`r>390r_g`$.
For the typical jet with no atypical pinch instabilities, the energy and velocity structure of the jet follow
$$ฯต(\theta )=ฯต_0e^{\theta ^2/2\theta _0^2},$$
(73)
where $`ฯต_00.18`$ and $`\theta _08^{}`$. The total luminosity per pole is $`L_j0.023\dot{M}_0c^2`$, where $`10\%`$ of that is in the โcoreโ peak Lorentz factor region of the jet within a half-opening angle of $`5^{}`$. Also, $`\mathrm{\Gamma }_{\mathrm{}}`$ is approximately Gaussian
$$\mathrm{\Gamma }_{\mathrm{}}(\theta )=\mathrm{\Gamma }_{\mathrm{},0}e^{\theta ^2/2\theta _0^2},$$
(74)
where $`\mathrm{\Gamma }_{\mathrm{},0}3\times 10^3`$ and $`\theta _04.3^{}`$. Also, $`\mathrm{\Gamma }`$ is approximately Gaussian
$$\mathrm{\Gamma }(\theta )=\mathrm{\Gamma }_0e^{\theta ^2/2\theta _0^2},$$
(75)
where $`\mathrm{\Gamma }_05`$ and $`\theta _011^{}`$. The outer sheathโs ($`\theta 0.2`$) seed photon temperature as a function of radius is
$$T_{\gamma ,seed}50\mathrm{k}\mathrm{e}\mathrm{V}\left(\frac{r}{5\times 10^3r_g}\right)^{1/3}.$$
(76)
## 7. Discussion
For GRB jets, the injected Poynting flux probably dominates the injected energy flux provided by neutrino annihilation. This poses problems for the classic neutrino-driven fireball model. Super-efficient neutrino emission mechanisms with an average neutrino energy of $`210`$MeV are required in order for the neutrino annihilation energy to be as large as the energy provided by the BZ effect. However, the BZ effect itself might operate in a super-efficient mode once flux has accumulated near the black hole (Narayan et al., 2003). This vertical field threading the disk leads up to $`5`$ times larger luminosity (McKinney & Gammie, 2004), in which case an average neutrino energy of $`1000`$MeV is required to compete with the BZ-effect.
For GRBs, equation 42 shows that slightly less rapidly rotating black holes would generate slightly less Lorentz factors but achieve a lower luminosity. This is consistent with the observation that harder long-duration bursts have higher luminosity, and so suggests that small changes in the stellar rotation rate might determine the hardness of long-duration bursts.
The fact that blazars are less luminous for increasing hardness could be explained by the dependence on the jet Lorentz factor on the pair creation physics. Blazars could have similar black hole spin, but the hardness of their emission is determined by the jet Lorentz factor. Lower luminosity systems load the jet with less pairs and so the Lorentz factor is larger. Compton drag of environment or disk reflected seed photons can then upscatter to very large energy, such as observed in TeV-emitting BL-Lac objects.
Our results suggest that radiatively efficient x-ray binaries, such as GRS1915+105, may only exhibit a relativistic Poynting-baryon jet. In particular, such a jet is relativistic only in the low-hard state when the disk is geometrically thick.
## 8. Conclusions
Primarily two types of relativistic jets form in black hole (and perhaps neutron star) systems. The Poynting-dominated jet region is composed of field lines that connect the rotating black hole to large distances. Since the ideal MHD approximation holds very well, the only matter that can cross the field lines are neutral particles, such as neutrinos, photons, and free neutrons.
The primary differences between GRBs, AGN, and black hole x-ray binaries is the pair-loading of the Poynting-dominated jet, a similar mass-loading by free neutrons in GRB-type systems, the optical depth of the jet, and the synchrotron cooling timescale of the jet.
For GRB-type systems the neutron diffusion flux is sufficiently large to be dynamically important, but small enough to allow $`\mathrm{\Gamma }1001000`$. Beyond $`r10r_g`$ many of the electron-positron pairs annihilate, so the Poynting-dominated jet is dominated in mass by electron-proton pairs from collision-induced neutron decay. Most of the energy is provided by the BZ effect instead of neutrino-annihilation.
For AGN and x-ray binaries, the density of electron-positron pairs established near the black hole primarily determines the Lorentz factor at large distances. Radiatively inefficient AGN, such as M87, achieve $`2\mathrm{\Gamma }_{\mathrm{}}10`$ and are synchrotron cooling limited. The lower the $`\gamma `$-ray radiative efficiency of the disk, the more energy per particle is available in the shock-zone. Radiatively efficient systems such as GRS1915+105 likely have no Poynting-lepton jet due to strong pair-loading and destruction by Comptonization by the plentiful soft photons for x-ray binaries with optically thick jets. However, all these systems have a mildly relativistic, baryon-loaded jet when in the hard-low state when the disk is geometrically thick, which can explain jets in most x-ray binary systems.
In an companion paper McKinney (2005b), a GRMHD code, HARM, with pair creation physics was used to evolve many black hole accretion disk models. The basic theoretical predictions made in this paper that determine the Lorentz factor of the jet were numerically confirmed. However, Poynting flux is not necessarily directly converted into kinetic energy, but rather Poynting flux is first converted into enthalpy flux into a โmagnetic fireballโ due to shock heating. Thus, at large distances the acceleration is primarily thermal, but most of that thermal energy is provided by shock-conversion of magnetic energy. In GRB systems this magnetic fireball leads to thermal acceleration over an extended radial range. The jets in AGN and x-ray binaries release this energy as synchrotron and inverse Compton emission and so the jet undergoes negligible thermal acceleration beyond $`r10^210^3r_g`$.
Based upon prior numerical (McKinney, 2005b) and this theoretical work, basic conclusions for collapsars include:
1. Black hole energy, not neutrino energy, typically powers GRBs.
2. Poynting-dominated jets are mostly loaded by $`e^{}e^+`$ pairs close to the black hole, and by $`e^{}p`$ pairs for $`r10r_g`$.
3. BZ-power and neutron diffusion primarily determines Lorentz factor.
4. Variability is due to toroidal field instabilities.
5. Poynting flux is converted into enthalpy flux and leads to the formation of a โmagnetic fireball.โ
6. Patchy jet develops $`10^2\mathrm{\Gamma }_{\mathrm{}}10^3`$, as required by internal shock model.
7. Random number of patches ($`<1000`$ for 30 second burst) and so random number of pulses.
8. Energy structure of jet is Gaussian with $`\theta _08^{}`$.
9. Core of jet with $`\theta _j5^{}`$ can explain GRBs.
10. Extended slower jet component with $`\theta _j25^{}`$ can explain x-ray flashes.
11. Coronal outflows with $`\mathrm{\Gamma }1.5`$ may power supernovae (by producing, e.g., $`{}_{}{}^{56}\mathrm{Ni}`$) with $`M0.1\mathrm{M}_{}`$ processed by corona.
Based upon prior numerical (McKinney, 2005b) and this theoretical work, basic conclusions for AGN or x-ray binaries include:
1. Poynting-dominated jets $`e^{}e^+`$ pair-loaded unless advect complicated field.
2. $`\gamma `$-ray radiative efficiency, and so pair-loading, determines maximum possible Lorentz factor.
3. Poynting-lepton jet is collimated with $`\theta _j5^{}`$.
4. Extended slow jet component with $`\theta _j25^{}`$.
5. For fixed accretion rate, variability is due to toroidal field instabilities.
6. Poynting flux is shock-converted into enthalpy flux.
7. In some AGN, shock heat in transonic transition lost to synchrotron emission and limits achievable Lorentz factor to $`2\mathrm{\Gamma }10`$ (e.g. in M87).
8. Coronal outflows produce broad inner-radial jet features in AGN together with well-collimated jet component (e.g. in M87).
9. In some x-ray binaries, Compton drag loads Poynting-lepton jets and limits Poynting-lepton jet to $`\mathrm{\Gamma }2`$ or jet destroyed.
10. In some x-ray binaries, Poynting-lepton jet optically thick and emits self-absorbed synchrotron.
11. Coronal outflows have collimated edge with $`\mathrm{\Gamma }1.5`$.
12. Coronal outflows may explain all mildly relativistic and nonrelativistic jets in radiatively efficient systems (most x-ray binaries).
For AGN and X-ray binaries, the coronal outflow collimation angle is strongly determined by the disk thickness. The above assumed $`H/R0.2`$ near the black hole and $`H/R0.6`$ far from the black hole, while $`H/R0.9`$ (ADAF-like) is perhaps more appropriate for some systems. The sensitivity of these results to $`H/R`$ is left for future work.
## Acknowledgments
I thank Avery Broderick for an uncountable number of inspiring conversations. I also thank Charles Gammie, Brian Punsly, Amir Levinson, and Ramesh Narayan, with whom each I have had inspiring conversations. This research was supported by NASA-ATP grant NAG-10780 and an ITC fellowship.
## Appendix A Pair Creation Notes
The electron-positron pair plasma that forms may annihilate itself into a fireball if the pair annihilation rate is faster than the typical rate of the jet ($`c^3/GM`$) near the black hole. Also, if the pair annihilation timescale is shorter than the dynamical time, then pair annihilation would give a collisional term in the Boltzmann equation. From the pair annihilation rate given by equation D7, one finds that $`t_{pa}GM/c^3`$ for AGN and marginally so for x-ray binaries. Thus, pairs mostly do not annihilate, and so formally the pair plasma that forms in the low-density funnel region is collisionless so that the Boltzmann equation should be solved directly. Plasma instabilities and relativistic collisionless shocks are implicitly assumed to keep the pairs in thermal equilibrium so the fluid approximation remains mostly valid, as is a good approximation for the solar wind (see, e.g. Feldman & Marsch 1997; Usmanov et al. 2000). This same approximation has to be invoked for the thick disk state in AGN and x-ray binaries, such as for the ADAF model (McKinney, 2004). For regions that pair produce slower than the jet dynamical time, each pair-filled fluid element has a temperature distribution that gives an equation of state with $`P=\rho _{0,e^{}e^+}k_bT_e/m_e`$ rather than $`P=(11/12)aT^4`$, where $`a`$ is the radiation constant. So most of the particles have a Lorentz factor of $`\mathrm{\Gamma }_eu/(\rho _{0,e^{}e^+}c^2)`$ and little of the internal energy injected is put into radiation. This also allow the use of a single-component approximation. A self-consistent Boltzmann transport solution is left for future work.
On the contrary for GRB systems, due to the relatively high density of pairs, the time scale for pair annihilation is $`t_{pa}GM/c^3`$ along the entire length of the jet. Thus a pair fireball forms and the appropriate equation of state is that of an electron-positron-radiation fireball. Thus, formally the pair fireball rest-mass density is not independent of the pair fireball internal energy density. However, because the pairs are well-coupled to the radiation until a much larger radius of $`r10^810^{10}r_L`$, the radiation provides an inertial drag on the remaining pair plasma. That is, the relativistic fluid energy-momentum equation is still accurate. So the effective rest-mass density is $`\rho _0+u`$ ($`u`$ the total internal energy of the fireball), and so the effective rest-mass is independent of the cooling of the fireball until the fireball is optically thin (see, e.g., Mรฉszรกros & Rees 1997).
For GRB systems, the mass conservation equation is reasonably accurate. Even though the electron-positron pairs annihilate, the rest-mass of pairs injected is approximately that of the pairs that are injected due to Fick-diffusion of neutrons (see next section). The annihilation energy from electron-positron pairs contributes a negligible additional amount of internal energy, so can be neglected, especially compared to the Poynting energy flux that emerges from the black hole. Thus, the rest-mass can always be assumed to be due to baryons rather than the electron-positron pairs. This also suggests that the neutrino annihilation is a negligible effect if the BZ power is larger than the neutrino annihilation power.
In summary, the rest-mass evolution discussed in McKinney (2005b) are accurate for GRB, AGN, and marginally so for x-ray binaries. This is despite the lack of Boltzmann transport for the collisional system, or a collisional term due to pair annihilation.
### A.1. Baryon Contamination
Notice that some fraction of baryons contaminate the jet due to neutron diffusion and subsequent collisional cascade into an electron-proton plasma (Levinson & Eichler, 2003). They estimate the diffusion using Fickโs law. First, the role of ambipolar diffusion is considered (see, e.g. Shu 1992, chpt. 27). The drift velocity is
$$v_{drift,pn}\frac{B^2}{4\pi \gamma _{pn}n_pm_pn_nm_nL},$$
(A1)
where $`Lr(H/R)`$ is the typical field radius of curvature induced by disk turbulence and
$$\gamma _{pn}=\frac{\sigma _{pn}v_{rel,pn}}{m_p+\rho /m_n}$$
(A2)
is the drag coefficient and $`\sigma v40c\times 10^{27}`$. The drift velocity can also be written as
$$v_{drift,pn}/c\left(\frac{b^2}{\rho c^2}\right)\left(\frac{m_nc}{\sigma v\rho r(H/R)}\right).$$
(A3)
Assuming all the diffused neutrons are converted to protons+electrons and carried with the outflow, then the diffusion flux is
$$F=\rho v_{drift,pn}=\frac{b^2}{\rho c^2}\frac{m_nc^2}{\sigma vr(H/R)}.$$
(A4)
and the mass flux across an axisymmetric conical outflow is
$$\dot{M}_{inj,ambi}=2\pi _0^{r_{out}}Fr๐r$$
(A5)
and so
$$\dot{M}_{inj,ambi}=2\pi (Fr)r_{out}$$
(A6)
GRB numerical GRMHD models show that the coronal region next to the Poynting-dominated jet has a time-averaged value of $`b^2/(\rho c^2)1`$ and the turbulent induced eddies occur when the disk has $`H/R0.2`$ (McKinney & Gammie, 2004; McKinney, 2004). This gives a mass flux vs. radius of
$$\dot{M}_{inj,ambi}10^{14}\left(\frac{rr_{stag}}{r_g}\right)\dot{M}_{acc},$$
(A7)
where pairs that enter inside $`rr_{stag}`$ fall into the black hole and do not load the jet. For a recombination radius of $`r_n2\times 10^9`$cm (Levinson & Eichler, 2003), this gives
$$\dot{M}_{inj,ambi}4\times 10^{10}\dot{M}_{acc}.$$
(A8)
This can be compared to the result of Levinson & Eichler (2003) for Fick-diffusion, where the mass injection rate of free-streaming particles (their eq. 7) is
$$\dot{M}_{inj,Fick}3\times 10^7\left(\frac{r}{r_g}\right)^{2/3}\dot{M}_{acc}$$
(A9)
for their collapsar model with jet half-opening angle of $`\theta _j=0.1`$, specific enthalpy $`h1`$, $`L_j10^{51}\mathrm{erg}\mathrm{s}^1`$, a neutron thermal velocity of $`v0.1c`$. The thermal velocity is based upon the near-funnel coronal value of $`u/(\rho _0c^2)0.010.1`$ as measured from GRMHD numerical models. Notice that this ratio is typically $`0.01`$ in the corona, but is $`0.1`$ at the edge, so we use $`0.1`$ since the Fick diffusion is based upon the edge values. Notice that they used $`vc`$. For a recombination radius of $`r_n2\times 10^9`$cm after which no more free neutrons exist, one has that
$$\dot{M}_{inj,Fick}7\times 10^5\dot{M}_{acc}.$$
(A10)
Hence, Fick diffusion dominates ambipolar diffusion. The role of reconnection, between the corona/coronal wind and jet, in loading the jet with baryons is left for future work.
The characteristic timescale for moving these pairs is $`t_g(r/r_g)`$ and characteristic length is $`r_g(r/r_g)`$, so a characteristic density vs. radius is
$$\rho _{pe^{}}\frac{\dot{M}_{inj,Fick}t_g}{r_g^3}\left(\frac{r}{r_g}\right)^23\times 10^7\left(\frac{r}{r_g}\right)^{4/3}\rho _{0,disk}.$$
(A11)
Notice that this is comparable to the rest-mass in pairs given by equation 18. Thus, as the fireball decays in pair rest-mass, the rest-mass quickly becomes dominated by neutrons diffusing across the magnetic wall between the corona and funnel. Hence, the baryon conservation law holds and the approximations used here hold well. A pair-annihilation term is only needed to account for the contribution to the internal energy. Since $`f_\rho 8f_h`$, this contribution is a $`10\%`$ effect and is not expected to affect the results of the numerical models of McKinney (2005b).
### A.2. Pair-Radiation Equation of State
The total amount of comoving energy put into the thermal fireball is
$`u_{0,tot}=\rho _{0,e^{}e^+}c^2+u_{0,e^{}e^+}+u_{0,\gamma }=`$
$`AT^4{\displaystyle _0^{\mathrm{}}}๐x{\displaystyle \frac{x^2\sqrt{x^2+\stackrel{~}{m}^2}}{e^{\sqrt{x^2+\stackrel{~}{m}^2}}+1}}+u_{0,\gamma },`$ (A12)
where $`u_{0,\gamma }=1.62348AT^4`$, $`A=4.66244\times 10^{15}\mathrm{erg}\mathrm{cm}^3\mathrm{K}^4`$, $`xpc/k_bT`$, $`\stackrel{~}{m}m_ec^2/k_bT`$, $`p`$ is the momentum in the fluid frame, and the rest-number density of photons is $`n_\gamma =20.2944\mathrm{cm}^3T^3`$. The rest-mass in pairs is
$$\rho _{0,e^{}e^+}=BT^3_0^{\mathrm{}}๐x\frac{x^2}{e^{\sqrt{x^2+\stackrel{~}{m}^2}}+1},$$
(A13)
where $`B=3.07589\times 10^{26}\mathrm{g}\mathrm{cm}^3\mathrm{K}^3`$, and the number density of pairs is $`n_{0,e^{}e^+}=\rho _{0,e^{}e^+}/m_e`$. Notice that the GRMHD equations of motion relate the comoving energy to energy at infinity by
$$u_{tot}=u_{0,tot}u^tu_t+p_{gas}=(f_\rho +f_h+f_m)e_{\nu \overline{\nu },ann}$$
(A14)
and $`\rho _{e^{}e^+}=\rho _{0,e^{}e^+}u^t`$, where
$$p_{gas}=\frac{AT^4}{3}_0^{\mathrm{}}๐x\frac{x^2\sqrt{x^2+\stackrel{~}{m}^2}}{e^{\sqrt{x^2+\stackrel{~}{m}^2}}+1}+p_{0,\gamma },$$
(A15)
and $`p_\gamma =u_{0,\gamma }/3`$. This gives sufficient information to solve for $`u^t`$ and $`u_t`$ or $`\dot{\rho }_{e^{}e^+}=/t(\rho _{0,e^{}e^+}u^t)`$. The study of McKinney (2005b) uses a fixed $`\gamma `$-law gas equation of state with $`\gamma =4/3`$ to model the typically radiation-dominated system.
## Appendix B Ideal MHD Quantities Conserved along each Flow Line
Kerr spacetime is stationary and axisymmetric with 2 Killing vectors $`\xi _t^\mu =\frac{}{_t}=\delta _t^\mu `$ and $`\xi _\mu ^\varphi =\frac{}{\varphi }=\delta _\varphi ^\mu `$ that satisfy $`_\xi (๐ )=0`$, where $``$ is the Lie derivative and $`๐ `$ is the metric. For a vector $`X^\mu `$, tensors $`๐`$ that obey $`_X(๐)=0`$ are conserved along $`X`$. In particular, for $`X=\xi `$, such a tensor is a physical quantity independent of the ignorable coordinates $`t`$ and $`\varphi `$. For $`X=u`$, the 4-velocity, the tensor is conserved along each flow line. One can derive a set of conserved flow quantities that are associated with the Killing symmetries (Bekenstein & Oron, 1978). The below summarizes those results that are key to this paper. This presentation is necessary for the discussion regarding the determination of the Lorentz factor of the jet.
In the ideal MHD approximation one can show that $`_u(\xi ^\mu A_\mu )=0`$, where $`F_{\mu \nu }A_{\nu ,\mu }A_{\mu ,\nu }`$ defines the vector potential $`A_\mu `$. Thus for an unsteady axisymmetric flow the $`\varphi `$ component of the magnetic vector potential ($`A_\varphi =\xi _\varphi ^\mu A_\mu `$) is conserved along each flow line, while for a steady non-axisymmetric flow the electric potential ($`A_t=\xi _t^\mu A_\mu `$) is conserved along each flow line. Bekenstein & Oron (1978) also show that there are four other independent conserved scalar quantities for an axisymmetric, stationary fluid. These correspond to the first integrals of Maxwellโs equation $`{}_{}{}^{^{}}F_{;\mu }^{\mu \nu }=0`$ and the conservation equation $`T_{\nu _;\mu }^\mu =0`$ with the use of the continuity equation $`(\rho _0u^\mu )_{;\mu }=0`$. These first integrals are associated with the projection of the 2 Killing vectors and the magnetic field ($`b^\mu `$) on these equations of motion.
For a degenerate ($`{}_{}{}^{^{}}F_{}^{\mu \nu }F_{\mu \nu }=e^\mu b_\mu =0`$), stationary and axisymmetric plasma, one finds
$$A_{\varphi ,\theta }A_{t,r}A_{t,\theta }A_{\varphi ,r}=0.$$
(B1)
It follows that one may write
$$\frac{A_{t,\theta }}{A_{\varphi ,\theta }}=\frac{A_{t,r}}{A_{\varphi ,r}}\mathrm{\Omega }_F(r,\theta )$$
(B2)
where $`\mathrm{\Omega }_F`$ is usually interpreted as the โrotation frequencyโ of the electromagnetic field (this is Ferraroโs law of isorotation; see e.g. Frank, King, & Raine 2002, ยง9.7 in a nonrelativistic context). Notice that $`\mathrm{\Omega }_FF_{tr}/F_{r\varphi }=F_{t\theta }/F_{\theta \varphi }`$. One can show in the ideal MHD approximation that
$$\mathrm{\Omega }_F=v^\varphi B^\varphi v^\theta /B^\theta =v^\varphi B^\varphi v^r/B^r$$
(B3)
is conserved along each flow line. The first term corresponds to fluid rotation and the second term corresponds to the slip along the toroidal component of a field line. This yields $`F_{\mu \nu }`$ in terms of the free functions $`\mathrm{\Omega }_F,A_\varphi `$, and $`B^\varphi `$. Thus, $`F_{rt}\sqrt{g}E_r=\mathrm{\Omega }_FA_{\varphi ,r}`$, $`F_{\theta t}\sqrt{g}E_\theta =\mathrm{\Omega }_FA_{\varphi ,\theta }`$, $`F_{r\theta }=\sqrt{g}B^\varphi `$, $`F_{\varphi r}=\sqrt{g}B^\theta =A_{\varphi ,r}`$, and $`F_{\theta \varphi }=\sqrt{g}B^r=A_{\varphi ,\theta }`$. The diagonal components are zero and $`F_{\varphi t}E_\varphi =0`$, where $`B^i{}_{}{}^{^{}}F_{}^{it}`$, $`E_iF_{it}/\sqrt{g}`$ such that $`B^iE_i={}_{}{}^{^{}}F_{}^{\mu \nu }F_{\mu \nu }/(4\sqrt{g})EB`$. Thus for fixed poloidal magnetic field, $`\mathrm{\Omega }_F`$ is a measure of the electric field. With the Faraday written in terms of $`B^i`$ and $`\mathrm{\Omega }_F`$, the electromagnetic field automatically satisfies the source-free Maxwell equations.
Using similar constraints on $`F^{\mu \nu }=A^{\nu ,\mu }A^{\mu ,\nu }`$ and with $`E^iF^{it}\sqrt{g}`$, $`B_i{}_{}{}^{^{}}F_{it}^{}`$, one can show that $`\sqrt{g}F^{rt}E^r=\tau _\theta B_\theta `$, $`\sqrt{g}F^{\theta t}E^\theta =\tau _rB_r`$, $`\sqrt{g}F^{\varphi t}E^\varphi =\tau _\varphi B_\varphi `$, $`\sqrt{g}F^{r\theta }=B_\varphi `$, and $`\sqrt{g}F^{r\varphi }=B_\theta `$, $`\sqrt{g}F^{\theta \varphi }=B_r`$. Here there is only one independent quantity among the three $`\tau _i`$โs that are set by $`{}_{}{}^{^{}}F_{\mu \nu }^{}F^{\mu \nu }=0`$ and that the flow and metric are stationary and axisymmetric (i.e. $`E_\varphi =0`$). One can show that $`\tau _\varphi =B_rB_\theta (\tau _r\tau _\theta )/B_\varphi ^2`$ and solve for another by using $`E_\varphi =F^{\alpha \beta }g_{\alpha t}g_{\beta \varphi }=0`$.
It is interesting to note that in Boyer-Lindquist coordinates $`\tau _\varphi =0`$ and $`\tau _r=\tau _\theta \tau `$ and then the contravariant and covariant Faraday take on the same simple form with
$$\tau =\frac{g^{t\varphi }g^{tt}\mathrm{\Omega }_F}{g^{\varphi \varphi }g^{t\varphi }\mathrm{\Omega }_F}=\frac{\mathrm{\Omega }_F/\mathrm{\Omega }_{ZAMO}1}{\mathrm{\Omega }_0\mathrm{\Omega }_F},$$
(B4)
where $`\mathrm{\Omega }_{ZAMO}g^{t\varphi }/g^{tt}=2ar/A`$ is the angular frequency of a zero angular momentum observer (ZAMO) and
$$\mathrm{\Omega }_0g^{\varphi \varphi }/g^{t\varphi }=\frac{a^2\mathrm{\Delta }/\mathrm{sin}^2\theta }{2ar}.$$
(B5)
Notice that if $`\mathrm{\Omega }_F=\mathrm{\Omega }_{ZAMO}`$, then $`\tau =0`$ and so $`E^i=0`$. So the difference between the dragging of inertial frames and the field rotation frequency generates the electric field $`E^i`$. Also in Boyer-Lindquist $`B_r=B^r|g|g^{\theta \theta }g^{\varphi \varphi }(\mathrm{\Omega }_F/\mathrm{\Omega }_01)`$, where $`|g|g^{\theta \theta }g^{\varphi \varphi }=(\mathrm{\Sigma }2r)/\mathrm{\Delta }`$. Also, $`B_\theta =B^\theta |g|g^{rr}g^{\varphi \varphi }(\mathrm{\Omega }_F/\mathrm{\Omega }_01)`$, where $`|g|g^{rr}g^{\varphi \varphi }=\mathrm{\Sigma }2r`$. Also, $`B_\varphi =B^\varphi |g|g^{rr}g^{\theta \theta }`$, where $`|g||\mathrm{Det}(g_{\mu \nu })|`$ and $`|g|g^{rr}g^{\theta \theta }=\mathrm{\Delta }\mathrm{sin}^2\theta `$.
Using the definition of $`F^{\mu \nu }`$ given above one can define $`A_\varphi `$ in terms of the poloidal $`B^r`$ and $`B^\theta `$ giving
$$A_\varphi (l_f)A_\varphi (l_i)=_{l_i}^{l_f}(\sqrt{g}(B^rd\theta B^\theta dr))$$
(B6)
over the line segment from $`l_i`$ to $`l_f`$. Thus given the poloidal field components, then $`A_\varphi `$ can be determined up to a constant. The contours of constant $`A_\varphi `$ represent time-dependent poloidal magnetic field surfaces for any $`\varphi `$. Shown in Cartesian coordinates, beyond a few gravitational radii from even a $`j=1`$ black hole, the density of lines represents the field strength in the lab frame. Near the horizon where the intrinsic volume of space is larger than in Minkowski space-time, the density of field lines in such a Cartesian plot overestimates the lab frame field strength by factors of $`2`$.
For an inviscid fluid flow of magnetized plasma, the energy and angular momentum flux per unit rest-mass flux
$$E=T_t^r/(\rho _0u^r)=T_t^\theta /(\rho _0u^\theta )$$
(B7)
and
$$L=T_\varphi ^r/(\rho _0u^r)=T_\varphi ^\theta /(\rho _0u^\theta ),$$
(B8)
respectively, are conserved along each flow line. For unmagnetized flows $`E`$ is conserved for any stationary flow, while $`L`$ is conserved for any axisymmetric flow. If the ideal MHD approximation ($`e^\mu =0`$) holds, then the magnetic flux per unit rest-mass flux
$$\mathrm{\Phi }=\frac{B^r}{\rho _0u^r}=\frac{B^\theta }{\rho _0u^\theta }=\frac{B^\varphi }{\rho _0(u^\varphi u^t\mathrm{\Omega }_F)}$$
(B9)
is conserved along each flow line. This also implies that $`(\mathrm{\Phi }+b^t)/\rho _0=b^r/v^r=b^\theta /v^\theta `$. Bekenstein & Oron (1978) also show that
$$\mathrm{\Psi }=E+\mathrm{\Omega }_FL=h(u_t+\mathrm{\Omega }_Fu_\varphi ),$$
(B10)
is conserved along each flow line. One can reduce $`E`$ and $`L`$ to forms such as
$$E=\zeta u_t\mathrm{\Psi }\mathrm{\Phi }b_t/h=hu_t+\mathrm{\Phi }\mathrm{\Omega }_FB_\varphi $$
(B11)
$$L=\zeta u_\varphi +\mathrm{\Psi }\mathrm{\Phi }b_\varphi /h=hu_\varphi +\mathrm{\Phi }B_\varphi $$
(B12)
where $`h=(\rho _0+u_g+p)/\rho _0`$ is the gas specific enthalpy and $`\zeta =h+b^2/\rho _0`$ is the total specific enthalpy. Notice that for an isentropic flow that $`d(\rho _0+u)/d\rho _0=h`$ and so
$$dp/(h\rho _0)=dh/h.$$
(B13)
Clearly any ratio of these conserved flow quantities is also conserved along each flow line (e.g. the energy flux per unit magnetic flux ($`E/\mathrm{\Phi }`$)). Any axisymmetric, stationary flow solution can be written in terms of the 6 independent quantities $`A_\varphi `$, $`B^\varphi `$, $`\mathrm{\Omega }_F`$, $`\mathrm{\Phi }`$, $`E`$, and $`L`$, where the single function $`A_\varphi `$ determines the dependent quantities $`B^r`$ and $`B^\theta `$. Entropy is a dependently conserved quantity when one writes the rest-mass density and enthalpy in terms of the entropy and another conserved quantity. In general, the solution is determined once the conserved flow quantities are set by the boundary conditions and the other quantities are set by the equations of motion, either directly or using the Grad-Shafranov approach (for a review see, e.g., Levinson 2005). The limitations of, and hence extensions to, the ideal MHD approximation are described in Meier (2004).
Note that for an axisymmetric stationary degenerate fluid the Boyer-Lindquist coordinate components are related to the Kerr-Schild coordinates by $`B^r[\mathrm{BL}]=B^r[\mathrm{KS}]`$ , $`B^\theta [\mathrm{BL}]=B^\theta [\mathrm{KS}]`$,$`B^\varphi [\mathrm{BL}]=B^\varphi [\mathrm{KS}]B^r[\mathrm{KS}](a2r\mathrm{\Omega }_F)/\mathrm{\Delta }`$, while for a general fluid $`B_r[\mathrm{BL}]=B_r[\mathrm{KS}]+(aB_\varphi )/\mathrm{\Delta }`$, $`B_\theta [\mathrm{BL}]=B_\theta [\mathrm{KS}]`$, and $`B_\varphi [\mathrm{BL}]=B_\varphi [\mathrm{KS}]`$.
## Appendix C Ideal MHD Fluid Forces
First, to investigate the spatial acceleration of the fluid along a flow line $`A_\varphi `$ one requires a unit length space-like vector that satisfies $`V^\mu A_{\varphi ,\mu }=0`$ since $`A_\varphi `$ is conserved along a flow line. The unit-length magnetic field satisfies these properties since $`B^\mu A_{\varphi ,\mu }=0`$. To study the poloidal acceleration along a field line, the toroidal magnetic field component it projected out to obtain the unit vector
$$B_p^{\widehat{\mu }}=NB_p^\mu =N(B^\mu \frac{(\omega _\nu ^\varphi B^\nu )\xi _\varphi ^\mu }{|\omega ^\varphi ||\xi _\varphi |})=N(B^\mu B^\varphi \delta _\varphi ^\mu )$$
(C1)
where $`N=1/\sqrt{B_p^\mu B_p^\nu g_{\mu \nu }}=1/\sqrt{B^iB^jg_{ij}}`$ with $`\{i,j\}=\{r,\theta \}`$ only, and $`\omega ^\varphi `$ is the $`\varphi `$ basis one-form. Therefore the projection of an acceleration along a poloidal projection of each flow line is
$$a_{A_\varphi }=a_\mu B_p^{\widehat{\mu }}=a_rB^{\widehat{r}}+a_\theta B^{\widehat{\theta }}$$
(C2)
Second, to investigate the spatial collimation of the fluid, a unit length space-like vector that is perpendicular to the poloidal field line and perpendicular to the $`\varphi `$-direction (i.e. $`\xi _\varphi `$) is required. This vector is
$$C^{\widehat{\mu }}=N_CC^\mu =N_Cฯต_{\alpha \beta }^\mu B_p^{\widehat{\alpha }}\xi _\varphi ^\beta ,$$
(C3)
where $`N_C=1/\sqrt{C^\mu C^\nu g_{\mu \nu }}`$ and $`ฯต_{}^{\mu }{}_{\alpha \beta }{}^{}`$ is the spatial permutation tensor. Thus $`C^{\widehat{r}}=NN_CB^\theta `$ and $`C^{\widehat{\theta }}=NN_CB^r`$. Therefore the projection of an acceleration in the collimation direction is
$$a_{coll}=a_\mu C^{\widehat{\mu }}=N_C(a_rB^{\widehat{\theta }}a_\theta B^{\widehat{r}}).$$
(C4)
Notice that for $`\theta <\pi /2`$ if $`a_{coll}<0`$, then the flow is collimating toward the polar axis. For $`\theta >\pi /2`$ if $`a_{coll}>0`$, then the flow is collimating toward the polar axis.
Third, there are many interesting frames to measure the acceleration. The acceleration away from geodesics is obtained from the projection of $`P^{\mu \nu }=g^{\mu \nu }+u^\mu u^\nu `$ on $`_\gamma T_\mu ^\gamma =0`$, giving Eulerโs equations for the deviation from geodesic motion
$$\rho _0ha_\mu ^G=P_{\mu }^{}{}_{}{}^{\alpha }p_{,\alpha }+J^\alpha F_{\alpha \mu }$$
(C5)
where $`a_G^\mu =u_{;\nu }^\mu u^\nu =u^\alpha (u_{,\alpha }^\mu +\mathrm{\Gamma }_{}^{\mu }{}_{\beta \alpha }{}^{}u^\beta )`$ is the โgeodesic accelerationโ away from the geodesic motion ($`a_G^\mu =0`$). Note that $`\mathrm{\Gamma }`$ here is the connection coefficient and not the Lorentz factor. This comoving geodesic acceleration โhidesโ the effect of gravity on the fluid. One could instead focus on the coordinate acceleration $`a_C^\mu =u^\nu u_{,\nu }^\mu `$, which represents the change in the 4-velocity in the momentarily comoving frame. From the geodesic equation of motion,
$$a_C^\mu =a_G^\mu \mathrm{\Gamma }_{}^{\mu }{}_{\alpha \beta }{}^{}u^\alpha u^\beta ,$$
(C6)
and so the coordinate acceleration along a flow line is
$$a_{A_\varphi }^C=a_\mu ^CB_p^{\widehat{\mu }}=a_{A_\varphi }^GB^{\widehat{\mu }}u^\alpha u^\beta \mathrm{\Gamma }_{\mu \alpha \beta },$$
(C7)
so the geodesic deviation and gravitational acceleration along each flow line can be studied separately.
Notice that for a stationary flow $`a_{A_\varphi }^C=0`$ where the poloidal velocity $`u^p=0`$. For a black hole with field lines that tie the black hole to large radii, there must exist a region where $`u^p=0`$. For an outflow to reach large distances, the region where $`u^p=0`$ separates those fluid elements that eventually fall into the black hole and those fluid elements that reach large distances. Levinson (2005) refers to this as the โinjection surface.โ If particles were created only due to reaching the Goldreich-Julian charge density, then this must be the location where particles emerge. For an injection region with negligible angular velocity $`u^\varphi 0`$, then
$$a_{A_\varphi }^GB^{\widehat{r}}\mathrm{\Gamma }_{rtt}+B^{\widehat{\theta }}\mathrm{\Gamma }_{\theta tt}$$
(C8)
determines the location of the stagnation surface. In Boyer-Lindquist coordinates
$$\mathrm{\Gamma }_{rtt}=\frac{r^2(a\mathrm{cos}\theta )^2}{\mathrm{\Sigma }^2}$$
(C9)
and
$$\mathrm{\Gamma }_{\theta tt}=\frac{2ra^2\mathrm{cos}\theta \mathrm{sin}\theta }{\mathrm{\Sigma }^2}.$$
(C10)
See section 3 for a discussion of injection physics.
### C.1. Forces in Lab Frame
The forces as written in the lab frame, rather than comoving frame, allow for simple understanding of the force dynamics. Equations C5 and C2 imply that
$$a_{A_\varphi }^G=NB^j(a_{j}^{G}{}_{}{}^{(MA)}+a_{j}^{G}{}_{}{}^{(EM)}),$$
(C11)
where for a stationary, axisymmetric flow
$$a_{j}^{G}{}_{}{}^{(MA)}=\frac{u_ju^ip_{,i}+p_{,j}}{\rho _0h}$$
(C12)
is the hydrodynamic acceleration. For an isentropic flow, equation B13 implies that
$$a_{j}^{G}{}_{}{}^{(MA)}=P_{j}^{}{}_{}{}^{i}(\mathrm{log}h)_{,i}=u_ju^i(\mathrm{log}h)_{,i}+(\mathrm{log}h)_{,j}.$$
(C13)
For an axisymmetric, stationary, ideal MHD fluid the electromagnetic acceleration $`a_{j}^{G}{}_{}{}^{(EM)}=J^\alpha F_{\alpha j}/(\rho _0h)`$ reduces to simply
$$a_{j}^{G}{}_{}{}^{(EM)}=\frac{B^\varphi B_{\varphi ,j}}{\rho _0h}.$$
(C14)
Thus the magnetic acceleration is due to the gradient of the toroidal magnetic field along a field line. Notice that from equation B11 or B12 that in the limit $`\rho _0h0`$ that $`B_\varphi E/(\mathrm{\Phi }\mathrm{\Omega }_F)=L/\mathrm{\Phi }`$, a conserved flow quantity, thus $`B^jB_{\varphi ,j}=0`$ implying the electromagnetic field is force-free. However, notice that the acceleration $`a_{j}^{G}{}_{}{}^{(EM)}`$ then becomes undefined.
Equations C5 and C4 imply that
$$a_{coll}^G=NN_CB^j(a_{k}^{G}{}_{}{}^{(MA)}+a_{k}^{G}{}_{}{}^{(EM)})ฯต_{}^{k}{}_{j}{}^{},$$
(C15)
where $`ฯต_{}^{k}{}_{j}{}^{}`$ is the poloidal permutation tensor and $`a_{k}^{G}{}_{}{}^{(MA)}`$ is given in equations C12 or C13. This hydrodynamic collimation is due to the pressure acceleration in the comoving frame along a field line but directed to collimate the flow. For an axisymmetric, stationary, ideal MHD fluid the electromagnetic collimation acceleration ($`C^{\widehat{j}}J^\alpha F_{\alpha j}/(\rho _0h)=C^{\widehat{j}}a_{j}^{G}{}_{}{}^{(EM)}`$) reduces to
$$a_{coll}^{G}{}_{}{}^{(EM)}=\frac{NN_C}{\rho _0h}\left(ฯต^{ab}(B^p)^2f[B_a]+ฯต_{a}^{}{}_{}{}^{b}B^aB^\varphi B_{\varphi ,b}\right),$$
(C16)
where $`f[B_a]\left(B_{a,b}\mathrm{\Omega }_F(\tau _aB_a)_{,b}\right)`$ and $`(B^p)^2(B^r)^2+(B^\theta )^2`$ and $`ฯต_{a}^{}{}_{}{}^{b}`$ is the poloidal permutation tensor. The last term on the right hand side of equation C16 represents the โhoop-stressโ that leads to collimation for nonrelativistic winds. The first two terms on the right hand side correspond, respectively, to the forces due to poloidal magnetic stresses and the electric field ($`E^i`$) gradients. The latter can collimate relativistic outflows.
## Appendix D Compton Scattering
In the lab frame, seed photons are Compton upscattered if the energy of the photon $`E_{seed}\mathrm{\Gamma }_em_ec^2`$ for an electron Lorentz $`\mathrm{\Gamma }_e`$. The upscattering continues until the lab frame photon energy exceeds electron energy. In the lab frame, each scatter gives the photon a new energy $`E_{scat}4\mathrm{\Gamma }_e^2E_{seed}`$ if $`E_{seed}<m_ec^2/\mathrm{\Gamma }_e`$ and $`E_{scat}\mathrm{\Gamma }_em_ec^2`$ otherwise. Two limiting scenarios are if the photon crosses the jet or if the photon travels parallel and within the jet.
The optical depth to Compton scattering for a photon in the rest frame of the jet electrons is
$$\tau =_l^{}\left(\frac{\rho _{0,e^{}e^+}}{m_e}\right)\sigma _T๐l^{},$$
(D1)
which in the lab frame gives
$$\tau =_l\left(\frac{\rho _{0,e^{}e^+}}{m_e}\right)\sigma _T\mathrm{\Gamma }_e(1\beta \mathrm{cos}\theta )๐l,$$
(D2)
where $`\beta =v/c`$, $`\mathrm{\Gamma }_e=(1\beta ^2)^{1/2}`$, and $`\theta =\pi /2`$ corresponds to perpendicular interactions and $`\theta =0`$ to parallel. Across the jet
$$\tau _{}=_{\theta _j}^{\theta _j}\left(\frac{\rho _{0,e^+e^{}}}{m_e}\right)\sigma _T\mathrm{\Gamma }_er๐\theta ,$$
(D3)
while for along the jet
$$\tau _{}=_{r_0}^{\mathrm{}}\left(\frac{\rho _{0,e^{}e^+}}{m_e}\right)\sigma _T\mathrm{\Gamma }_e(1\beta )๐r,$$
(D4)
where $`\beta =\left(\frac{\mathrm{\Gamma }_e^21}{\mathrm{\Gamma }_e^2}\right)^{1/2}`$. For $`\mathrm{\Gamma }_e1`$, $`\mathrm{\Gamma }_e(1\beta )(2\mathrm{\Gamma }_e)^1`$, so
$$\tau _{}_{r_0}^{\mathrm{}}\left(\frac{\rho _{0,e^+e^{}}}{m_e}\right)\left(\frac{\sigma _T}{2\mathrm{\Gamma }_e}\right)๐r,$$
(D5)
where $`\sigma _T`$ is the Thomson scattering cross-section and $`r_0`$ is some fiducial starting radius in the jet (see Rybicki & Lightman 1979; Longair 1992).
Similar calculations can be used to estimate the forward or backwards pair-production or pair-annihilation optical depths. In the lab frame, for a photon gas moving with $`\mathrm{\Gamma }`$ with one photon energy $`E=\mathrm{\Gamma }E^{}`$ and another $`>(\mathrm{\Gamma }m_ec^2)^2/E`$, then the photons annihilate with a cross section at $`E0.5`$MeV of $`\sigma \sigma _T`$ and the Klein-Nishina corrections for photons with $`E0.5`$MeV modify the cross section such that $`\sigma E^1`$. For a spectrum (number per unit time per unit area per unit energy) $`fE^\alpha `$, then for $`\alpha =2`$, the average cross section is $`\sigma 0.06\sigma _T`$. In this case the relevant proper density is $`n_\gamma `$ when for a photon beam, where $`n_\gamma `$ is the number density of (typically fewer) high energy photons. See also Lithwick & Sari (2001).
For a beam of electrons with velocity $`\beta `$, $`\sigma \sigma _T/\beta `$ for nonrelativistic electrons and for $`\mathrm{\Gamma }1`$
$$\sigma \frac{3\sigma _T}{8\mathrm{\Gamma }}(\mathrm{log}2\mathrm{\Gamma }1).$$
(D6)
The pair annihilation rate is
$$t_{pa}^1\sigma v\left(\frac{\rho _{0,e^{}e^+}}{m_e}\right),$$
(D7)
which can be compared to some dynamical time to determine if pair annihilation is important. |
warning/0506/cond-mat0506326.html | ar5iv | text | # Self-propelled running droplets on solid substrates driven by chemical reactions
## I Introduction
The movement of droplets in external gradients fascinates scientists and layman alike at least since Newtonโs description Newt1730hab2 of Hauksbeeโs experiment with drops of orange oil that move between two non-parallel glass plates towards the point of smallest plate distance Hauk1710 . In another example a drop of liquid freely immersed in another liquid subject to a temperature gradient will move towards the higher temperature region due to Marangoni forces caused by surface tension gradients Vela98 . A drop sitting on a solid substrate also moves in a temperature Broc89 or wettability Gree78 ; Raph88 ; ChWh92 gradient. Especially the Marangoni force is already used to manipulate droplets, for example in light induced drop movement ION00 . Similar concepts of a directed movement of small amounts of soft matter in a given gradient are also realized in models of cell motility JJP03 .
However, even more intricate are situations where matter spontaneously starts a directed movement in initially homogeneous settings. Small pieces of camphor that move on a liquid surface by emitting a surfactant have also been studied for centuries Vent1799 ; Toml1869 ; Rayl1890b ; Haya02 . More recently oil droplets containing volatile additives and interacting droplets of different volatile oils have been reported to move on solid surfaces due to the solutal Marangoni effect caused by evaporation/condensation CMS64 . Also intricate is the spontaneous movement of a juxtaposed pair of droplets with different wetting properties, so called bi-slugs, along a capillary tube BiQu00 . This effect was already mentioned by Marangoni Mara1871 . Another example are drops immersed in a second liquid. If the drops contain a soluble surfactant undergoing an isothermal chemical reaction at their surface the drops may start to move RRV94 . Apart from chemical reactions, drop movement can also be driven by surface phase transitions Rieg03 ; YoPi05\_pre . The movement is possible because such active drops change their surrounding and produce a gradient that drives their motion.
Recent experiments also found chemically driven running droplets on solid substrates BBM94 ; DoOn95 ; LeLa00 ; LKL02 ; SMHY05 ; Sumi05pre ; Mage03 . In these cases small droplets of solution are put on partially wettable substrates. A droplet changes the substrate by adsorption or desorption of a solute rendering the substrate underneath the droplet less wettable than the bare substrate. The radial symmetry still assures an equilibrium position that is, however, unstable. As a result, fluctuations break the symmetry and the drop starts to move in a self-sustained manner. In the course of its movement it changes the substrate and leaves a less wettable trail behind (see Fig. 1). We distinguish two types of experiments: In case I the substrate is changed in an irreversible way whereas in case II it recovers its initial state after the droplet has passed. The droplet does not rest again and follows a random trajectory until trapping itself between self-produced less wettable patches (type I) DoOn95 ; LKL02 . Alternatively, a running drop may perform a periodic movement (type II) SMHY05 ; Sumi05pre . The droplet often moves with nearly constant velocity for a rather long time (tens of minutes). However, in all cases the motion ceases when the overall system reaches thermal equilibrium. Next, we shortly recall the specific results of different sets of experiments, yielding running droplets.
In Ref. DoOn95 droplets of $`n`$-alkanes ($`n`$-octane and $`n`$-dodecane) are used that contain 1,1,2,2-tetrahydroperfluorodecyltrichlorosilane (CF<sub>3</sub>(CF<sub>2</sub>)<sub>7</sub>(CH<sub>2</sub>)<sub>2</sub>SiCl<sub>3</sub>). The silane molecules form dense crafted monolayers on silicon or glass and render these surfaces hydrophobic. The droplet motion is stalled when no further hydrophilic surface is available. No limitation due to silane depletion inside the droplet is observed. For millimeter sized drops (smaller than the capillary length) droplet velocities between 1 mm/s and 10 cm/s are observed, depending on liquid viscosity, silane concentration and droplet size. The velocity increases with the silane concentration and the droplet size. The reaction rate is estimated to be 1100 s<sup>-1</sup>mol<sup>-1</sup>. Typical reaction times are given as 0.01โ0.2 s.
Another experiment is performed in a chemically different system LeLa00 ; LKL02 using millimeter-sized droplets of a nonpolar solution of $`n`$-alkylamine (1 mM solution of C<sub>6</sub>NH<sub>2</sub> or C<sub>18</sub>NH<sub>2</sub>) in decahydro-naphthalene ($`\gamma =41`$ mJ/m<sup>2</sup>, $`\mu =0.001`$ Ns/m<sup>2</sup>). Therein, silicon substrates with microprinted high-energy surfaces are employed, that expose a dense packing of carboxylic acid functionalities (CO<sub>2</sub>H). The amines dissolved in the droplet adsorb at the substrate and produce a surface of lower energy exposing methyl groups. The effect of different adsorbates LeLa00 and reaction kinetics LKL02 is investigated. The velocity decreases with increasing droplet size.
The only example of a type II experiment SMHY05 features droplets of oil (5 mM iodine solution of nitrobenzene saturated with potassium iodide) on glass substrates immersed in aqueous solution of a cationic surfactant (1 mM stearyl trimethyl ammonium chloride). In this system the stearyl trimethyl ammonium ion (STA) absorbs at the glass substrate outside the oil droplet and renders the glass lyophilic. When the oil droplet moves on top of the coating the STA desorbs into the oil. In this way a wettability contrast between front and back of the moving droplet is created and sustained. In contrast to the type I case the substrate recovers its lyophilic state soon after the droplet has passed, because โnewโ STA absorbs at the glass. The aqueous phase can be seen as an infinite reservoir of adsorbent. Movement only finally stops when the oil droplet is saturated with STA.
Similar phenomena can be seen in metallurgic systems where droplets of liquid metals or alloys react with the metallic substrate, for instance, by alloying. The layer between the droplet and the substrate may be less wettable than the bare substrate resulting in the migration of reactive islands. This was studied for tin islands on copper surfaces that move and leave tracks of bronze behind SBH00 . Contrariwise, the layer can be more wettable than the bare substrate. This is typical for the related process of reactive spreading (also called reactive wetting) where a droplet of liquid on a (nearly) non-wettable surface starts to spread after forming a more wettable layer underneath LaEu96 ; Yost98 ; WBR98 ; VMHE99 ; SCT00 ; WeGr02 . Variants are possible in which during spreading the substrate becomes less wettable in the center of the drop ZWT98 . However, reactive spreading processes do normally not result in running droplets (but see the โsuddenly displacedโ droplets in Ref. KRE99 ).
For type I experiments an implicit equation for the velocity $`v`$ of the droplet was derived Brde95 ; deGe98 from a simple theoretical argument. Based on a balance of friction force and driving capillary forces one obtains $`v=C\mathrm{tan}\theta ^{}(1\mathrm{exp}(rL/v))`$, where $`r`$ is the reaction rate, $`L`$ the size of the droplet and $`C`$ a constant. The dynamic contact angles at the advancing and receding ends of the droplet are then assumed to be identical ($`\theta ^{}`$), i.e. the droplet profile is approximated by a spherical cap with $`\theta ^{}`$ given by $`\mathrm{cos}\theta ^{}=(\mathrm{cos}\theta _e^a+\mathrm{cos}\theta _e^r)/2`$. The static contact angles at the advancing edge $`\theta _e^a`$ and at the receding edge $`\theta _e^r>\theta _e^a`$ are different due to the chemical gradient. The expression for the velocity is found from a first order reaction on the substrate that yields chemical concentrations $`\alpha _a=0`$ and $`\alpha _r=1\mathrm{exp}(rL/v)`$ at the respective ends of the droplet. The expression for the velocity predicts a monotone increase of the droplet velocity with the droplet length $`L`$ and the reaction rate $`r`$ in line with experimental observations DoOn95 . The droplet velocity in the limiting case of a saturated chemical reaction was also given in MiMe97 . However, the experiments of Ref. LKL02 show the opposite trend; the velocity decreases with increasing drop sizes and effective reaction rate. In the framework of Ref. Brde95 the decrease is explained as a result of the flattening of the drops by gravity. Related works discuss running droplets in a random medium deGe99 and the forced wetting of a plate immersed into a reactive fluid deGe97 .
In this paper we propose and analyze dynamical models for self-propelled running droplets for both, type I and type II experiments. Our models consist of coupled evolution equations for the film thickness profile and the substrate coverage. Thereby the wettability of the substrate is modeled by a coverage dependent disjoining pressure. The different types of experiments are taken into account by incorporating different reaction terms. The models are capable of reproducing the different experimentally found regimes. A short account of a variant of the model for the type I experiments, i.e. for the case of a irreversibly changed substrate was recently presented TJB04 . There it was found that the dynamic contact angles at the advancing and receding ends of the droplet are not identical but rather resemble the corresponding static ones. Even for a strong driving force the deviation between dynamic and static angles is only about 10 percent. The bifurcation leading from sitting to running droplets for a finite reaction rate was identified as a drift-pitchfork bifurcation. Here we analyze the type I model in more detail and extend our analysis to the type II model.
In section II the type I and type II dynamical models are presented. Results for the two models are discussed in section III and IV, respectively. Stationary running droplets and sitting droplets are characterized in dependence of the control parameters reaction rate, droplet volume, diffusion constant of the adsorbate field and desorption rate of the adsorbate (for type II model only). The resulting families of solutions and their linear stability are used to derive phase diagrams that describe the existence regions for running and sitting droplets. To further elucidate the drift-pitchfork bifurcation that mediates the transition from sitting to running droplets also in the type II model we analyze the linear modes that destabilize sitting droplets close to the bifurcation. Before we conclude in section V, we use numerical simulations to illustrate and analyze periodic droplet motion that is possible in the type II model. The simulations are compared with our continuation results and the type II experiment of Ref. SMHY05 .
## II Model
### II.1 Evolution equations
The evolution of a thin liquid film on a horizontal smooth solid substrate is described by an equation for the thickness profile $`h(x,y,t)`$ derived from the Navier-Stokes equations using long-wave or lubrication approximation ODB97
$$_th=\left\{\frac{h^3}{3\eta }p\right\}.$$
(1)
The parameters $`\gamma `$ and $`\eta `$ are the surface tension and viscosity of the liquid, respectively. They define the viscous time scale $`\tau _v=\gamma L/\eta `$, where $`L`$ is a typical length for the system. The change in time of the film thickness profile equals the gradient of a flow that results as the product of a mobility and a pressure gradient. The mobility $`h^3/3\eta `$ corresponds to a parabolic velocity profile
$$u(x,z)=\left(\frac{z^2}{2}zh\right)p$$
(2)
within the film. The velocity component orthogonal to the substrate $`w`$ can be obtained using the continuity equation
$$_xu+_zw=\mathrm{\hspace{0.17em}0}.$$
(3)
The pressure
$$p=\gamma \mathrm{\Delta }h+\mathrm{\Pi }(h,\varphi )$$
(4)
contains the curvature (or Laplace) pressure $`\gamma \mathrm{\Delta }h`$ and the disjoining pressure $`\mathrm{\Pi }(h,\varphi )`$. The latter comprises effective molecular interactions between the film surface and the substrate and accounts for the wetting properties of the substrate deGe85 ; Hunt92 ; Isra92 . As discussed in detail below, here a mathematically simple function $`\mathrm{\Pi }(h,\varphi )`$ is used that is common in the literature. In many situations the qualitative outcome only depends on very general characteristics of the disjoining pressure TNPV02 . The used form allows for solutions of Eq. (1) that represent static (i.e. sitting) droplets with a finite mesoscopic equilibrium contact angle. The disjoining pressure is chosen such that the droplets coexist with an ultrathin precursor film.
The evolution of the density of the adsorbed layer on the substrate determines the wettability and is modeled by a reaction-diffusion equation for the dimensionless field $`0<\varphi (x,y,t)<1`$
$$_t\varphi =R(h,\varphi )+d^{}\mathrm{\Delta }\varphi ,$$
(5)
where the function $`R(h,\alpha )`$ describes adsorption or desorption on the substrate. The second term allows for a diffusion of the chemical species along the substrate. For simplicity we assume that the adsorbate has the same diffusion constant on the bare and the droplet covered substrate. However, assuming different constants would not change the presented results qualitatively. Note that for type I experiments $`\varphi `$ directly corresponds to the substrate coverage. However, in a type II experiment the droplet dissolves a more wettable coating. In this case the coverage corresponds to $`1\varphi `$. Also here one can use Eq. (5), however, the signs of adsorption and desorption term are switched. With this convention, in both cases the wettability decreases with increasing $`\varphi `$.
One can neglect the dynamics of the concentration field of the chemical in the bulk of the droplet by assuming a fast equilibration of the solute concentration within the moving droplet as compared to the reaction at the substrate. The fast equilibration is caused by diffusion and convective motion within the droplet whereas the latter is driven by the lateral movement of the droplet along the substrate. It corresponds to the limit of a small Damkรถhler number Da$`=rc_0/(D/L)`$ giving the ratio of reaction velocity at the substrate and diffusion velocity in the droplet Prob94 . The parameter $`r`$ is a typical reaction rate, $`c_0`$ stands for a typical concentration of the chemical species in the droplet and $`D`$ is the diffusion constant in the droplet.
### II.2 Reaction term
Corresponding to the two different sets of experiments that the model shall describe we use two different reaction terms. For type I the initial coverage of the substrate is zero (i.e. $`\varphi =0`$) and we only allow for adsorption underneath the droplet
$$R_1(h,\varphi )=r_{in}\xi (h)\left(1\varphi \right).$$
(6)
The reaction saturates at $`\varphi =1`$ because this is the maximal possible coverage. The function $`\xi (h)`$ represents a (smooth) step function that approaches one and zero inside and outside the droplet, respectively. For type II the โrestโ state of the substrate without any droplet is the fully covered substrate (again corresponding to $`\varphi =0`$, see above section II.1) and we allow for desorption underneath the droplet and adsorption outside the droplet
$$R_2(h,\varphi )=r_{in}\xi (h)(1\varphi )r_{out}[1\xi (h)]\varphi .$$
(7)
The time scales of the reactions at the substrate inside and outside the droplet are defined by the effective rate constants $`r_{in}`$ and $`r_{out}`$, respectively. Note that they have the dimension s<sup>-1</sup>. The function $`\xi (h)`$ may be the step function
$$\xi _1=\mathrm{\Theta }(hh_c),$$
(8)
or the smooth function
$$\xi _2=\{\mathrm{tanh}[(hh_c)/\mathrm{\Delta }]+1\}/2.$$
(9)
The value of $`h_c`$ is chosen slightly larger than the thickness of the precursor film and the maximal drop height is always $`h_{max}h_c`$. A small value of $`\mathrm{\Delta }h_{max}`$ ensures that the switch between a predominant adsorption reaction and a predominant desorption reaction occurs over a small film thickness range. Note, that $`\xi _2\xi _1`$ for $`\mathrm{\Delta }0`$. Changes in the details of the reaction term do not affect the results qualitatively.
### II.3 Disjoining pressure
For the disjoining pressure we use throughout the present work
$$\mathrm{\Pi }(h)=\frac{2S_ld_0^2}{h^3}+\left(1\frac{\varphi }{g}\right)\frac{5S_sd_0^5}{h^6}$$
(10)
where $`d_0=0.158`$ nm is the Born repulsion length that defines a lower cut-off for the film thickness. The parameter $`S_l`$ and $`S_s`$ are the long- and short-range components of the total spreading coefficient $`S=S_l+S_s`$ (for $`\varphi =0`$). We use $`S_l<0`$ and $`S_s>0`$ corresponding to a destabilizing long-range van der Waals and a stabilizing short-range interaction disjpress . For $`\varphi =0`$ the pressure allows for drops sitting on a stable precursor film as can be seen studying the corresponding densities of the excess surface energy $`f`$, related to the disjoining pressures by $`\mathrm{\Pi }=_hf`$.
The short-range interaction contains the influence of the coating and crucially influences the static contact angle Shar93 . To account for the varying wettability caused by the different substrate coverage we let the short-range part of the spreading coefficient $`S_s`$ depend linearly on the field $`\varphi (x,y,t)`$. The signs are chosen in a way that for both types of experiments $`g>0`$ assures that an increase in $`\varphi `$ corresponds to a lower wettability, i.e. to a larger equilibrium contact angle $`\theta _e`$ given by $`\mathrm{cos}\theta _e=S/\gamma +1`$ Shar93 . The constant $`g`$ relates the coverage to the wettability and therefore defines the magnitude of the possible wettability gradient. Note, that our Eq. (10) corresponds to the linear relation between $`\mathrm{cos}\theta _e`$ and $`\varphi `$ assumed in Ref. DoOn95 ; LKL02 ; Brde95 .
### II.4 Dimensionless equations
We rewrite Eqs. (1) to (10) by introducing scales $`3\gamma \eta /l\kappa ^2`$, $`\sqrt{\gamma l/\kappa }`$, and $`l`$ for $`t`$, $`(x,y)`$, and $`h`$, respectively. Then one obtains $`2|S_l|d_0^2/l^3`$ for the scaled spreading coefficient $`\kappa `$. As length scale $`l`$ we use the value of the film thickness where the local free energy $`f(h)`$ has its minimum, i.e., where $`\mathrm{\Pi }(h)=0`$. This gives $`l=(5S_s/2|S_l|)^{1/3}d_0`$ and implies that the ratio of the strength of the two antagonistic effective molecular interactions is only an implicit parameter of the system. The length $`l`$ also corresponds to the thickness of the precursor film for a sitting drop on the bare substrate.
Defining the dimensionless overall reaction rate $`r=3r_{in}\gamma \eta /l\kappa ^2`$, diffusion constant $`d=3d^{}\eta /\kappa l^2`$, and ratio of reaction rates $`s=r_{in}/r_{out}`$, we obtain from Eqs. (1) to (10) the dimensionless coupled evolution equations for the thickness profile $`h`$ and the field $`\varphi `$
$`_th`$ $`=`$ $`\left\{h^3\left[\mathrm{\Delta }h+\mathrm{\Pi }(h,\varphi )\right]\right\}`$ (11)
$`_t\varphi `$ $`=`$ $`rR(h,\varphi )+d\mathrm{\Delta }\varphi `$ (12)
with the disjoining pressure
$$\mathrm{\Pi }(h,\varphi )=\frac{1}{h^3}+\left(1\frac{\varphi }{g}\right)\frac{1}{h^6}$$
(13)
and the options for the reaction term
$`R_1(h,\varphi )`$ $`=`$ $`\xi (h)\left(1\varphi \right)`$ (14)
$`R_2(h,\varphi )`$ $`=`$ $`\xi (h)(1\varphi )s[1\xi (h)]\varphi .`$ (15)
To give an impression of the influence of the coverage we give in Fig. 2 droplet profiles for different values of the coverage $`\varphi `$. Thereby, $`\varphi `$ is assumed to be independent of time and constant along the substrate. Increasing the ratio $`\varphi /g`$ clearly leads to increasing contact angles and decreasing droplet length.
In the following we explore the full non-linear behavior of the dynamical models presented above. The main interest lies in an exploration of the existence regions of qualitatively different solutions depending on control parameters like reaction rate $`r`$, ratio of desorption and adsorption rates $`s`$, and droplet volume $`V`$. It was already shown that a type I model allows for three-dimensional running droplets that move with constant velocity $`v`$ and droplet shape TJB04 . Here, we restrict our attention to two-dimensional drops to be able to explore a large part of the parameter space. To determine such droplets we therefore use Eqs. (11,12) replacing $``$ by $`_x`$. Then continuation techniques DKK91 ; DKK91b ; AUTO97 are employed to calculate two-dimensional running droplets moving with constant speed. This is achieved by switching to the comoving frame $`xvt`$ and imposing appropriate boundary conditions. In the comoving frame running droplets correspond to steady solutions. Integration in time of the full system is also used in some cases. Details on all numerical techniques used can be found in the Appendix.
## III Results for model I (only adsorption)
We start with model type I, where adsorption takes place only underneath the droplet, i.e. we combine the evolution equations for the film thickness profile (11) and the field of adsorbate coverage (12) with the reaction term (14) and the disjoining pressure (13).
### III.1 Thickness and coverage profiles
First we focus on the characterization of the changing solution behavior in dependence of the overall reaction rate $`r`$. Without diffusion along the substrate ($`d=0`$) one finds unstable sitting droplets and stable running droplets for all $`r`$. The droplets and the coverage profile move with constant velocity $`v`$ and constant shape. We emphasize that this state corresponds to a subtle dynamical equilibrium given that coverage profile and droplet velocity depend on each other.
Fig. 3 shows for two different reaction rates profiles of moving droplets \[(a) and (b)\] where the streamlines correspond to contour lines of the stream function
$$\psi (x,z)=\left(\frac{z^3}{2}\frac{3z^2h}{2}\right)_xpvz,$$
(16)
and indicate the flow in the comoving coordinate system. Note that the velocity fields are obtained by $`(u,w)=(_z\psi ,_x\psi )`$. Also shown are the corresponding profiles for the coverage $`\varphi `$ \[(c) and (d)\].
The two sets of profiles belong to two qualitatively different regimes that are prominently visible in the profiles of the coverage (corresponding to the results obtained for a different disjoining pressure in TJB04 ).
### III.2 Dependence on reaction rate
For low reaction rates \[Fig.3 (a) and (c)\] the coverage starts to increase at the advancing contact zone and continues to increase up to the receding contact zone where it is still well below the saturation value of $`\varphi =1`$. We call this the non-saturated regime. For high reaction rates \[Fig.3 (b) and (d)\] the coverage starts to increase at the advancing contact zone as in the non-saturated regime. However, it increases much faster and reaches the saturation value $`\varphi =1`$ already underneath the droplet. At the receding contact zone it is always at the saturation value. We call this the saturated regime. The droplets in the two regimes behave qualitatively different when the reaction rate is changed as shown in Fig. 4.
For low reaction rates in the non-saturated regime the coverage does not reach the saturation value $`\varphi =1`$ at the rear of the droplet. This implies that the driving wettability gradient between front and rear of the droplet has not yet reached its maximally possible value. Therefore an increase in the reaction rate leads to a steeper increase in the spatial profile of the coverage underneath the droplet. In consequence a larger wettability gradient results implying a larger velocity. However, again it is subtle to determine the dynamic equilibrium. First, the higher velocity reduces the contact time of the droplet with a given point of the substrate. Second, it also reduces the length of the droplet (see below). Both effects restrict the growth of the wettability gradient. The resulting increase of the velocity with increasing reaction rate can be seen in Fig. 4 (a) for $`r<0.0001`$. Note that it does only weakly depend on the diffusion along the substrate.
Fig. 4 (b) shows the corresponding dynamical contact angles at the advancing and receding edges of the moving droplets. Below $`r0.0001`$ both angles increase, implying a decrease in the length of the droplet (that has constant volume). Because in the non-saturated regime the coverage in the contact zone at the advancing edge does effectively not depend on $`r`$, the increase in the advancing contact angle is caused by the larger velocity only. Seeing the full dynamic contact angle at a moving contact line as a superposition of a static (or equilibrium) part and a dynamic part, one can say that only the dynamic part of the advancing contact angle increases with $`r`$. The increase at the front is in line with the increase with velocity found for sliding drops on an incline (see for example Ref. Thie02 ). However, at the receding edge the coverage in the contact zone depends strongly on $`r`$. Therefore the rather strong increase in the receding contact angle with increasing $`r`$ is caused by changes in both, the static and the dynamic part. For a larger $`r`$ the wettability at the back is smaller, i.e. the โlocal equilibrium contact angleโ increases. However, the dynamic receding contact angle normally decreases with increasing velocity Thie02 . This implies that the increase in the receding angle found here results from an increase of the static part only. It is even counteracted by a decrease of the dynamic part. Also in the contact angles, a relatively small $`d`$ has no visible influence in the non-saturated regime.
The maximum wettability gradient, i.e. driving force, and therefore the maximum droplet velocity is reached for the bare substrate ($`\varphi =0`$) at the advancing edge and maximum coverage ($`\varphi =1`$) at the receding edge. The point of largest velocity \[Fig. 4 (a)\] corresponds to the point of the largest difference between the receding and advancing contact angles \[Fig. 4 (b)\]. We call the corresponding reaction rate $`r_{max}`$.
Next we discuss the behavior in the saturated regime ($`r>0.0001`$ in Fig. 4). For larger reaction rates the coverage at the rear of the droplet remains at its saturation value. However, the driving force and in consequence the velocity decrease slightly with increasing $`r`$ (cf. also Ref. TJB04 ). Note, that for small $`d`$ in Fig. 4 (a) this is barely visible. The decrease from $`r_{max}`$ to $`r=0.1`$ corresponds to about 2% of the maximal velocity.
As detailed next this slight decay is caused by the dynamics in the advancing contact zone. In the saturated regime the time scale of the reaction is short compared to the one of the droplet movement. Therefore the coverage increases very steeply underneath the droplet leading to an elevated coverage already in the advancing contact zone (see Fig. 3 b). The resulting decrease of the overall wettability gradient felt by the droplet causes the slight decrease of the velocity. In this regime the slight increase of the advancing contact angle results from an increase of the static part (because of the decreasing wettability) that overcomes a decrease of the dynamic part (decreasing velocity). The slight increase of the receding angle comes from the dynamic part only.
We emphasize here that Fig. 4 (b) clearly shows that the advancing contact angles are always smaller than the receding ones. As was already shown in Ref. TJB04 also here the differences between the static and the dynamic contact angles at the front and the rear, respectively, are an order of magnitude smaller than the difference between the two static (or the two dynamic) contact angles.
### III.3 Phase diagrams
#### III.3.1 Influence of Diffusion
Let us finally discuss the influence of diffusion along the substrate on the dependencies shown in Fig. 4. The result without diffusion (not shown) coincides nearly perfectly with the shown curves for $`d=0.001`$. Only the decay in the saturated regime is slightly slower. We mention here that the decay also depends slightly on the used cutoff $`h_c`$ (cf. Eq. (8) and Ref. TJB04 ). Increasing the diffusion has a rather small influence on the non-saturated regime but changes the saturated regime even qualitatively. The velocity decreases with the reaction rate because the adsorbate produced underneath the droplet close to the advancing edge diffuses onto the bare substrate ahead of the moving droplet. This effectively reduces the wettability there, implying a slower movement. In fact, running droplets cease to exist above a reaction rate $`r_{dp}`$ where the velocity drops to zero again, corresponding to a supercritical bifurcation for large diffusion (see curve for $`d=1.0`$). The bifurcation structure can, however, be more involved. For a smaller $`d=0.1`$ one finds that the branch of running droplets joins the branch of sitting drops in a subcritical bifurcation at $`r_{dp}`$, i.e. decreasing $`v`$ the running droplet branch first turns back in a saddle-node bifurcation at $`r_{sn}`$ and then joins the sitting droplet branch at $`r_{dp}`$. We indicate the definition of $`r_{sn}`$ and $`r_{dp}`$ in Fig. 4 (a) using the curve for $`d=0.1`$ as example. However, the subcritical part of the branch is difficult to see. We refer the reader to the discussion of model II in section IV for more details. Between $`r_{dp}`$ and $`r_{sn}`$ both the large velocity running droplet and the sitting droplet are linearly stable and correspond to metastable states. The branch emerging from the subcritical bifurcation is linearly unstable and corresponds to threshold solutions that separate the two stable solutions. The prominent features ($`r_{max}`$, $`r_{dp}`$ and $`r_{sn}`$) of the solution branches shown in Fig. 4 can be used to determine phase diagrams by continuation techniques AUTO97 . The phase diagrams show existence regions for the different types of solutions in the parameter space spanned, for instance, by the reaction rate and the diffusion constant (see Fig. 5). To this end we followed the maximum of the $`v(r)`$ dependence at $`r_{max}`$, which marks the transition between the non-saturated and saturated regime for running droplets (dashed line in Fig. 5). Continuation of the loci of the saddle-node bifurcation at $`r_{sn}`$ and the drift-pitchfork bifurcation at $`r_{dp}`$ gives the border of the existence region for stable running droplets (dotted line) and stable sitting droplets (solid line), respectively. The region between the two latter borders corresponds to a coexistence region where running and sitting droplets are metastable. Note, that unstable sitting drops do also exist everywhere in the existence region of the running droplets. Such sitting droplets are steady solutions of the governing equations. However, because they are unstable, even infinitely small perturbations will grow with a characteristic rate. In consequence the droplets start to move and adopt the shape and speed of the stable running droplet solution.
#### III.3.2 Influence of volume
Practically, it is difficult to change the reaction rate over orders of magnitude. Experiments normally only cover a much smaller range DoOn95 ; LKL02 . Nevertheless, the analysis of the $`v(r)`$ dependence (Fig. 4) for droplets of fixed volume gives a very good characterization of the general system behavior. An experimentally important parameter is the size or volume of the droplet. The velocity may increase DoOn95 or decrease LKL02 with increasing volume depending on the parameter regime TJB04 . We show in Fig. 6 the dependency of velocity on droplet volume (logarithmic scale) for a variety of reaction rates. Also there one can well distinguish a saturated and a non-saturated regime. In the non-saturated regime the droplet velocity increases with increasing size, whereas in the saturated regime it decreases. For small droplets one always finds a non-saturated regime whereas droplets of a very large size are always in a saturated regime (in Fig. 6 the latter is not yet reached for the curve with $`r=10^5`$). The explanation for this behavior is closely related to the one given for the $`v(r)`$ dependence. For small droplets the reaction does not reach saturation until the rear of the droplet has passed. This implies that an increase in droplet size gives more time for the reaction leading to a larger value of $`\varphi `$ at the back and, in consequence, to a larger velocity.
For very large droplets the reaction has enough time to reach saturation ($`\varphi =1`$) at the rear, i.e. an increase in size does not increase the driving wettability gradient. It may even decrease slightly due to an increase of $`\varphi `$ in the contact zone at the advancing edge (see above). However, a larger droplet has a larger viscous dissipation leading to a decreasing velocity with increasing size. Inspecting Fig. 6 shows that for $`r510^4`$ and $`V=\mathrm{30\hspace{0.17em}000}`$ (i.e. the volume used in Fig. 4) one is well in the saturated regime.
The volume $`V_{max}`$ where the maximal velocity is obtained can be followed in the space spanned by droplet volume and reaction rate. The obtained existence region for non-saturated and saturated running droplets are shown in Fig. 7.
Having studied the type I model describing experiments where the passing droplet irreversibly changes the substrate we next report on results for our type II model where the substrate recovers. In contrast to type I this allows for the description of a periodic droplet movement as experimentally observed in Ref. SMHY05 .
## IV Results for model II (adsorption and desorption)
This section is devoted to the type II model that extends and generalizes the model of type I discussed above. Model II accounts for experimental situations where both, the droplet and its surrounding medium are able to change the substrate by adsorption or desorption. In this way a moving droplet makes the substrate less wettable, but after the droplet has passed the substrate may relax to its initial state. Such systems are modeled by extending the reaction kinetics for the $`\varphi `$-field (14) by an additional desorption term, i.e. we combine the evolution equations for the film thickness profile (11) and the adsorbate field (12) with the reaction term (15) and the disjoining pressure (13). We will refer to this set of equations as type II model. The reaction term (15) is chosen such that adsorption and desorption of $`\varphi `$ take predominantly place underneath and outside the droplet, respectively.
### IV.1 Thickness and coverage profiles
As above for model I, we use continuation techniques to calculate running droplet solutions of model II as solutions stationary in a comoving frame. Examples of the resulting droplet profiles along with the profiles of the $`\varphi `$-field are shown for a very small diffusion $`d`$ in Figs. 8 and 9. Thereby, Fig. 8 displays the results for several values of the reaction rate $`r`$ for a fixed ratio $`s`$ of the desorption to adsorption rates. The drop profiles change only slightly but the $`\varphi `$-profile undergoes prominent changes. As in model I there is practically no $`\varphi `$-field in front of the droplet ($`\varphi =0`$). The field increases (i.e. the wettability decreases) underneath the droplet as in model I, but in contrast the $`\varphi `$-field decreases behind the droplet due to the desorption. The concentration of the $`\varphi `$-field at the rear of the droplet is increasing with $`r`$ until it reaches saturation. Because $`r`$ is the overall reaction rate also the desorption of $`\varphi `$ behind the droplet becomes faster, i.e. the tail of the $`\varphi `$-field becomes shorter.
Fig. 9 displays droplet and $`\varphi `$ profiles for a fixed reaction rate $`r`$ but different values for the desorption to adsorption ratio $`s`$. The droplet shape and the form of the $`\varphi `$-field underneath the droplet are almost independent of $`s`$ as one would expect. However the decay of the elevated $`\varphi `$ value occurring behind the droplet is strongly influenced. For a small $`s`$, i.e. slow desorption as compared to adsorption, $`\varphi `$ decays slowly approaching qualitatively the behavior of model I. Increasing $`s`$ reduces drastically the length of the $`\varphi `$ tail behind the droplet. The length of the tail is very important when studying the periodic movement of a droplet on a finite stripe-like substrate (see below).
### IV.2 Phase diagrams
In the following the existence regions of running and sitting droplets are determined in their dependence on the control parameters reaction rate, desorption to adsorption ratio, diffusion constant and droplet volume. This is done along the lines explained in section III. Continuation gives branches of stationary solutions depending on one control parameter. On these branches special points that separate qualitatively different behavior are identified and followed in the space of parameters. The linear stability of the stationary solutions is also determined. Details on the used techniques can be found in the Appendix.
#### IV.2.1 Influence of the desorption/adsorption ratio
Studying first the influence of $`r`$ and $`s`$ we present in Fig. 11 (a) and (b) branches of stationary solutions in dependence of $`r`$ for different values of $`s`$ characterized by their velocity and the resulting existence regions of running and sitting droplets in the $`rs`$ parameter plane, respectively. Beside the moving droplets there exist sitting droplets (steady states) for all values of $`r`$ and $`s`$. They are symmetrical with respect to the droplet maximum and may be stable or unstable. An example is shown in Fig. 10.
The curve for $`s=1`$ in Fig. 11 (a) is used to better illustrate special values of the reaction rate introduced above in section III. They mark the loci of the maximum of the $`v(r)`$ dependence at $`r_{max}`$, of the saddle-node bifurcation at $`r_{sn}`$ and of the drift-pitchfork bifurcation at $`r_{dp}`$. As before, continuation in $`r`$ of these loci when changing $`s`$ gives the phase diagram Fig. 11 (b). In particular one obtains the border between the non-saturated and saturated regime for running droplets ($`r_{max}`$), the border of the existence region for stable running droplets ($`r_{sn}`$), and the border of the existence region for stable sitting droplets ($`r_{dp}`$). In the small region where both, running and sitting droplets are stable, one finds metastability. There noise, for instance, in the form of substrate inhomogeneities may lead to intermittent droplet movement. The sitting droplets existing in the gray shaded area of Fig. 11 (b) are all unstable.
Inspection of Fig. 11 reveals the rather strong influence of the desorption/adsorption ratio $`s`$ especially on the transition from the non-saturated to the saturated regime and the slowing down in the saturated regime. The extension in $`r`$ of the saturated regime is drastically reduced with increasing $`s`$. For decreasing $`s`$ the maximum of the $`v(r)`$ curve slowly transforms into a plateau and reaches for $`s=0`$ the form typical for model I. Comparing the strong decrease of the velocity with increasing $`r`$ in the saturated regime \[Fig. 11 (a)\] to the slight decrease found for $`s=0`$ (model I, Fig. 4 (a), curve for $`d=0.001`$) poses the question why $`s`$ has such a strong influence.
The influence of $`s`$ parallels the influence of the diffusion $`d`$ in model I. There for large $`d`$ the adsorbate diffuses in front of the advancing droplet rendering the substrate there less wettable. Thereby it decreases the overall wettability gradient between front and rear of the droplet. Here, for large $`s`$ the $`\varphi `$-field is removed behind the rear of the droplet rendering the substrate there more wettable. That also implies a reduction of the overall wettability gradient.
#### IV.2.2 The drift-pitchfork bifurcation
To elucidate the mechanism of the transition between moving and sitting drops we focus for a moment on the drift-pitchfork bifurcation at $`r_{dp}`$ that separates the metastable and the running droplet region. This type of bifurcation is well known from reaction-diffusion (see, for instance, KTB92 ; KrMi94 ; HaMe94 ; OBSP98 ; Boed03 and references therein) and hydrodynamic systems (see, for instance, PrJo88 ; CGG89 ; KnMo90 ; RiPa92 and references therein), where it mediates the transition between steady and travelling structures. In our system it breaks the reflectional symmetry of the sitting droplets leading to moving asymmetric droplets and acompanying travelling asymmetric adsorbate profiles. At the bifurcation a real eigenvalue switches sign, i.e. at $`r_{dp}`$ the velocity of the moving drops is zero. Beyond the bifurcation sitting droplets become unstable and start to move slowly. That makes it unlike a Hopf bifurcation (associated with the zero crossing of the real part of a pair of complex eigenvalues) where the travelling structure has a finite velocity at the bifurcation (for waves on flowing thin liquid films see the discussion in Ref. ThKn04 ). The bifurcation may be subcritical or supercritical (see Figs. 4, 11 and 12) and in consequence the bifurcating branch corresponds to unstable or stable moving drops. Approaching the bifurcation their respective velocity goes to zero as $`\sqrt{|rr_{dp}|}`$ providing an unambiguous signature of the drift-pitchfork bifurcation.
The basic mechanism of the drift-pitchfork bifurcation is connected to the behaviour of the neutral (or Goldstone) mode related to the translational symmetry of the system. This mode with eigenvalue zero is obtained by analyzing the linear stability of the stationary solutions. In general, each continuous symmetry is related to such a neutral mode. In a sense, the neutral modes are the modes that are โclosestโ to zero. This implies that a perturbation or modulation of these modes in an additional degree of freedom (if existing) gives modes that are probable candidates to cross zero and become instability modes. For instance, the transversal (fingering) instability of a liquid front results from a transversal modulation of the longitudinal translational neutral mode ThKn03 . Also one of the two coarsening modes of two liquid droplets is the combination of translational neutral modes of the individual droplets directed in opposite directions MPBT05 .
Here, the drift instability is associated with a mode representing a relative shift between the translational modes of the height and the $`\varphi `$ profiles. Right at the bifurcation this mode corresponds exactly to the translational neutral mode. Although elsewhere one can still identify the translational mode when only looking at the sub-mode for the height profile or the one for the $`\varphi `$ profile, looking at the complete mode one realizes that the relative weight of the two sub-modes is shifted in favor of one of them. This corresponds to the introduction of a relative shift between the two fields. The relative shift breaks the overall reflection symmetry and leads to the movement of the drops.
#### IV.2.3 Influence of diffusion
Next we discuss the influence of $`r`$ and $`d`$ presenting in Fig. 12 dependencies of droplet velocities on reaction rate for different diffusion constants and the resulting phase diagram in the $`rd`$ plane for a fixed desorption/adsorption ratio $`s`$. One notes first that the strength of diffusion has nearly no influence on the non-saturated regime and the transition value $`r_{max}`$ \[the dashed line in Fig. 12 (b) is practically vertical\]. Note that Fig. 12 focuses on relatively small diffusion. For large diffusion the droplets will also stop as shown for model I in Fig. 5. However, the saturated regime does depend on $`d`$ quantitatively as well as qualitatively. For fixed $`r`$, with increasing $`d`$ the stable running droplets become slower. In parallel its existence range in $`r`$ shrinks slowly. As in model I this behavior is mainly caused by the diffusion of the $`\varphi `$-field in front of the advancing droplet. There it increases the coverage thereby reducing the overall wettability gradient, i.e. slowing down the droplets. The qualitative change concerns the character of the drift-pitchfork bifurcation. With increasing $`d`$ it becomes less subcritical and at a critical $`d_c`$ it becomes supercritical, i.e. there exists no coexistence region for sitting and running drops any more. For small diffusion $`d<d_c`$ the existence range of stable sitting drops shrinks with increasing $`d`$ whereas for larger diffusion $`d>d_c`$ it grows.
#### IV.2.4 Influence of volume
Finally, in Fig. 13 we present the phase diagram for the dependence on reaction rate and droplet volume for fixed desorption/adsorption ratio $`s`$ and diffusion constant $`d`$. The locations of the saddle-node, the drift-pitchfork bifurcation and the velocity maximum are all shifted slightly towards smaller $`r`$ when increasing the droplet volume. The range in $`r`$ of the non-saturated regime shrinks slightly, but the range of the saturated regime and the metastable region remain practically constant, they are only shifted towards smaller $`r`$.
Obviously the two contact regions and the substrate outside the droplet are not affected by a change in the droplet volume. However, an increase in droplet size increases the viscous forces and therefore stalls the droplet movement at smaller reaction rates.
### IV.3 One-dimensional numerical simulations
Finally, we employ numerical simulations of running droplets and show that the type II model is able to describe different experimentally found modi of periodic droplet movement SMHY05 ; Sumi05pre . To perform the simulations we use the routine โd02cjcโ provided by the NAG library NAG . It is based on a variable order, variable step size Adamโs method. The simulations either utilize periodic boundary conditions (to model droplets on a ring-like track SMHY05 ) or boundary conditions that mimick non-wettable borders of an otherwise wettable channel (to model droplets on finite stripe-like tracks) SMHY05 ; Sumi05pre . Simulations are started from steady droplet solutions which develop in the absence of a chemical field, i.e. imposing $`\varphi =0`$. Then the droplet movement is initiated by breaking the symmetry of the $`\varphi `$-field by imposing a small gradient and starting the adsorption/desorption reaction. After a short initial transient the running droplets follow periodic trajectories that do not depend on details of the initial symmetry breaking. For the periodic boundary condititions the initial solution was one stationary droplet, whereas in the case of non-wettable boundaries the initial solution were two stationary droplets.
In the case of periodic boundary conditions the droplets move with constant speed and shape after an initial phase. Fig. 14 (a) shows space-time plots of the evolution of the film thickness for different reaction rates $`r`$. The droplet velocity increases with the reaction rate as expected from our continuation results. Fig. 14 (b) shows a comparison of the droplet velocities obtained in the simulations and by continuation. The values for the simulations are estimated after the droplets have reached a constant speed and shape. One finds that both velocities match fairly well. However, for higher reaction rates the simulations slightly overestimate the velocities compared to the continuation results. This results from the lower and equidistant discretization used in the simulation in time.
Experiments with droplets in a finite wettable channel found โregular rhythmic motionโ SMHY05 or different types of โshuttling motionโ along with slowing and stopping behavior Sumi05pre . The here performed simulations show a smooth transition between the different types of periodic movements depending on our control parameters.
In general, the droplets in a wettable channel with non-wettable walls move periodically between the two walls as shown for different reaction rates $`r`$ in the space-time plots in Fig. 15 (a). The initial solution of two stationary droplets quickly coarses upon starting the chemical reactions. The prevailing droplet oscillates between the channel walls with a frequency that depends on $`r`$.
The droplet movement can be classified into two identical but antisymmetric halfcycles (i.e. with different signs in the velocity, and profiles that are related by reflection). Fig. 15 (b) shows the droplet velocity depending on time during one halfcycle exemplary for two reaction rates $`r`$. We find that each half cycle typically contains three phases, distinguishable by their different velocities. After meeting the non-wettable boundary (phase I) the droplet velocity is very low or even zero in the case of very small reaction rates $`r`$. In this phase, the $`\varphi `$-field, which has been produced by thof e passing droplet has first to decrease, until the droplet can return on its own path. This phase of very small velocity is followed by a short phase of a very high velocity (phase II) until the droplet returns to a medium velocity (phase III), that is kept until it meets the opposite wall and the next half cycle begins. The subphases can also be very well distinguished in Fig. 16, were we show hidden line plots of the film thickness profiles (top) and coverage profiles (bottom) for one period of droplet movement. One can clearly see that in phase I (velocity close to zero), after the droplet has encountered the non-wettable wall, the concentration in the $`\varphi `$-field is very high and drops steeply in phase II when the droplet starts moving again. Similar subphases of droplet motion have also been observed in Ref. Sumi05pre . In the context of that work the continuous transition from Fig. 15 (a) between $`r=0.002`$ (top) and $`r=0.0001`$ (bottom) can be seen as their transition between โshuttling motionโ and โintermittent shuttling motionโ.
A simple bead-spring model based on a mechanical analogy is used in Ref. Sumi05pre . It models the different experimentally observed regimes varying the wettability of the walls. Although it well captures the overall behaviour it can not resolve the more hydrodynamic aspects of the motion like the flow field inside the running droplets or the dynamical contact angles.
## V Conclusions
In the present work we have developed and analysed two models for chemically-driven self-propelled running droplets on solid substrates. Such moving droplets were described for several experimental systems using different liquids, substrates and reacting substances BBM94 ; DoOn95 ; LeLa00 ; LKL02 ; SMHY05 ; Sumi05pre . The movment of the droplets is driven by a self-produced wettability gradient that is perpetuated with the droplet itself by means of a desorption or adsorption reaction underneath the droplet.
Two types of experimental systems were reported on: (I) adsorption underneath the droplet decreases the wettability of the substrate irreversibly DoOn95 ; LeLa00 ; LKL02 ; and (II) desorption underneath the droplet removes a more wettable coating that is recovered behind the droplet through an adsorption from a surrounding medium SMHY05 ; Sumi05pre . We have described the droplet dynamics for both types of experimental systems using coupled evolution equations for the profiles of the film thickness and the adsorbate coverage. The equations can be derived from the Navier-Stokes equations using a long-wave or lubrication approximation ODB97 and assuming a small Damkรถhler number. The two types of experiments have been mapped onto two types of models. In the type I model a wettability-decreasing reaction takes place underneath the droplet. In the type II model this mechanism is extended by additionally introducing a wettability-increasing reaction that takes place at the substrate outside the droplet.
The wettability of the substrate enters both models through a disjoining pressure supplementing the Laplace pressure in the thin film equation. We have chosen here a disjoining pressure consisting of a long-range destabilizing part $`h^3`$ and a short-range stabilizing part $`h^6`$ used, for instance, in Ref. PiPo04 to study coarsening in dewetting. The long-range part corresponds to van der Waals interaction and is not influenced by the adsorbate. All the influence of the adsorbate goes into the short-range part that in the simplest case selected here depends linearly on the adsorbate coverage.
Using continuation techniques and numerical simulations we have analyzed the solution behaviour of both models. Thereby we have focused on stationary running and sitting droplets in two dimensions. Both models display a transition from a non-saturated to a saturated regime with increasing reaction rate. The transition is also obtained when increasing the droplet volume. In the non-saturated regime an increase in the reaction rate leads to a larger wettability gradient implying a larger droplet velocity. In the saturated regime an increase in the reaction rate does not increase the coverage at the rear of the droplet, i.e. it does neither lead to a larger wettability gradient nor to a larger droplet velocity. However, it has turned out that the driving force and in consequence the velocity are decreasing slightly with increasing reaction rate. This effect is due to a rise in the adsorbat concentration in the advancing contact zone. A similar behavior occurs when the droplet volume increases. There, however, in the saturated regime the velocity clearly decreases with increasing volume because the constant driving force (wettability gradient) is counteracted by an increasing viscous friction. The latter dependencies found for the non-saturated and the saturated regime correspond very well to experimental results of Ref. DoOn95 (Fig. 5) and Ref. LKL02 \[Fig. 7 (a)\], respectively. To our knowledge there exist, however, no experimental results for a physico-chemical system that show the qualitative change between the two regimes in dependence of the droplet velocity on its volume. With the combination of materials used in Ref. LKL02 this should be possible because their Fig. 6 (a) shows the transition from the non-saturated to saturated regime for increasing solute concentration within the droplet. This corresponds directly to the dependence on the reaction rate shown here in Fig. 4.
Allowing for diffusion of the adsorbate along the substrate does not affect the results in the non-saturated regime too much. However, it leads to a stronger decrease of the velocity in the saturated regime because adsorbate is transported to the substrate in front of the running droplet thereby decreasing the overall wettability gradient. This may even lead to a transition towards sitting droplets. This transition occurs either continuously through a supercritical drift-pitchfork bifurcation or discontinuously through a saddle-node bifurcation. In the latter case a metastable parameter region exist where running and sitting drops can coexist.
We have determined the existence regions of running droplets depending on the reaction rate, the diffusion constant in the surface coating, the droplet volume and the desorption/adsorption ratio (only for model II). Both models show a strong dependence of the existence regions of sitting and running droplets on the diffusion constant and only a week dependence on the droplet volume. The strong dependence on the diffusion constant results mainly from the diffusion of the $`\varphi `$-field in front of the advancing droplet, thereby greatly affecting the driving force and the droplet velocity. The type II model also shows a strong dependence on the desorption/adsorption ratio, which is due to a decreased substrate coverage in close proximity behind the droplet, effectively decreasing the driving force. The primary effect of the desorption reaction is the possibility for the droplet to return on its own path, which we have illustrated by numerical simulations of droplets moving on a wettable stripe with non-wettable borders in an oscillatory manner, which has also been observed experimentally SMHY05 ; Sumi05pre . Finally, we have shown that different modes of periodic movement called โshuttling motionโ and โintermittent shuttling motionโ in Ref. Sumi05pre are covered by the presented model.
The analysis of both models has shown that for the chemically-driven running droplets the advancing contact angle is always smaller than the receding one. As was already shown in Ref. TJB04 also here the differences between the static and the dynamic contact angles at the front and the rear are one order of magnitude smaller than the difference between the two dynamic contact angles. This challenges the assumption of equal dynamic contact angles at the front and the rear that was used in Ref. Brde95 ; deGe98 to develop a simple description of self-propelled running droplets. A simple quantitative theory should instead be based on the assumption that the respective dynamic contact angles equal the different static contact angles at the front and rear.
The proposed type I and type II model based on thin film or lubrication theory can very well reproduce the main features of droplet motion that have been observed experimentally. However, they fail to reproduce the reported damped oscillations in the droplet shape overlaying the continuous droplet movement Sumi05pre . These oscillations could be on the one hand the result of a weakly inhomogeneous surface, since the authors of Ref. Sumi05pre themselves suggested that they have no full control of the experimental surface properties. On the other hand the oscillations could be a sign of a Hopf bifurcation along the solution branch of running droplets. We are well aware that the models presented in this paper are minimal models that, however, reproduce qualitatively most aspects of the behaviour of chemically self-propelled droplets. Refinements of the presented theory could include the viscous motion of the ambient medium as present in the type II experiments. This extension can still be based on the lubrication approximation, for instance, along the lines of the two-layer systems studied in Refs. MPBT05 ; PBMT04 . A second important extension could cover the case of a higher Damkรถhler number. Such a model has to include a description of the transport of the solute within the droplet.
## Appendix A Numerical techniques
### A.1 Continuation
Sitting and running droplets are steady solutions in the laboratory and comoving frame, respectively. They can be followed in parameter space using numerical continuation techniques DKK91 ; DKK91b , for instance, employing the continuation software AUTO97 AUTO97 . The following section highlights some technical details of the employed techniques.
The basic idea behind continuation is that unknown solutions of an algebraic system for a certain set of control parameters are obtained by iterative techniques from known solutions nearby in parameter space. Differential equations of the form
$$u^{}(x)=f(u(x),p)\text{with}f,u๐^n$$
(17)
subject to initial, boundary and integral constraints are discretized in space and then the resulting algebraic system is solved iterativly. Here the dash indicates the first derivative with respect to $`x`$ and $`p`$ denotes the set of control parameters. The presence of boundary conditions and/or integral conditions requires the presence of free parameters which are determined simultaneously and are part of the solution to the differential equation. The package AUTO97 is limited to the continuation of ordinary differential equations (ODEโs), thus it can only be used to compute droplet solutions in two dimensions. As an example we consider the continuation of stationary running droplets, steady in a comoving frame. After transforming Eqs. (11) and (12) into the comoving frame with velocity $`v`$ and integrating the resulting time-independent thin-film equation we have the system of ODEโs
$`h_1^{}`$ $`=`$ $`h_2`$ (18)
$`h_2^{}`$ $`=`$ $`h_3`$ (19)
$`h_3^{}`$ $`=`$ $`{\displaystyle \frac{\mu vh_1}{h_1^3}}\mathrm{\Pi }_x(h_1,\varphi _1)`$ (20)
$`\varphi _1^{}`$ $`=`$ $`\varphi _2`$ (21)
$`\varphi _2^{}`$ $`=`$ $`{\displaystyle \frac{1}{d}}\left(R(h_1,\varphi _1)+v\varphi _2\right),`$ (22)
where $`\mu `$ is an integration constant. It has the physical meaning of a mean flow in the comoving frame. $`h_1`$, $`h_2`$, $`h_3`$, $`\varphi _1`$ and $`\varphi _2`$ denote $`h`$, $`_xh`$, $`_{xx}h`$, $`\varphi `$ and $`_x\varphi `$, respectively. The dashed quantities denote first derivatives with respect to $`x`$. The system is flux conservative, thus we need to specify the integral condition
$$0=\frac{1}{L}h_1๐x\overline{h}$$
(23)
where $`L`$ is the system length and $`\overline{h}`$ is the mean film thickness. Furthermore, for a complete description of the system we introduce periodic boundary conditions for the film thickness
$`h_1(0)`$ $`=`$ $`h_1(L)`$ (24)
$`h_2(0)`$ $`=`$ $`h_2(L)`$ (25)
$`h_3(0)`$ $`=`$ $`h_3(L)`$ (26)
and the mixed boundary conditions
$`0`$ $`=`$ $`\xi (h_1(0))\lambda _1+\{\xi (h_1(0))+s[1\xi (h_1(0))]\}\times `$ (27)
$`\left(\varphi _2(0)\varphi _1(0)\lambda _1\right)`$
$`0`$ $`=`$ $`\xi (h_1(0))\lambda _2+\{\xi (h_1(L))+s[1\xi (h_1(L))]\}\times `$ (28)
$`\left(\varphi _2(L)\varphi _1(L)\lambda _2\right)`$
for the $`\varphi `$-field (see appendix A.2). Since translation of a running droplet solution in the $`x`$-direction also yields a valid solution one needs to introduce a pinning condition in the form of an additional boundary or integral condition, which will not be specified further in here.
As mentioned earlier, the number of boundary conditions (NBCD) and integral conditions (NINT) imposes a constraint on the number of parameters (NPAR), that have to be varied during the continuation process. Specifically NPAR=NBCD+NINT-NDIM+1, where NDIM is the dimensionality of the system of ODEโs. This condition leaves us here with three free parameters, which are the principal continuation parameter (e.g. the reaction rate $`r`$), the mean flow $`\mu `$ and the droplet velocity $`v`$. AUTO97 uses the method of Orthogonal Collocation for discretizing solutions, where the solution is approximated by piecewise polynomials with 2โ7 collocation points per mesh interval. The mesh is adaptive as to equidistribute the discretization error. Having specified the ODE system in standard form with boundary and integral conditions AUTO97 then tries to find stationary solutions to the discretized system, by using a combination of Newton and Chord iterative methods. Once the solution has converged AUTO97 proceeds along the solution branch by a small step in the parameter space defined by the free continuation parameters and restarts the iteration.
The challenge usually is, to provide AUTO97 with a nonuniform starting solution for the continuation. For our purpose it is sufficient to start the continuation close to the point of the primary bifurcation point, where the stable uniform solution of the ODE system (18)โ(22) or a system similar to (18)โ(22) with periodic boundary conditions undergoes a Hopf-bifurcation. In the vicinity of the bifurcation one can determine analytically small-amplitude sinusoidal stationary traveling waves. By selectively using reaction, boundary or integral conditions as primary continuation parameter one finally computes fully nonlinear solutions for the film thickness and the coverage.
AUTO97 is not only able to follow solution branches but can also detect bifurcations, like saddle-node bifurcation or branching points, and can then follow these bifurcations in parameter space.
### A.2 Boundary conditions
The use of periodic boundary conditions does not rule out interactions between droplets in consecutive periods either via the film thickness or the $`\varphi `$-field. Interactions through the $`\varphi `$-field arise if the desorption of the coating is slow compared to the drop movement. One way to avoid this problem is to use very large periods, such that the $`\varphi `$-field has enough time to recover. This approach is successful in suppressing interactions but is very costly from a computational point of view. An alternative approach is to use other than periodic boundary conditions for the $`\varphi `$-field. The following paragraphs illustrate this approach.
We generally consider a running droplet steady in a comoving frame, that has its maximum approximately in the center of the computational domain of length $`L`$. We assume that the film thickness profile obeys periodic boundary conditions $`h_bh(0)=h(L)`$ and $`_zh(0)=_zh(L)1`$. The latter restriction ensures that the minimal film thickness is close to equilibrium and the period is large compared to the droplet.
#### A.2.1 Model I
In model I we assume that there is no chemical reaction taking place outside the droplet. Therefore, outside the computational domain for $`x0`$ and $`xL`$ the following ordinary differential equation holds
$$0=d\varphi _{xx}+v\varphi _x,$$
(29)
which has the general solution
$$\varphi (x)=c+c^{}e^{\frac{v}{d}x}$$
(30)
with $`c`$ and $`c^{}`$ being yet undetermined constants. For simplicity we assume that $`v>0`$, i.e. the droplet is moving to the right. In front of the droplet it is assumed that $`\varphi 0`$ as $`x\mathrm{}`$. Furthermore we assume that the $`\varphi `$-field adopts a finite value behind the droplet such that $`\varphi \varphi _{\mathrm{}}`$ as $`x\mathrm{}`$. First we consider the boundary at $`x=0`$ with the boundary value for the $`\varphi `$-field $`\varphi (0)=\varphi _0`$ and the first derivative of the $`\varphi `$-field $`\varphi _x(0)=\varphi _{x0}`$. Since $`\varphi _{\mathrm{}}`$ is a finite value we find $`c^{}=0`$, which leads to the Neumann boundary condition
$$\varphi _{x0}=0.$$
(31)
Next we consider the boundary at $`x=L`$ with the boundary value for the $`\varphi `$-field $`\varphi (L)=\varphi _L`$ and the first derivative of the $`\varphi `$-field $`\varphi _x(L)=\varphi _{xL}`$. Integrating (29) we find
$$0=d\varphi _x+v\varphi +c^{\prime \prime },$$
(32)
where $`c^{\prime \prime }`$ is an integration constant. Since $`\varphi 0`$ as $`x\mathrm{}`$ and also $`\varphi _x0`$ as $`x\mathrm{}`$ we find $`c^{\prime \prime }=0`$, which leads us to the mixed boundary condition
$$0=\varphi _Lv+\varphi _{xL}d.$$
(33)
#### A.2.2 Model II
We assume for the $`h`$-field outside the computational domain with $`x0`$ or $`xL`$ $`h(x)=h(0)=h(L)=h_b=const.`$. Then the concentration profile for the $`\varphi `$-field for $`x0`$ or $`xL`$ can be computed analytically by solving a linear 2<sup>nd</sup> order ordinary differential equation in $`x`$ of the form
$$0=R_2(h_b,\varphi )+d\varphi _{xx}+v\varphi _x,$$
(34)
which has the general solution
$$\varphi (x)=c+c^{}e^{\lambda _1x}+c^{\prime \prime }e^{\lambda _2x},$$
(35)
where
$$c=\frac{\xi (h_b)}{\xi (h_b)+s\left[1\xi (h_b)\right]}$$
(36)
$$\lambda _{1,2}=\frac{v}{2d}\pm \sqrt{\frac{v^2}{4d^2}+\frac{r}{d}\left\{\xi (h_b)+s\left[1\xi (h_b)\right]\right\}}$$
(37)
and $`c^{}`$ and $`c^{\prime \prime }`$ are yet undetermined constants. For $`x\pm \mathrm{}`$ $`\varphi `$ adopts a small but finite value $`\varphi _{\mathrm{}}`$.
First we consider the boundary at $`x=0`$. This leaves us with $`c^{\prime \prime }=0`$. Using the boundary values for the $`\varphi `$-field $`\varphi (0)=\varphi _0`$ and the first derivative of the $`\varphi `$-field $`\varphi _x(0)=\varphi _{x0}`$ at the boundary $`x=0`$ the following system of equations holds
$`0`$ $`=`$ $`\varphi _0cc^{}`$ (38)
$`0`$ $`=`$ $`\varphi _{x0}c^{}\lambda _1.`$ (39)
Solving for $`c^{}`$ one finds the following mixed boundary condition
$$0=\xi (h_b)\lambda _1+\left\{\xi (h_b)+s\left[1\xi (h_b)\right]\right\}\left(\varphi _{x0}\varphi _0\lambda _1\right).$$
(40)
The same procedure can be performed at the boundary $`x=L`$ with $`\varphi (L)=\varphi _L`$ and $`\varphi _x(L)=\varphi _{xL}`$ yielding the condition
$$0=\xi (h_b)\lambda _2+\left\{\xi (h_b)+s\left[1\xi (h_b)\right]\right\}\left(\varphi _{xL}\varphi _L\lambda _2\right)$$
(41)
In the continuation algorithm the two periodic boundary conditions for the $`\varphi `$-field were substituted by the two mixed boundary conditions obtained above.
## Acknowledgments
We thank E. Knobloch for his comments on the drift-pitchfork bifurcation, and acknowledge support through the EU RTN โUnifying principles in non-equilibrium pattern formationโ (Contract MRTN-CT-2004โ005728). |
warning/0506/cond-mat0506351.html | ar5iv | text | # Ferromagnetism and metal-like transport in antiferromagnetic insulator heterostructures.
## Abstract
Strained Pr<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub>/La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub>/Pr<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> trilayers were grown on (001)-SrTiO<sub>3</sub> substrates using the pulsed-laser deposition technique. The coupling at the interfaces of several trilayers has been investigated from magnetization and electronic transport experiments. An increase of La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> layer thickness induces a magnetic ordering in the strain layers and at the interfaces leading to ferromagnetic behavior and enhanced coercivity, while resistivity shows metal-like behaviors. These effects are not observed in the parent compounds, which are antiferromagnetic insulators, opening a path, to induce artificially some novel properties.
Over the past few years, perovskite type manganites such as $`R_xA_{1x}MnO_3`$ ($`R`$=rare earth elements and $`A`$=alkaline earth elements) have been extensively investigated because of their colossal magneto-resistance properties , a huge decrease in the resistance with applied magnetic field. This material has been explored in the form of thin films including multilayer structures, made by several combinations among ferromagnetic ($`FM`$), antiferromagnetic ($`AFM`$), paramagnetic ($`PM`$) in addition to insulator and/or metal, whose properties are different from their single layer structures. Their physical properties have been attributed to the structural and magnetic modifications at the interfaces of the two constituents of the multilayers. For example, canted spin arrangements of ferromagnetic layer in La<sub>0.6</sub>Sr<sub>0.4</sub>MnO<sub>3</sub>/La<sub>0.6</sub>Sr<sub>0.4</sub>FeO<sub>3</sub> and La<sub>0.6</sub>Sr<sub>0.4</sub>MnO<sub>3</sub>/SrTiO<sub>3</sub> superlattice, formation of interfacial ferromagnetism in CaMnO<sub>3</sub>/CaRuO<sub>3</sub> superlattice and disorder interfacial phase of structural and magnetic origin in La<sub>0.7</sub>Ca<sub>0.3</sub>MnO<sub>3</sub>/LaNiO<sub>3</sub> superlattice have been observed. These examples have confirmed the importance of magnetic interfaces which, in fact, has already been revealed in other metallic systems. However, spin ordering is also observed at interfaces of the superlattices consisting of antiferromagnetic layers of LaFeO<sub>3</sub> and LaCrO<sub>3</sub>. The manganites show interesting properties like the charge/orbital ordering ($`CO/OO`$) that occurs in some half-doped compounds. This $`CO/OO`$ behavior corresponds to an ordering of charges and orbitals in two different Mn sublattices (i.e., a long-range ordering of $`Mn^{3+}`$ and $`Mn^{4+}`$ ions) below the $`CO/OO`$ temperature, when the materials is cooled down to low temperature. This $`CO/OO`$ state is highly an insulating state, but can be destroyed (i.e., inducing a metallic behavior) by the application of a magnetic field. Similar two prototype compounds are Pr<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> ($`PCMO`$) and La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> ($`LCMO`$) exhibits insulator like behaviour with CO/OO temperature ~175K and ~180K, respectively (bulk $`LCMO`$ does not show $`FM`$ behavior). Here, we report the interface effect between $`PCMO`$ and $`LCMO`$, in detail, through transport and magnetic measurements of $`PCMO`$/$`LCMO`$/$`PCMO`$ trilayer structures grown on (001)-oriented SrTiO<sub>3</sub> substrates (STO).
Thin films and heterostructures of $`PCMO`$ and $`LCMO`$ were fabricated by the pulsed laser deposition ($`PLD`$) technique using a KrF laser ($`\lambda =248`$ $`mm`$) on STO substrates. The samples were grown at $`720`$ $`{}_{}{}^{}C`$ in an oxygen ambient of $`300`$ $`mtorr`$. The deposition rates (typically $`\mathrm{~}`$ $`0.38A/pulse`$) of $`PCMO`$ and $`LCMO`$ were calibrated for each laser pulse of energy density $`\mathrm{~}3J/cm^2`$. After the deposition the chamber was filled to $`400torr`$ of oxygen at a constant rate, and then the samples were slowly cooled down to room temperature at the rate of $`20^{}C/\mathrm{min}`$. Trilayer structures comprising of $`50`$-($`unit`$ $`cell`$, $`u.c.`$) $`PCMO`$/$`n`$-($`u.c.`$) $`LCMO`$/$`10`$ -($`u.c.`$) $`PCMO`$, with $`n`$ taking integer values from $`1`$ to $`18`$, were thusly made. To reduce the substrate-induced strain, we have deposited a thicker bottom layer than the top layer, while the top layer is made thinner to have the effect of interface in the transport measurements. The structural study was done by x-ray diffraction (XRD) using a Seifert XRD 3000P (Cu, K$`\alpha 1`$, $`\lambda =0.15406nm`$). Special arrangement for the large intensity of x-ray beam and large surface area of the sample was used during the $`\theta `$ \- $`2\theta `$ scan. The resistivity ($`\rho `$) was measured using a four-probe method with in-plane current. The measurements are done by putting silver contact pads, with a separation of 6 mm between the voltage electrodes of the sample with lateral dimension close to $`3\times 10`$ mm<sup>2</sup>. Magnetotransport and magnetization measurements were performed with magnetic field aligned along the \[$`100`$\] direction of the substrate. The samples were cooled to a desired temperature from room temperature in the absence of electric and magnetic field to perform transport measurements.
All sample shows ($`00l`$) fundamental Braggโs reflections of the substrate and the constituents of the heterostructure indicate the epitaxial growth of the trilayers. The pseudocubic lattice parameter of $`STO`$ ($`3.905A`$) is larger than the lattice parameter of $`PCMO`$ ($`3.802A`$) and $`LCMO`$ ($`3.83A`$) provides in-plane tensile strain for their epitaxial growth. The $`\theta `$ \- $`2\theta `$ scan close to the fundamental ($`001`$) diffraction peak of the substrate for the samples with $`n=5`$ and $`18`$ are shown in the Fig. 1. The out-of-plane lattice parameter of various samples with different $`LCMO`$ thicknesses are shown in the inset of Fig. 1. It also includes the bulk lattice parameters of $`LCMO`$ and $`PCMO`$. The $`c`$-axis lattice parameter of these sample increases monotonically with the increase in $`LCMO`$. The $`c`$-axis lattice parameter of the sample with $`n=18`$ is lower than the lattice parameter of both $`LCMO`$ and $`PCMO`$, which suggests the presence of substrate-induced strain state in these samples. As a consequence, this modification in the structure of $`PCMO`$ and $`LCMO`$ with the interfaces in the heterostructures are expected to effect the magnetic as well as transport properties of the constituents.
Both materials, $`PCMO`$ and $`LCMO`$, are insulators with an $`AFM`$ behavior in the range of our measurements. Similar insulator-like temperature dependent resistivity $`\rho `$($`T`$) is observed in the thin films of $`LCMO`$ and $`PCMO`$. In Fig.2, we show the zero-field cooled $`\rho (T)`$ in presence of $`0T`$ and $`7T`$ in-plane magnetic field for three heterostructures. Though the constituents of the samples are insulator the resistivity of these samples was calculated using their actual dimensions. As we cool the sample with 5 u.c. thick $`LCMO`$ layer, the zero-field resistivity remains insulator-like down to $`100`$ $`K`$. Below $`100K`$, the resistance of the sample is high and it is limited by the input impedance of PPMS (Physical Properties Measurement System). However, for the sample with 10 u.c. thick $`LCMO`$, the resistivity below room temperature is insulator-like down to $`10K`$. As the $`LCMO`$ layer thickness increases 18 u.c. and the sample is cooled from room temperature the resistivity shows thermally activated behavior down to $`150K`$, shows metal-like behavior in the temperature range of $`150K`$ to $`30K`$ and an upturn below $`30K`$. In presence of $`7T`$ magnetic field, the $`\rho `$($`T`$) of the superlattice with 5 u.c. $`LCMO`$ is similar to that of the zero-field $`\rho `$($`T`$) of the sample with 18 u.c. thick $`LCMO`$. Qualitatively similar in-field $`\rho `$($`T`$) with the broader metal-like window and higher metal-insulator transition temperature is observed for higher LCMO layer thickness. The temperature dependent magnetoresistance ($`MR`$) of these heterostructures are shown in the inset of Fig.2. The $`MR`$ of both samples is decreasing as the temperature increasing, similar to the bulk materials.
To understand this transport behavior of these heterostructures, we have measured their magnetic properties. The temperature dependent field cooled (0.01T) magnetization of three heterostructures are shown in the Fig.3a. The magnetization is shown after the diamagnetic correction of $`STO`$. On heating from $`10K`$, the heterostructure with 3 u.c. thick $`LCMO`$ shows a sharp antiferromagnetic transition at $`30K`$ and a ferromagnetic-to-paramagnetic transition at $`250K`$. As the $`LCMO`$ layer thickness increases the antiferromagnetic behavior suppress and these samples show the same Curie temperature (T<sub>C</sub>). At $`10K`$, the magnetization of a thin film of $`PCMO`$ is 1.51$`\times `$10<sup>-3</sup> emu/gauss/cm<sup>3</sup> which is same order of magnetization (9.033$`\times `$10<sup>-3</sup> emu/gauss/cm<sup>3</sup>) observed for the heterostructure with 18 u.c. thick $`LCMO`$.
The zero-field-cooled magnetic hysteresis loop measured at $`10K`$ of the samples with $`LCMO`$ layer thickness above 5 u.c. shows ferromagnetic behavior. The magnetic hysteresis loop of the sample with 18 u.c. thick $`LCMO`$ layer is shown in the Fig. 3b. The magnetization of the sample increases gradually with the increase in magnetic field and does not show a clear saturation. We have extracted the $`M_S`$ (saturation magnetization) taking into account of the weak diamagnetic response of the substrate by extrapolating the linear part of the hysteresis loop to $`\mu _0H=0`$. The extracted value of $`M_S`$ is $`1.98`$ $`\mu _B/Mn`$. This value of $`M_S`$ is small compare the theoretical value of a ferromagnetic phase of $`(Pr,La)_{0.5}Ca_{0.5}MnO_3`$ composition ($`3.5`$ $`\mu _B/Mn`$). Though this sample does not exhibits a clear saturation magnetization, but shows a significant coercive field ($`H_C`$) ( ~0.07tesla). However, the shape of the hysteresis loop does not change as we cool the sample in presence of magnetic field. The field-cooled and zero-field-cooled magnetic behavior of these sample does not indicate the presence of FM cluster as expected from the $`FM/AFM`$ exchange bias system. The interesting point is the origin of the $`FM`$ behavior, where both materials are $`AFM`$. Ferromagnetic ordering of spin at the interfaces of the superlattices consisting of two antiferromagnet $`LaCrO_3`$ and $`LaFeO_3`$ has already been realized, with the ordering along the $`[111]`$ direction. These samples show clear saturation in the magnetic hysteresis loop and the author have explained $`FM`$ ordering due to the Goodenough-Kanamori rules. However, to the best of our knowledge, the metal-like transport in an heterostructure using insulators as the constituent has not been reported so far.
Several interesting behaviors have been observed in manganites by changing average $`A`$-site ionic radius. This doping process provides modification in the MnO<sub>6</sub> octahedra and formation of $`FM`$ metallic phase in the $`AFM`$ insulating matrix. As seen in the inset of Fig.1, both $`PCMO`$ as well as $`LCMO`$ are in the strain state due to the lattice mismatch between them and with the substrate. We believe that this effective strain modify the structure as well as the spin configurations in the $`PCMO`$, $`LCMO`$ and at their interface. However, the modification at the $`PCMO/LCMO`$ interface is expected to be more due to the interfacial-induced strain and the presence of $`A`$-site ion La, Pr and Ca. This interfacial stress might induce a spin re-orientation, which will modify the Jahn-Teller distortion of the $`MnO_6`$ octahedra compare to the bulk materials and it may results in spin ordering and/or spin canting at the interfaces, though interfacial magnetic modification like spin ordering, spin frustration and spin canting has been observed in different $`FM/AFM`$ systems . The lower value of T<sub>N</sub> in the sample with 5 u.c. thick $`LCMO`$, the increase in normalized magnetization and the same value of T<sub>C</sub> with the increase in $`LCMO`$ layer thickness suggest the possibility of the presence of $`FM`$ domain at the interfaces due to the weakening of $`CO/OO`$ ordered state, which is expected as $`CE`$ type ordering is most susceptible to disorder. Since these samples show remarkable H<sub>C</sub>, we believe that the presence of $`FM`$ domains at the interfaces and the coupling at the $`FM`$ domain boundary between the $`FMAFM`$ may responsible for the remarkable value of H<sub>C</sub>. The existence of $`FM/AFM`$ interaction should be realized in the field-cooled hysteresis loop, but perhaps we could not able to observed this effect due to the small interfacial volume. However, this effect was realized once the number of interface is increasing.
The heterostructures with $`LCMO`$ layer thickness larger than 10 u.c. show a clear metal-insulator transition in the $`\rho `$($`T`$). Also the resistivity as well as the activation energy in the insulator-like region decreases with the increase in $`LCMO`$ layer thickness. For example, the activation energy of the heterostructure with 5 u.c. thick $`LCMO`$ is ~0.12eV which is lower than the thin film of $`PCMO`$ (~0.29eV) and $`LCMO`$ (~0.32eV) on $`STO`$. However, the activation energy of these sample in the paramagnetic state of the $`\rho `$($`T`$) is same. The in-field resistivity of these sample show thermally activated behavior below 30 K, which we attribute to the grain boundary-like tunneling. This also suggest the enhancement of $`FM`$ phase at the expense of $`CO/OO`$ phase in the $`LCMO`$ with the grain boundary modification. Since the resistance of the sample with $`n=5`$ is insulator-like in the entire temperature range, we believe that the electronic and magnetic transport in these samples are due to spin polarized tunneling and percolative conduction as described in the conduction process of bulk $`PCMO`$ and $`LCMO`$. Though there may be some other mechanisms responsible for the metal-like conduction in the $`\rho `$($`T`$), we attribute it mainly to the presence of $`FM`$ domains near or at the interfaces and the substrate-induced lattice distortion. The presence of $`FM`$ domains may partially open double exchange conducting percolative path, which induced the metal-like transport in these samples.
In conclusion, we fabricated epitaxial Pr<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub>/La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub>/Pr<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> trilayers, where the parent compounds are $`AFM`$ insulator. We observed ferromagnetic and metal-like behavior with the increase in the La<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> layer thickness. We proposed that the distribution of magnetic order and magnetic moments near the interface are mainly responsible for the ferromagnetic behavior. The presence of interfacial spin ordering may opens the double exchange percolative conducting path to facilitate metal-like behavior in these heterostructures.
We greatly acknowledge financial support of Centre Franco-Indien pour la Promotion de la Recherche Avancee/Indo-French Centre for the Promotion of Advance Research (CEFIPRA/IFCPAR) under Project N2808-1 and the Ministรจre de la Jeunesse et de lโEducation Nationale (2003/87). We thank Dr. H.W. Eng for discussions.
Figure 1: Typical room temperature $`\mathrm{\Theta }2\mathrm{\Theta }`$ x-ray diffraction pattern around the (001) Braggโs peak of (50 u.c.) $`PCMO/`$(n u.c.) $`LCMO/`$(10 u.c.) $`PCMO`$ trilayer with (a) $`n=5`$ and (b) $`n=18`$ grown on (001)-oriented STO. The inset shows the out-of-plane lattice parameter of various trilayer structures with different $`LCMO`$ layer thicknesses. The bulk lattice parameters of $`PCMO`$ and $`LCMO`$ are also indicated.
Figure 2: Zero-field-cooled $`\rho (T)`$ in the presence of $`0`$ tesla (open symbols) and $`7`$ $`tesla`$ (full symbols) magnetic field for different trilayers. The inset depicts the $`MR`$ ($`MR=\frac{\rho (0)\rho (7T)}{\rho (7T)}`$) for the heterostructures with $`n`$ $`=`$ $`5`$ and $`18`$. Note that under magnetic field, all curves show an insulator-to-metal transition.
Figure 3(a): Field-cooled temperature dependent magnetization of different heterostructures at $`10K`$ at 0.01 tesla magnetic field. (b) Zero-field-cooled magnetic field dependent magnetization at $`10K`$ of the heterostructure with $`18`$ $`u.c.`$ thick $`LCMO`$. |
warning/0506/hep-ph0506255.html | ar5iv | text | # Proton Decay in a 6D ๐โข๐โข(10) model
## 1 $`SO(10)`$ model in 6D
Our starting point is an $`SO(10)`$ gauge theory in 6D with $`N=1`$ supersymmetry compatified on the orbifold $`T^2/(Z_2^I\times Z_2^{PS}\times Z_2^{GG})`$ $`^{\mathrm{?},\mathrm{?}}`$. The theory has four fixed points, $`O_I`$, $`O_{PS}`$, $`O_{GG}`$ and $`O_{fl}`$, located at the four corners of a โpillowโ corresponding to the two compact dimensions. At $`O_I`$ only supersymmetry is broken whereas at the other fixed points, $`O_{PS}`$, $`O_{GG}`$ and $`O_{fl}`$, also the gauge group $`SO(10)`$ is broken to its three GUT subgroups G<sub>PS</sub>= $`SU(4)\times SU(2)\times SU(2)`$, G$`{}_{GG}{}^{}=SU(5)\times U(1)_X`$ and flipped $`SU(5)`$, G$`{}_{fl}{}^{}=SU(5)^{}\times U(1)^{}`$, respectively. The intersection of all these GUT groups yields the standard model group with an additional $`U(1)`$ factor, G$`{}_{SM^{}}{}^{}=SU(3)\times SU(2)\times U(1)_Y\times U(1)_Y^{}`$, as unbroken gauge symmetry below the compactification scale. The field content of the theory is strongly constrained by imposing the cancellation of irreducible bulk and brane anomalies $`^\mathrm{?}`$. We study the model proposed in $`^\mathrm{?}`$, containing 3 16-plets $`\psi _i`$, $`i=1\mathrm{}3`$, as brane fields and 6 10-plets, $`H_1,\mathrm{},H_6`$, and 4 16-plets, $`\mathrm{\Phi },\mathrm{\Phi }^c,\varphi ,\varphi ^c`$, as bulk hypermultiplets. Vacuum expectation values of $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^c`$ break the surviving $`U(1)_{BL}`$. The electroweak gauge group is broken by expectation values of the anti-doublet and doublet $`H_{u/d}`$ contained in $`H_1`$ and $`H_2`$.
We choose the parities of $`H_5`$, $`H_6`$ and $`\varphi ,\varphi ^c`$ such that their zero modes
$$L=\left(\begin{array}{c}\nu _4\hfill \\ e_4\hfill \end{array}\right),L^c=\left(\begin{array}{c}\nu _4^c\hfill \\ e_4^c\hfill \end{array}\right),D^c=d_4^c,D=d_4$$
(1)
act as a (vectorial) fourth generation of d-quarks and leptons and mix with the three generations of brane fields, located on the three branes where $`SO(10)`$ is broken to its three GUT subgroups. This leads to a characteristic pattern of mass matrices of the lopsided type as described in $`^\mathrm{?}`$. In particular the hierarchy and mixing angles for the quark sector can be accounted for and GUT relations do not hold for all the generations due to the presence of split multiplets. See $`^\mathrm{?}`$ for the details and $`^\mathrm{?}`$ for the explicit form of the mixing matrices.
## 2 Short review of 4D proton decay in supersymmetry
Proton decay arises from effective 4-fermion operators joining three quarks and a lepton, which are of dimension 6. It can therefore be a bit puzzling to hear about โdimension 4โ or โdimension 5โ operators, as happens in supersymmetric models. So I will stop a little and review the terminology before discussing the dominant contribution in our case.
In supersymmetric models, there are contributions to proton decay from superpotential terms, either renormalizable or obtained integrating out heavy states, from kinetic terms and also from supersymmetry breaking terms. The first type of contributions, are usually classified according to the dimension of the superpotential terms that break the baryon or lepton number. In general we have
* dimension four operators:
$$W=\lambda LLE^c+\lambda ^{}LQD^c+\lambda ^{^{\prime \prime }}U^cD^cD^c;$$
(2)
they are renormalizable and give very rapid proton decay via an intermediate scalar squark (therefore the effective 4-fermion operator is just suppressed by $`\frac{\lambda ^{}\lambda ^{^{\prime \prime }}}{m_{susy}^2}`$, where $`m_{susy}^2`$ is the typical squark mass) and have to be excluded by a discrete symmetry, usually R-parity.
* dimension five operators:
$$W=\frac{1}{M_{H_C}}[\frac{1}{2}Y_{qq}Y_{ql}QQQL+Y_{ue}Y_{ud}E^cU^cU^cD^c;]$$
(3)
they arise e.g. from integrating out the heavy colored Higgs triplets with mass $`M_{H_C}`$ and allow the decay via a loop of scalar superpartners. They produce effective 4-fermion operators that scale as $`\frac{1}{M_{H_C}m_{susy}}`$, are color antisymmetric and therefore must be also flavour non-diagonal, so that the dominant channel results in $`pK^+\overline{\nu }`$. Such operators give the dominant contribution in the simple supersymmetric SU(5) case, and they can give proton decay even above the present limit $`^\mathrm{?}`$.
* โrealโ dimension 6 operators, arising from the fermion kinetic terms and mediated by the gauge multiplet; they are therefore not of the chiral type and cannot be written as superpotential terms. They do not involve sparticles and are therefore independent of $`m_{susy}`$, apart for the weak dependence coming from determining the GUT scale by RGEs. As an example, in $`SU(5)`$ we have the exchange of the $`๐ณ`$ leptoquark gauge bosons with masses $`M_๐ณ`$. The effective vertex is give by the Fermi-type coupling:
$$_{eff}=\frac{g_5^2}{2M_๐ณ^2}ฯต_{\alpha \beta \gamma }\overline{u^c}_{\alpha ,i}\gamma ^\mu Q_{\beta ,i}\left[\overline{e^c}_j\gamma _\mu Q_{\gamma ,j}\overline{d^c}_{\gamma ,k}\gamma _\mu L_k\right]+\text{h.c}.,$$
(4)
where $`i,j`$ and $`k`$ count the generations. With Fierz reordering, one can write the operators in the usual form as
$$_{eff}=\frac{g_5^2}{M_๐ณ^2}ฯต_{\alpha \beta \gamma }\left[\overline{e^c}_j\overline{u^c}_{\alpha ,i}Q_{\beta ,i}Q_{\gamma ,j}\overline{d^c}_{\alpha ,k}\overline{u^c}_{\beta ,i}Q_{\gamma ,i}L_k\right]+\text{h.c.}.$$
(5)
Note that these operators scale as $`M_๐ณ^2`$ and are proportional to a gauge coupling and not a Yukawa; still some flavour dependence arises from the quark and lepton mixing matrices. Assuming quarks and leptons to be embedded into multiplets according to their hierarchy, the dominant decay channel is the one involving only first generation fermions, i.e. $`p\pi ^0e^+`$.
* dimension 6 operators coming from supersymmetry breaking, e.g. the one mediated by intermediate Higgs scalars mixing via the soft SUSY breaking mass terms; they are usually much more suppressed compared to the previous ones and are usually neglected.
## 3 6D proton decay
In our 6D orbifold model we have a residual 4D N=1 supersymmetry and we could have in principle dimension 4 and dimension 5 proton decay. Luckily in extra-dimensional models, we can exclude them both with appropriate choice of R-symmetry, forbidding the $`\lambda ^{},\lambda ^{\prime \prime }`$ couplings and also the $`\mu `$ term $`^\mathrm{?}`$. Note also that in extra-dimensional models, the two heavy triplet Higgs bosons become massive together with their N=2 superpartners, not directly with each other, so the mixing between them is generated only by supersymmetry breaking. So in general the dominant contribution to proton decay in orbifold models comes from dimension 6 operators $`^\mathrm{?}`$. Moreover in our model, since the first generation of $`u`$ quarks are confined to live on the fixed point where $`SU(5)`$ is unbroken, we can use $`SU(5)`$ language to describe the operators, even if the bulk symmetry is $`SO(10)`$ <sup>b</sup><sup>b</sup>b A small effect from the other gauge bosons can arise from brane derivative operators $`^\mathrm{?}`$..
### 3.1 Effective operator in 6D
In our orbifold model there is an important difference compared to the 4D case, we have to take into account the presence of a Kaluza-Klein tower of X gauge bosons with masses given by
$$M_๐ณ^2(n,m)=\frac{(2n+1)^2}{R_5^2}+\frac{(2m)^2}{R_6^2}$$
(6)
for $`n,m=0`$ to $`\mathrm{}`$. The lowest possible mass is $`M_๐ณ(0,0)=1/R_5`$, as given by the $`SU(5)`$ breaking parity $`^{\mathrm{?},\mathrm{?}}`$. Note that if we define the 4D gauge coupling as the effective coupling of the zero modes, the KK modes interact more strongly by a factor $`\sqrt{2}`$ due to their bulk normalization.
To obtain the low energy effective operator, we have then to sum over the Kaluza Klein modes. We can define
$$\frac{1}{(M_๐ณ^{eff})^2}=2\underset{n,m=0}{}\frac{1}{M_๐ณ^2(n,m)}=2\underset{n,m=0}{}\frac{R_5^2}{(2n+1)^2+\frac{R_5^2}{R_6^2}(2m)^2};$$
(7)
taking the limit $`R_6/R_50`$, we regain the finite 5D result $`^\mathrm{?}`$,
$`2{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{R_5^2}{(2n+1)^2}}={\displaystyle \frac{\pi ^2R_5^2}{4}}.`$ (8)
But in 6D the summation shows a logarithmic divergence; since our theory is non-renormalizable and valid only below the scale $`M_{}`$, where the theory becomes strongly coupled and 6D gravity corrections are no more negligible, we regulate the sum with the cut-off $`M_{}`$, and obtain formally
$`{\displaystyle \frac{1}{(M_๐ณ^{eff})^2}}{\displaystyle \frac{\pi }{4}}R_5R_6\left[\mathrm{ln}\left(M_{}R_5\right)+C\left({\displaystyle \frac{R_5}{R_6}}\right)+๐ช\left({\displaystyle \frac{1}{R_{5/6}M_{}}}\right)\right].`$ (9)
In the case $`R_5=R_6=1/M_c`$ the expression can be approximated by
$`{\displaystyle \frac{1}{(M_๐ณ^{eff})^2}}{\displaystyle \frac{\pi }{4M_c^2}}\left[\mathrm{ln}\left({\displaystyle \frac{M_{}}{M_c}}\right)+2.3\right],`$ (10)
which agrees within 1% with explicit discrete sum for $`\frac{M_{}}{M_c}=10\mathrm{}50`$.
## 4 Flavour structure
Another important difference in 6D is the non-universal coupling of the $`๐ณ`$ gauge bosons. In fact, due to the parities and the $`SO(10)`$ breaking pattern, their wavefunctions must vanish on the fixed points $`O_{PS}`$ and $`O_{fl}`$, and therefore no coupling arises with the charm and top quark or to the brane states $`d_{2,3}^c,l_{2,3}`$. We have in principle couplings to the bulk states $`d_4^c,d_4`$ and $`l_4,l_4^c`$, but in this case, the charge current interaction always mixes the light states with the heavy KK modes and it is therefore irrelevant for the low energy proton decay $`^\mathrm{?}`$. So the kinetic coupling in Eqn. (5) arises only for the $`1st`$ flavour eigenstate, not for all flavours as in the usual $`4D`$ case.
Proton decay involves only the light quark states and the operators containing the combinations $`uud`$ and $`udd`$. Starting in the basis where the up-quark Yukawa is diagonal, we have to rotate the down-type quarks and the leptons from the weak into the mass eigenstates and single out the contributions for the lightest generation. We have then
$`d_L=U_L^dd_L^{},e_L=U_L^ee_L^{},\nu _L=U_L^\nu \nu _L^{},d_R=U_R^dd_R^{},e_R=U_R^ee_R^{},`$ (11)
where the prime denotes mass eigenstates. Since the up quarks are diagonal, $`U_L^d`$ coincides with the CKM-matrix. We can write the proton decay operators of Eqn. (5) in mass eigenstates as
$`_{eff}`$ $`=`$ $`{\displaystyle \frac{g_5^2}{(M_๐ณ^{eff})^2}}ฯต_{\alpha \beta \gamma }[2\overline{e^c}_k^{}\left(U_R^e\right)_{k1}\overline{u^c}_{\alpha ,1}d_{\beta ,m}^{}\left(U_L^d\right)_{1m}u_{\gamma ,1}`$ (13)
$`+\overline{d^c}_{\alpha ,l}^{}\left(U_R^d\right)_{l1}\overline{u^c}_{\beta ,1}(u_{\gamma ,1}\left(U_L^e\right)_{1j}e_j^{}d_{\gamma ,m}^{}\left(U_L^d\right)_{1m}\left(U_L^\nu \right)_{1j}\nu _j^{})]+\text{h.c.}`$
Note again that due the orbifold construction only the first weak eigenstates couple to the $`๐ณ`$ bosons, instead of all of them. So the proton decay in this 6D model has naturally different branching ratios compared to a 4D model with the same mixing matrices.
## 5 Results
### 5.1 Bound on $`M_c`$
Considering the dominant channel $`pe^+\pi ^0`$, a lower bound on the compactification scale can be derived from the SuperKamiokande limit on the lifetime. We have in fact
$$\mathrm{\Gamma }K_{had}^{\pi ^0}\frac{\pi ^2}{16}\frac{M_{}^4}{M_c^4}\left(\mathrm{ln}\left(\frac{M_{}}{M_c}\right)+2.3\right)^2\left[\mathrm{\hspace{0.17em}4}V_{ud}^4+\frac{\stackrel{~}{M}_2^{d\mathrm{\hspace{0.17em}2}}}{\stackrel{~}{M}_1^{d\mathrm{\hspace{0.17em}2}}+\stackrel{~}{M}_2^{d\mathrm{\hspace{0.17em}2}}}\frac{\stackrel{~}{M}_2^{e\mathrm{\hspace{0.17em}2}}}{\stackrel{~}{M}_1^{e\mathrm{\hspace{0.17em}2}}+\stackrel{~}{M}_2^{e\mathrm{\hspace{0.17em}2}}}\right],$$
(14)
where $`K_{had}^{\pi ^0}=1.87\times 10^{40}`$ sec<sup>-1</sup> contains a factor $`M_{}^4`$, the hadronic matrix element, kinematical factors, gauge coupling and the running of the operator from the high to the proton scale $`^\mathrm{?}`$. With $`M_{}=10^{17}\text{GeV}`$ and $`\stackrel{~}{M}_2^{d,e}=0`$, the limit $`\tau 6.9\times 10^{33}`$ yields $`M_c0.89\times 10^{16}\text{GeV}`$, not far from the 4D GUT scale.
### 5.2 Rates and branching ratios
We calculate the branching ratios for dimension 6 proton decay in our model and in some lopsided 4D models with a similar flavour structure discussed in $`^\mathrm{?}`$. As expected, we find sizable differences in many channels, most noticeably in $`p\mu ^+K^0`$ due to the absence of direct coupling of the second generation weak eigenstates to the $`๐ณ`$ gauge bosons. Even changing the unknown high energy parameters does not modify the picture: if we vary the heavy masses $`\stackrel{~}{M}_j/\stackrel{~}{M}=0.1,\mathrm{},1`$ and $`\stackrel{~}{\mu }_3^{d,e}/\mu _3=1,\mathrm{},5`$, we still find $`BR(p\mu ^+K^0)5\%`$.
## 6 Conclusions
We have studied dimension 6 proton decay in a particular orbifold model, where the flavour eigenstates are placed at different fixed points. We have found two very interesting results. First the predicted decay rate is enhanced compared to 4D, even if still compatible with the experimental bounds. In fact, for a cut-off scale $`M_{}=10^{17}`$ GeV, we set a strong lower bound on the compactification scale $`M_c0.89\times 10^{16}\text{GeV}`$, which means that the range of validity of our model is more restricted also than that of 5D orbifold models $`^\mathrm{?}`$. Secondly, the peculiar flavour structure can give striking signatures in the branching ratios for proton decay, suppressing strongly the decay into $`\mu ^+K^0`$. This is due to the localization of states in the extra-dimension and the consequent non-universal coupling of the GUT bosons to the fermions.
## Acknowledgments
It is a pleasure to thank my collaborators Takehiko Asaka, Wilfried Buchmรผller, David Emmanuel-Costa and especially Sรถren Wiesenfeldt for discussions and for the fruitful and enjoyable collaboration. I would also like to thank the organizers of Moriond EW โ05 for the exciting atmosphere at the workshop and for partial local support.
## References |
warning/0506/gr-qc0506066.html | ar5iv | text | # Half polarized ๐โข(1) symmetric vacuum spacetimes with AVTD behavior
## 1 Introduction.
A rigorous study of the singularities in cosmological solutions of the vacuum Einstein equations has been hampered by the fact that the generic such solution is expected to have a singularity of the oscillatory type predicted by Belinsky, Lifshitz and Khalatnikov \[BLK\]. There is currently no satisfactory mathematical method for treating such oscillating singularities, at least when the spacetimes under study are spatially inhomogeneous. For this reason much effort has gone into the study of families of solutions which have milder cosmological singularities such as those of AVTD (asymptotically velocity term dominated) type, for which rigorous, so called Fuchsian, methods are available.
To suppress the oscillatory behavior expected for the generic solution, one can
* (i) introduce suitable matter sources such as scalar fields and study the solutions of the associated non vacuum field equations \[AR\],
* (ii) study higher dimensional models motivated by string or supergravity theories wherein (for sufficiently high dimensions at least) the oscillations are naturally suppressed \[DHRW\]; or
* (iii) remain in 3+1 dimensions but impose a combination of symmetry and polarization conditions in order to achieve the desired AVTD behavior.
For the case of $`U(1)`$ symmetric vacuum solutions on the trivial $`S^1`$ bundle $`T^2\times R\times S^1T^2\times R`$ (with $`U(1)`$ symmetry imposed on the circular fibers) Isenberg and Moncrief \[IM\] have showed, using Fuchsian methods, that AVTD behavior is achieved provided the solutions considered are at least half polarized in a certain well defined sense. The half polarization condition includes, as a special case, the fully polarized solutions wherein the 3 planes orthogonal to orbits of the $`U(1)`$ isometry action are integrable and the vacuum 3+1 field equations reduce to a system of 2+1 dimensional Einstein equations coupled to a massless scalar field on the quotient manifold $`T^2\times R.`$ The more general (half polarized) solutions admit, in addition, half the extra (asymptotic) Cauchy data expected for a fully general, non polarized solution of the same ($`U(1)`$ symmetric) type. On the basis of numerical studies due to Berger and Moncrief the fully general, non polarized $`U(1)`$ symmetric vacuum solution on $`T^2\times R\times S^1T^2\times R`$ is expected to have an oscillatory singularity and hence not to be amenable to Fuchsian analysis \[BM\].
Choquet-Bruhat, Isenberg and Moncrief have extended the analysis given in \[IM\] to cover the case of polarized $`U(1)`$ symmetric vacuum solutions on manifolds of the more general type $`\mathrm{\Sigma }^2\times R\times S^1\mathrm{\Sigma }^2\times R,`$ where $`\mathrm{\Sigma }^2`$ is an arbitrary compact surface and the bundle (in view of the assumed polarization condition) is necessarily trivial. In the present paper the polarization restriction is eliminated in favor of an appropriate half polarization condition and the limitation to trivial $`S^1`$ bundles over the base $`\mathrm{\Sigma }^2\times R`$ is also removed. The present work thus demontrates the existence of a large family of vacuum $`U(1)`$ symmetric solutions of half polarized type defined on trivial and non trivial bundles over $`\mathrm{\Sigma }^2\times R`$ (with $`\mathrm{\Sigma }^2`$ an arbitrary compact surface) and having AVTD singularity behavior. The half polarization condition used in \[IM\] involved requiring one of the asymptotic functions to vanish. The half polarization condition which we find here necessary and sufficient for possible AVTD behavior can be understood in terms of the behavior of the VTD solutions to which our solutions converge as one approaches the singularity. Specifically a VTD solution is half polarized if and only if the set of geodesics in the Poincarรฉ plane which represent it (at different spatial points) all tend to the same point as $`t`$ appoaches the singularity.
## 2 Einstein equations.
A spacetime metric on a manifold $`V_4M\times R,`$ with $`M`$ an $`S^1`$ principal fiber bundle over a surface $`\mathrm{\Sigma },`$ reads, if it is invariant under the $`S^1`$ action on $`V_4,`$
$${}_{}{}^{(4)}ge^{2\varphi }\text{ }^{(3)}g+e^{2\varphi }(d\theta +A)^2$$
(2.1)
with $`\theta `$ a parameter on the (spacelike) circular orbit, $`\varphi `$ a scalar, $`A`$ a locally defined 1 - form and $`{}_{}{}^{(3)}g`$ a lorentzian metric, all on $`V_3:=\mathrm{\Sigma }\times R.`$
The vacuum 3+1 Einstein equations $`Ricci(^{(4)}g)=0`$ for such an $`S^1`$ symmetric metric on $`V_4`$ are known \[M86\], \[CB-M96\] to be equivalent<sup>2</sup><sup>2</sup>2If we choose an arbitrary harmonic 1 - form appearing in the solution to be zero. to the wave map equation from $`(V_3,^{(3)}g)`$ into the Poincarรฉ plane $`P=:(R^2,G),`$ $`\mathrm{\Phi }(\varphi ,\omega ):`$ $`V_3R^2,`$ where
$$G2(d\varphi )^2+\frac{1}{2}e^{4\gamma }(d\omega )^2,$$
(2.2)
coupled to the 2+1 Einstein equations for $`{}_{}{}^{(3)}g`$ on $`V_3`$ with the wave map as the source field. The scalar function $`\omega `$ on $`V_3`$ is linked to the differential $`F`$ of $`A`$ by the relation
$$d\omega =e^{4\varphi }F,\text{ with }F=dA.$$
(2.3)
Thus in local coordinates $`x^\alpha `$, $`\alpha =0,1,2,`$ on $`V_3,`$ with $`\eta `$ the volume form of $`{}_{}{}^{(3)}g=^{(3)}g_{\alpha \beta }dx^\alpha dx^\beta ,`$ one has
$$F_{\alpha \beta }\frac{1}{2}e^{4\varphi }\eta _{\alpha \beta \lambda }^\lambda \omega .$$
(2.4)
The wave map equations are, with $`{}_{}{}^{(3)}`$ the covariant derivative in the metric $`{}_{}{}^{(3)}g`$
$$g^{\alpha \beta }(^{(3)}_\alpha _\beta \varphi +\frac{1}{2}e^{4\varphi }_\alpha \omega _\beta \omega )=0$$
(2.5)
$$g^{\alpha \beta }(^{(3)}_\alpha _\beta \omega 4_\alpha \omega _\beta \varphi )=0.$$
(2.6)
The $`2+1`$ Einstein equations are, with โ.โ indicating a scalar product in the metric $`G`$
$${}_{}{}^{(3)}R_{\alpha \beta }^{}=_\alpha \mathrm{\Phi }._\beta \mathrm{\Phi }.$$
(2.7)
To solve these equations we choose for $`{}_{}{}^{(3)}g`$ a zero shift, we denote the lapse by $`e^\lambda `$ and we weigh by $`e^\lambda ,`$ without restricting the generality, the $`t`$ dependent space metric $`g=g_{ab}dx^adx^b,`$ $`a,b=1,2.`$ That is, we set
$${}_{}{}^{(3)}gN^2dt^2+g_{ab}dx^adx^b\text{ with }Ne^\lambda ,\text{ }g_{ab}e^\lambda \sigma _{ab},\text{ }$$
(2.8)
We denote by $`\sigma ^{ab}`$ the contravariant form of $`\sigma .`$ The extrinsic curvature of $`\mathrm{\Sigma }_t`$ in $`(V_3,^{(3)}g)`$ is
$$k_{ab}:=\frac{1}{2N}_tg_{ab}\frac{1}{2}(\sigma _{ab}_t\lambda +_t\sigma _{ab}).$$
(2.9)
The mean extrinsic curvature $`\tau `$ is therefore
$$\tau :=g^{ab}k_{ab}e^\lambda (_t\lambda +\frac{1}{2}\psi ),$$
(2.10)
where we have defined
$$\psi :=\sigma ^{ab}_t\sigma _{ab}.$$
(2.11)
The connection coefficients (Christoffel symbols) of $`{}_{}{}^{(3)}g`$ are found to be (note that $`{}_{}{}^{(3)}g_{}^{00}=e^{2\lambda },`$ $`{}_{}{}^{(3)}g_{}^{ab}=g^{ab}=e^\lambda \sigma ^{ab})`$
$${}_{}{}^{(3)}\mathrm{\Gamma }_{ab}^{c}=\mathrm{\Gamma }_{ab}^c(g)=\mathrm{\Gamma }_{ab}^c(\sigma )+\frac{1}{2}(\delta _b^c_a\lambda +\delta _a^c_b\lambda \sigma ^{cd}\sigma _{ab}_d\lambda )$$
(2.12)
$${}_{}{}^{(3)}\mathrm{\Gamma }_{00}^{0}=_t\lambda \text{}^{(3)}\mathrm{\Gamma }_{0a}^0=_a\lambda ,\text{ }^{(3)}\mathrm{\Gamma }_{00}^a=\sigma ^{ab}e^\lambda _a\lambda ,$$
(2.13)
$${}_{}{}^{(3)}\mathrm{\Gamma }_{ab}^{0}=e^\lambda k_{ab},\text{ }^{(3)}\mathrm{\Gamma }_{a0}^b=e^\lambda k_a^b.$$
(2.14)
In particular it holds that
$${}_{}{}^{(3)}g_{}^{\alpha \beta }{}_{}{}^{(3)}\mathrm{\Gamma }_{\alpha \beta }^{0}=\frac{1}{2}\psi e^{2\lambda }.$$
(2.15)
We see that the metric $`{}_{}{}^{(3)}g`$ is in harmonic time gauge if and only if $`\psi =0.`$
The Einstein equations split into constraints and evolution equations. We denote by $`S_\beta ^\alpha ^{(3)}R_\beta ^\alpha \frac{1}{2}\delta _\beta ^\alpha {}_{}{}^{(3)}R`$ the Einstein tensor of $`{}_{}{}^{(3)}g,`$ by $`T_\beta ^\alpha `$ the stress energy tensor of $`\mathrm{\Phi },`$ and we set $`\mathrm{\Sigma }_\beta ^\alpha S_\beta ^\alpha T_\beta ^\alpha .`$ The constraints are:
$$C_0\mathrm{\Sigma }_0^0\frac{1}{2}\{R(g)k.k+\tau ^2e^{2\lambda }_t\mathrm{\Phi }._t\mathrm{\Phi }g^{ab}_a\mathrm{\Phi }._b\mathrm{\Phi }\}=0$$
(2.16)
and (indices raised with $`g^{ab},`$ $``$ the covariant derivative in the metric $`g)`$
$$C_ae^\lambda \mathrm{\Sigma }_a^0\{_bk_a^b_a\tau +e^\lambda _a\mathrm{\Phi }._t\mathrm{\Phi }\}=0.$$
(2.17)
The evolution equations are, with $`N=e^\lambda ,`$
$$N(^{(3)}R_a^b\rho _a^b)_tk_a^b+N\tau k_a^b^b_aN+NR_a^bN_a\mathrm{\Phi }.^b\mathrm{\Phi }=0.$$
(2.18)
In order to obtain a first order system in the Fuchsian analysis that we will make, we introduce auxiliary unknowns $`\mathrm{\Phi }_t,\mathrm{\Phi }_a,\sigma _c^{ab}`$ which are identified with the first partial derivatives of $`\mathrm{\Phi }`$ and the covariant derivative of $`\sigma `$ with respect to a given $`t`$ independent metric $`\stackrel{~}{\sigma }.`$ These new unknowns satisfy the evolution equations
$$_t\mathrm{\Phi }=\mathrm{\Phi }_t,$$
(2.19)
$$\text{ }_t\mathrm{\Phi }_a=_a\mathrm{\Phi }_t,$$
(2.20)
$$_t\sigma _c^{ab}=\stackrel{~}{}_c_t\sigma ^{ab}\text{ }$$
(2.21)
where, by the definitions of $`\sigma `$ and $`k,`$
$$_t\sigma ^{ab}=2e^{2\lambda }k^{ab}+\sigma ^{ab}_t\lambda .$$
(2.22)
The function $`\lambda `$ is not left unknown, but rather is determined by a gauge condition from its VTD value.
## 3 VTD equations and solutions.
The Velocity Terms Dominated equations are obtained by dropping the space derivatives in the equations.
We denote by a tilde quantities which are independent of $`t,`$ and we denote VTD solutions using a hat.
### 3.1 Einstein evolution VTD solutions.
In order to obtain a global (on $`\mathrm{\Sigma }`$) formulation we choose a VTD metric which remains in a fixed conformal class over $`\mathrm{\Sigma }`$ as $`t`$ evolves; we set
$$\widehat{\sigma }_{ab}=\stackrel{~}{\sigma }_{ab}\text{ and }\widehat{g}_{ab}=e^{\widehat{\lambda }}\stackrel{~}{\sigma }_{ab}.$$
(3.1)
Then
$$\widehat{\psi }=0,\text{ }\widehat{g}^{ab}=e^{\widehat{\lambda }}\stackrel{~}{\sigma }^{ab},\text{ }_t\widehat{g}_{ab}=e^{\widehat{\lambda }}\stackrel{~}{\sigma }_{ab}_t\widehat{\lambda },$$
(3.2)
and the definition of $`k`$ gives that
$$\widehat{k}_{ab}=\frac{1}{2}\stackrel{~}{\sigma }_{ab}_t\widehat{\lambda },\text{ }\widehat{k}_a^b=\frac{1}{2}e^{\widehat{\lambda }}\delta _a^b_t\widehat{\lambda },\text{ }\widehat{\tau }:=\widehat{k}_a^a=e^{\widehat{\lambda }}_t\widehat{\lambda }.$$
(3.3)
Requiring that these VTD quantities satisfy the VTD evolution equations, we obtain
$$_t\widehat{k}_a^b=\widehat{N}\widehat{\tau }\widehat{k}_a^b.$$
(3.4)
Therefore, by straightforward computation
$$_{tt}^2\widehat{\lambda }=0;\text{ hence }\widehat{\lambda }=\stackrel{~}{\lambda }\stackrel{~}{v}t$$
(3.5)
with $`\stackrel{~}{\lambda }`$ and $`\stackrel{~}{v}`$ arbitrary functions on $`\mathrm{\Sigma },`$ independent of $`t.`$ Then we have
$$\widehat{k}_{ab}=\frac{1}{2}\stackrel{~}{v}\stackrel{~}{\sigma }_{ab},\text{ }\widehat{k}_a^b=\frac{1}{2}\stackrel{~}{v}e^{\widehat{\lambda }}\delta _a^b,\text{ }\widehat{\tau }=e^{\widehat{\lambda }}\stackrel{~}{v}.$$
(3.6)
### 3.2 Wave map VTD solutions.
The results for a VTD wave map are very different from the results obtained for a scalar function \[CBIM\]. If we drop space derivatives in the wave map equations we obtain geodesic equations in the target manifold, with $`t`$ the length parameter on these geodesics so long as the 2+1 metric is in harmonic time gauge. If we make the change of coordinates $`Y=e^{2\varphi }`$ in the target (which defines a diffeomorphism from $`R^2`$ onto the upper half plane $`Y>0),`$ the metric $`G`$ takes a standard form for the metric of a Poincarรฉ half plane; namely
$$G\frac{1}{2}\{\frac{d\omega ^2+dY^2}{Y^2}\},\text{ }Y=e^{2\varphi }.$$
(3.7)
The VTD, geodesic, equations written in this metric read, with a prime denoting the derivative with respect to $`t`$
$$\omega ^{\prime \prime }2Y^1\omega ^{}Y^{}=0,$$
(3.8)
$$Y^{\prime \prime }+Y^1\omega ^{}X^{}=0.$$
(3.9)
The general solution of these geodesic equations is represented in these coordinates, as is well known, by half circles<sup>3</sup><sup>3</sup>3We discard here the special case which corresponds to the polarized case, treated elsewhere, where these circles are centered at infinity. The geodesics are then the half lines X$`\omega =`$constant,. centered on the line $`Y=0;`$ specifically, with $`A`$ and $`B`$ arbitrary constants (that is, independent of $`t)`$, the solution takes the form
$$\widehat{\omega }=B+A\mathrm{cos}\theta ,\text{ }\widehat{Y}=A\mathrm{sin}\theta ,\text{ }0<\theta <\pi .$$
(3.10)
These functions $`\omega `$ and $`Y`$ satisfy the differential equations 3.8, 3.9 if and only if it holds that:
$$\frac{\theta ^{\prime \prime }}{\theta ^{}}=\frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }\theta ^{}.$$
(3.11)
Integrating this equation we have that, with $`\stackrel{~}{w}`$ independent of $`t`$
$$\theta ^{}=\stackrel{~}{w}\mathrm{sin}\theta .$$
(3.12)
Another integration gives that, with $`\stackrel{~}{\mathrm{\Theta }}`$ independent of $`t`$
$$\mathrm{tan}\frac{\theta }{2}=\stackrel{~}{\mathrm{\Theta }}e^{\stackrel{~}{w}t}.$$
(3.13)
If we now make the substitution $`A=e^{2\stackrel{~}{\varphi }}`$ and $`B=\stackrel{~}{\omega },`$ then 3.10 reads
$$\widehat{\varphi }=\stackrel{~}{\varphi }+\frac{1}{2}log(sin\theta ),\text{ }\widehat{\omega }=\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}cos\theta \text{ }.\text{ }$$
(3.14)
###### Remark 3.1
The set of above formulas is identical to the following one
$$\widehat{\omega }\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}\frac{1\stackrel{~}{\mathrm{\Theta }}^2e^{2\stackrel{~}{w}t}}{1+\stackrel{~}{\mathrm{\Theta }}^2e^{2\stackrel{~}{w}t}},\text{ }e^{2\widehat{\varphi }}\widehat{Y}=e^{2\stackrel{~}{\varphi }}\frac{2\stackrel{~}{\mathrm{\Theta }}e^{\stackrel{~}{w}t}}{1+\stackrel{~}{\mathrm{\Theta }}^2e^{2\stackrel{~}{w}t}}.$$
(3.15)
$`\widehat{Y}`$ tends to zero when $`t`$ tends to $`\mathrm{},`$ but $`\widehat{\omega }`$ tends to $`\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}.`$
### 3.3 Einstein constraint VTD solutions.
We deduce from 3.14 and 3.12 that
$$_t\widehat{\varphi }\frac{1}{2}\widehat{Y}^1\widehat{Y}^{}=\frac{1}{2}\frac{cos\theta \theta ^{}}{sin\theta }=\frac{1}{2}\stackrel{~}{w}cos\theta $$
(3.16)
and
$$e^{2\widehat{\varphi }}_t\widehat{\omega }=\theta ^{}=\stackrel{~}{w}sin\theta .$$
(3.17)
The Einstein VTD constraints reduce to the following equation
$$2\widehat{C}_0\widehat{k}.\widehat{k}+\widehat{\tau }^2e^{2\widehat{\lambda }}\{2(_t\widehat{\varphi })^2+\frac{1}{2}e^{4\widehat{\varphi }}(_t\widehat{\omega })^2\}=0.$$
(3.18)
We have, using 3.16 and 3.17,
$$2(_t\widehat{\varphi })^2+\frac{1}{2}e^{4\widehat{\varphi }}(_t\widehat{\omega })^2=\frac{1}{2}\stackrel{~}{w}^2.$$
(3.19)
We deduce therefore from 3.6 that the VTD constraint 3.18 is satisfied if and only if
$$\stackrel{~}{v}^2=\stackrel{~}{w}^2$$
(3.20)
## 4 Fuchsian expansion.
### 4.1 2+1 metric expansions.
For the unknowns $`\sigma `$ and $`k`$ we choose the following expansions, with the various $`\epsilon ^{}s`$ being positive numbers to be chosen later
$$\sigma ^{ab}=\stackrel{~}{\sigma }^{ab}+e^{\epsilon _\sigma t}u_\sigma ^{ab}$$
(4.1)
$$k_a^b=e^\lambda (\frac{1}{2}\stackrel{~}{v}\delta _a^b+e^{\epsilon _kt}u_{k,a}^b).$$
(4.2)
Then
$$\tau k_a^a=e^\lambda (\stackrel{~}{v}+e^{\epsilon _kt}u_{k,a}^a).$$
(4.3)
We take as a gauge condition
$$\lambda =\widehat{\lambda };\text{ hence }_t\lambda =\stackrel{~}{v}.$$
(4.4)
Comparing the expressions 4.4 and 2.10 for $`\tau ,`$ we find that this condition is equivalent to the gauge fixing requirement
$$e^{\epsilon _kt}u_{k,a}^a+\frac{1}{2}\psi =0.$$
(4.5)
### 4.2 Wave map expansion.
We expand $`\mathrm{\Phi }`$ near its VTD value; that is we set
$$\varphi =\widehat{\varphi }+e^{\epsilon _\varphi t}u_\varphi \text{ with }\widehat{\varphi }=\stackrel{~}{\varphi }+\frac{1}{2}log(sin\theta ),$$
(4.6)
while for $`\omega ,`$ for convenience of computation, we choose to set
$$\omega =\widehat{\omega }+e^{2\varphi }e^{\epsilon _\omega t}u_\omega \text{ with }\widehat{\omega }=\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}cos\theta ,$$
(4.7)
### 4.3 Expansion for first derivatives.
We expand the auxiliary unknowns near the values of the derivatives of the VTD solution. That is we set (see 3.16, 3.17)
$$\varphi _t=_t\widehat{\varphi }+e^{\epsilon _{\varphi _t}t}u_{\varphi _t}\frac{1}{2}\stackrel{~}{w}cos\theta +e^{\epsilon _{\varphi _t}t}u_{\varphi _t},$$
(4.8)
$$\omega _t=_t\widehat{\omega }+e^{2\varphi }e^{\epsilon _{\omega _t}}u_{\omega _t}e^{2\stackrel{~}{\varphi }}\stackrel{~}{w}sin^2\theta +e^{2\varphi }e^{\epsilon _{\omega _t}}u_{\omega _t}$$
(4.9)
The expansions of $`\varphi _a`$ and $`\omega _a`$ are defined similarly by setting
$$\varphi _a=_a\widehat{\varphi }+e^{\epsilon _\varphi ^{}}u_{\varphi _a},\text{ }\omega _a=_a\widehat{\omega }+e^{2\varphi }e^{\epsilon _\omega ^{}t}u_{\omega _a}.$$
(4.10)
We next compute $`_a\widehat{\varphi }`$ and $`_a\widehat{\omega }.`$ It follows from 3.14 that
$$_a\widehat{\varphi }=_a\stackrel{~}{\varphi }+\frac{cos\theta }{2sin\theta }_a\theta ,\text{ }_a\widehat{\omega }=_a\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}(2cos\theta _a\stackrel{~}{\varphi }sin\theta _a\theta ).$$
(4.11)
We compute $`_a\theta `$ using 3.13 and elementary properties of sine and cosine. We find that
$$_a\theta =\stackrel{~}{\mathrm{\Theta }}^1sin\theta _a(\stackrel{~}{\mathrm{\Theta }}t\stackrel{~}{w})$$
(4.12)
Therefore it holds that
$$_a\widehat{\varphi }=_a\stackrel{~}{\varphi }+\stackrel{~}{\mathrm{\Theta }}^1\frac{cos\theta }{2}_a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}t).$$
(4.13)
Then, writing $`_a\widehat{\omega }`$ as sum of a term independent of $`t`$ plus terms tending to zero when $`t`$ tends to infinity, we have
$$_a\widehat{\omega }_a(\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }})e^{2\stackrel{~}{\varphi }}[2(1cos\theta )_a\stackrel{~}{\varphi }\stackrel{~}{\mathrm{\Theta }}^1sin^2\theta _a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}t)]$$
(4.14)
For $`\sigma _c^{ab},`$ since $`\stackrel{~}{}_c\stackrel{~}{\sigma }^{ab}=0,`$ we set
$$\sigma _c^{ab}e^{\epsilon _\sigma ^{}t}u_{\sigma ^{},c}^{ab}.$$
(4.15)
## 5 Fuchsian system for the evolution equations.
Given the Fuchsian expansions of the previous section, the Einstein - wave map evolution system reads as a first order system for the set of unknowns $`U(u_\sigma ,`$ $`u_k,`$ $`u_\mathrm{\Phi },`$ $`u_{\mathrm{\Phi }_t},`$ $`u_{\mathrm{\Phi }^{},\text{ }}u_\sigma ^{}).`$
The differential system for $`U`$ is Fuchsian in a neighbourhood of $`t=+\mathrm{}`$ if it takes the form
$$_tULU=e^{\mu t}F(t,x,U,\stackrel{~}{}U)$$
(5.1)
with $`L`$ a linear operator independent of $`t`$ with non negative eigenvalues, $`\mu `$ a positive number and $`F`$ a set of tensor fields linear in $`\stackrel{~}{}U`$, continuous in $`t,`$ analytic in $`x`$ and $`U`$ and uniformly Lipshitzian in all its arguments in a neighbourhood of $`U=0`$, for $`t`$ large enough.
### 5.1 Einstein evolution equations.
#### 5.1.1 Equation for $`u_\sigma .`$
The Fuchsian expansion 4.2 for $`k`$ yields the following equation:
$$_tg^{ab}2Nk^{ab}2e^\lambda g^{ac}k_c^be^\lambda (\stackrel{~}{v}\sigma ^{ab}+2e^{\epsilon _kt}\sigma ^{ac}u_{k,c}^b).$$
(5.2)
Using $`g^{ab}e^\lambda \sigma ^{ab}`$ and $`_t\lambda =\stackrel{~}{v},`$ we have
$$_tg^{ab}e^\lambda (\stackrel{~}{v}\sigma ^{ab}+_t\sigma ^{ab}).$$
(5.3)
Combining these equations together with the Fuchsian expansion of $`\sigma `$ results in the equation:
$$_tu_\sigma ^{ab}\epsilon _\sigma u_\sigma ^{ab}=2e^{(\epsilon _\sigma \epsilon _k)t}\sigma ^{ac}u_{k,c}^b,$$
(5.4)
which is of Fuchsian type if $`\epsilon _k>\epsilon _\sigma >0.`$
#### 5.1.2 Equation for $`u_k.`$
The Fuchsian expansion of $`k`$ together with $`N=e^\lambda `$ and $`_t\lambda =\stackrel{~}{v}`$ imply by straightforward computation that
$$_tk_a^be^\lambda \{\frac{1}{2}\stackrel{~}{v}^2\delta _a^b+e^{\epsilon _kt}(\stackrel{~}{v}\epsilon _k)u_{k,a}^b+e^{\epsilon _kt}_tu_{k,a}^b\},$$
(5.5)
and
$$N\tau k_a^be^\lambda \{\frac{1}{2}\stackrel{~}{v}^2\delta _a^b+\stackrel{~}{v}e^{\epsilon _kt}u_{k,a}^b)+e^{\epsilon _kt}u_{k,c}^c(\frac{1}{2}\stackrel{~}{v}\delta _a^b+e^{\epsilon _kt}u_{k,a}^b)\}.$$
(5.6)
We see that $`e^\lambda \stackrel{~}{v}^2`$ disappears from the difference $`_tk_a^bN\tau k_a^b,`$ which motivates the choice of the Fuchsian expansion.
To write the evolution equation 2.18 for $`k`$ we now compute
$$^b_aNe^\lambda \sigma ^{bc}_c_ae^\lambda \sigma ^{bc}[_c\lambda _a\lambda +_a_c\lambda \mathrm{\Gamma }_{ac}^d(g)_d\lambda ].$$
(5.7)
On the other hand, since $`\mathrm{\Sigma }`$ is 2 dimensional and $`g`$ is conformal to $`\sigma `$ with a factor $`e^\lambda ,`$ we have that
$$NR_a^be^\lambda R_a^b\frac{1}{2}e^\lambda \delta _a^bR(g)=\frac{1}{2}\delta _a^b\{R(\sigma )\mathrm{\Delta }_\sigma \lambda \}.$$
(5.8)
From these results, if we define
$$f_a^b(t,u,u_x):=^b_aN+NR_a^bN_a\mathrm{\Phi }.^b\mathrm{\Phi }$$
then we calculate
$$f_a^b\sigma ^{bc}[_c\lambda _a\lambda +_a_c\lambda \mathrm{\Gamma }_{ac}^d(g)_d\lambda ]+\frac{1}{2}\delta _a^b[R(\sigma )\mathrm{\Delta }_\sigma \lambda ]\sigma ^{bc}\mathrm{\Phi }_a.\mathrm{\Phi }_c$$
(5.9)
We see that $`f_a^b`$ is at most a second order polynomial in $`t,`$ is analytic in $`x`$ when $`\stackrel{~}{v},\stackrel{~}{w},\stackrel{~}{\lambda },\stackrel{~}{\sigma }`$ are analytic; is linear in $`u;`$ and is analytic, bounded and Lipshitzian in $`u`$ for $`u`$ bounded and for large<sup>4</sup><sup>4</sup>4This restriction on t comes from the covariant components of $`\sigma `$ which remain bounded as long as $`\sigma ^{ab}`$ remains positive definite. $`t`$, except eventually for the last term which reads
$$\sigma ^{bc}\mathrm{\Phi }_a.\mathrm{\Phi }_c=2\sigma ^{bc}(\varphi _a\varphi _c+\frac{1}{2}e^{4\varphi }\omega _a\omega _c).$$
(5.10)
The expansion 4.10 of $`\varphi _a`$ shows that it does not cause problems for the boundedness of $`f_a^b`$. However the expansion of $`\omega _a`$ gives
$$e^{2\varphi }\omega _a=e^{2\varphi }_a(\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }})e^{2\varphi +2\stackrel{~}{\varphi }}[2(1cos\theta )_a\stackrel{~}{\varphi }+\stackrel{~}{\mathrm{\Theta }}^1sin^2\theta _a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}t)]$$
$$+e^{\epsilon _\omega ^{}t}u_{\omega _a}.$$
(5.11)
It follows from 4.6 that
$$e^{2(\stackrel{~}{\varphi }\varphi )}=\frac{e^{2\delta \varphi }}{sin\theta },\text{ with }\delta \varphi e^{\epsilon _\varphi t}u_\varphi .$$
Therefore, using ($`1cos\theta )/sin\theta =tan(\theta /2)`$ we have:
$$e^{2\varphi }\omega _a=e^{2\stackrel{~}{\varphi }}\frac{e^{2\delta \varphi }}{sin\theta }_a(\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }})e^{2\delta \varphi }[2tan\frac{\theta }{2}_a\stackrel{~}{\varphi }+\stackrel{~}{\mathrm{\Theta }}^1sin\theta _a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}t)]+e^{\epsilon _\omega ^{}t}u_{\omega _a}.$$
(5.12)
We see that $`e^{2\varphi }\omega _a`$ will increase like ($`sin\theta )^1`$ \- that is, like $`e^{\stackrel{~}{w}t}`$ \- as $`t`$ tends to infinity, except if
$$\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}=constant.$$
(5.13)
Condition 5.13 is a generalization of the condition imposed on the fields in \[IM\], with other notations, to obtain AVTD behaviour, in the case that $`\mathrm{\Sigma }`$ is a torus. Following the terminology of \[IM\] we call equation 5.13 the โhalf polarizationโ condition. Its geometric meaning is that the set of geodesics in the Poincarรฉ plane representing the VTD solution all tend to the same point of the axis $`Y=0`$ as $`t`$ tends to infinity.
After inserting the Fuchsian expansions and multiplying by $`e^{\lambda +\epsilon _k:t}`$ we find that the equation 2.18 takes the form
$$_tu_{k,a}^b\epsilon _ku_{k,a}^b\frac{1}{2}v\delta _a^bu_{k,c}^c=e^{\epsilon _kt}u_{k,c}^cu_{k,a}^b+e^{\lambda +\epsilon _kt}f_a^b(t,u,u_x).$$
(5.14)
Since $`\lambda =\stackrel{~}{\lambda }\stackrel{~}{v}t`$ and $`\stackrel{~}{v}=\stackrel{~}{w}`$ this system can take a Fuchsian form only if the functions $`\stackrel{~}{\omega }`$ and $`\stackrel{~}{\varphi }`$ satisfy the half polarization condition 5.13. To obtain the system in obviously Fuchsian form in that case, we split 5.14 into its trace and its traceless parts. For the trace part we have
$$_tu_{k,a}^a\epsilon _ku_{k,a}^a\stackrel{~}{v}u_{k,a}^a=e^{\epsilon _kt}u_{k,c}^cu_{k,a}^a+e^{\lambda +\epsilon _kt}f_a^a(t,u,u_x).$$
(5.15)
This equation takes Fuchsian form if and only if 5.13 is satisfied and $`\stackrel{~}{v}>\epsilon _k.`$ The same is verified for the traceless part $`{}_{}{}^{T}u_{k,a}^{b},`$ which satisfies an equation with left hand side
$$_t\text{ }^Tu_{k,a}^b\epsilon _k{}_{}{}^{T}u_{k,a}^{b}.$$
(5.16)
#### 5.1.3 Equation for $`u_\sigma ^{}.`$
Using the expansion of $`k`$ and the relation $`_t\lambda =\stackrel{~}{v},`$ we find that
$$_t\sigma ^{ab}=2e^{2\lambda }k^{ab}+\sigma ^{ab}_t\lambda =2e^{\epsilon _kt}\sigma ^{ac}u_{k,c}^b$$
(5.17)
The equation for $`\sigma _c^{ab}`$ gives therefore the following equation for $`u_{\sigma ^{}\text{ }}:`$
$$_tu_{\sigma ^{},c}^{ab}\epsilon _\sigma ^{}u_{\sigma ^{},c}^{ab}=2e^{(\epsilon _\sigma ^{}\epsilon _k)t}\stackrel{~}{}_c[\sigma ^{ac}u_{k,c}^b.]\text{ }$$
(5.18)
which is of Fuchsian type so long as $`\epsilon _\sigma ^{}<\epsilon _k.`$
### 5.2 Wave map equations.
#### 5.2.1 Equations for auxiliary variables.
The equations resulting from the introduction of the new variables $`\varphi _t,\omega _t`$ are
$$_t\varphi \varphi _t=0,\text{ }_t\omega \omega _t=0.$$
(5.19)
The first equation is of Fuchsian type for $`u_\varphi `$ if $`\epsilon _{\mathrm{\Phi }_t}>\epsilon _{\mathrm{\Phi }\text{ }},`$ since it reads
$$_tu_\varphi \epsilon _\varphi u_\varphi =e^{(\epsilon _{\mathrm{\Phi }_t}+\epsilon _\mathrm{\Phi })t}u_{\varphi _t}.$$
(5.20)
The second equation reads
$$[_tu_\omega +(2\varphi _t\epsilon _\omega )u_\omega ]\text{ }e^{(\epsilon _\mathrm{\Phi }\epsilon _{\mathrm{\Phi }_t})t}u_{\omega _t}=0.$$
(5.21)
We replace $`\varphi _t`$ by its value given in 4.8, which we write as follows
$$\varphi _t=\frac{1}{2}\stackrel{~}{w}+\frac{1}{2}\stackrel{~}{w}(1cos\theta )+e^{\epsilon _{\mathrm{\Phi }_t}t}u_{\varphi _t}.$$
(5.22)
Since $`1cos\theta `$ falls off to zero as $`e^{2\stackrel{~}{w}t},`$ the equation 5.21 is of Fuchsian type for $`u_\omega `$ if $`\stackrel{~}{w}>0`$ and $`\epsilon _{\mathrm{\Phi }_t}>\epsilon _\mathrm{\Phi }.`$
In the equations 2.20 to be satisfied by $`\varphi _a`$ and $`\omega _a,`$ the derivatives of the VTD terms disappear, due to the commutation of partial derivatives. The equation for $`\varphi _a`$ reads
$$_tu_{\varphi _a}\epsilon _\varphi ^{}u_{\varphi _a}=e^{(\epsilon _{\mathrm{\Phi }_t}\epsilon _\mathrm{\Phi }^{})t}(_au_{\varphi _t}t_aw),$$
(5.23)
while the equation for $`\omega _a`$ becomes, using the expressions for $`\omega _t`$ and $`\omega _a`$
$$_tu_{\omega _a}+(2\varphi _t\epsilon _\omega ^{})u_{\omega _a}=e^{(\epsilon _{\mathrm{\Phi }_t}\epsilon _\mathrm{\Phi }^{})t}(_au_{\omega _t}+2\varphi _au_{\omega _t}).$$
(5.24)
These equations are of Fuchsian type so long as $`\stackrel{~}{w}>0`$ and $`\epsilon _{\mathrm{\Phi }_t}>\epsilon _\mathrm{\Phi }^{}.`$
#### 5.2.2 Equation for $`u_{\mathrm{\Phi }_t}.`$
The first equation, 2.5, for the wave map reads
$$g^{\alpha \beta }(_\alpha _\beta \varphi +\frac{1}{2}e^{4\varphi }_\alpha \omega _\beta \omega )$$
$$e^{2\lambda }(_t\varphi _t+\frac{1}{2}e^{4\varphi }\omega _t\omega _t)+e^\lambda \sigma ^{ab}(_a\varphi _b+\frac{1}{2}e^{4\varphi }\omega _a\omega _b)+g^{\alpha \beta }\mathrm{\Gamma }_{\alpha \beta }^0\varphi _t=0$$
(5.25)
Using the Fuchsian expansions for $`\varphi _t`$ and $`\omega _t`$ together with $`\theta ^{}=\stackrel{~}{w}\mathrm{sin}\theta `$ and the value given in section 5.1.2 for $`e^{2(\stackrel{~}{\varphi }\varphi )}`$we find that:
$$_t\varphi _t+\frac{1}{2}e^{4\varphi }\omega _t\omega _t$$
$$e^{\epsilon _{\mathrm{\Phi }_t}t}(_tu_{\varphi _t}\epsilon _{\varphi _t}u_{\varphi _t})\frac{1}{2}\stackrel{~}{w}^2\mathrm{sin}^2\theta +\frac{1}{2}(e^{2\delta \varphi }\stackrel{~}{w}sin\theta +e^{\epsilon _{\omega _t}t}u_{\omega _t})^2$$
On the other hand, using the expansions for $`\sigma ^{ab},\varphi _a`$ and $`\omega _a`$ we find that:
$$e^\lambda \sigma ^{ab}(_a\varphi _b+\frac{1}{2}e^{4\varphi }\omega _a\omega _b)$$
$$e^\lambda (\stackrel{~}{\sigma }^{ab}+\delta \sigma ^{ab})\{_b_a\widehat{\varphi }+e^{\epsilon _\mathrm{\Phi }^{}t}_bu_{\varphi _a}+\frac{1}{2}(e^{2\varphi }_a\widehat{\omega }+e^{\epsilon _\omega ^{}t}u_{\omega _a})(e^{2\varphi }_b\widehat{\omega }+e^{\epsilon _\omega ^{}t}u_{\omega _b})\}$$
We recall that
$$_b_a\widehat{\varphi }_b[_a\stackrel{~}{\varphi }+\frac{cos\theta }{2}\stackrel{~}{\mathrm{\Theta }}^1_a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}t)],$$
while under the half polarization assumption
$$\stackrel{~}{\omega }+e^{2\stackrel{~}{\varphi }}=constant,$$
(5.26)
the product $`e^{2\varphi }_a\widehat{\omega }`$ is given by
$$e^{2\varphi }_a\widehat{\omega }=e^{2\delta \varphi }\stackrel{~}{\mathrm{\Theta }}^1sin\theta _a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}).$$
Finally we calculate
$$g^{\alpha \beta }\mathrm{\Gamma }_{\alpha \beta }^0\frac{1}{2}\psi e^{2\lambda }=e^{2\lambda \epsilon _kt}u_{k,a}^a.$$
(5.27)
Inserting these computations into the first wave map equation produces an equation of the form
$$_tu_{\varphi _t}\epsilon _{\mathrm{\Phi }_t}u_{\varphi _t}=e^{\mu t}f_{\varphi _t}(x,t,u,u)$$
(5.28)
which is of the Fuchsian type 5.1 (with $`\mu >0)`$ so long as $`\stackrel{~}{v}>\epsilon _{\mathrm{\Phi }_t},`$ and $`\epsilon _k>\epsilon _{\mathrm{\Phi }_t}.`$
Analogous computations show that the equation for $`u_{\omega _t}`$ is Fuchsian presuming these same inequalities hold.
### 5.3 Results for evolution.
As a consequence of the calculations above we have proven the following theorem.
###### Theorem 5.1
There exist a collection of positive numbers {$`\epsilon _\sigma ,\epsilon _\sigma ^{},\epsilon _k,\epsilon _\mathrm{\Phi },\epsilon _{\mathrm{\Phi }_t},\epsilon _\mathrm{\Phi }^{}\}`$ such that, given analytic asymptotic data on $`\mathrm{\Sigma },`$ $`\stackrel{~}{A}=\{\stackrel{~}{v}=\stackrel{~}{w},`$ $`\stackrel{~}{\lambda },`$ $`\stackrel{~}{\sigma },\stackrel{~}{\mathrm{\Theta }},\stackrel{~}{\varphi },\stackrel{~}{\omega }\},`$ the Einstein - wave map evolution system written in first order form for the unknown $`U,`$ which defines $`g,`$ $`k,`$ $`\mathrm{\Phi }`$ and auxiliary variables by the Fuchsian expansions of section 4, is a Fuchsian system for $`U`$ if and only if $`\stackrel{~}{\varphi }`$ and $`\stackrel{~}{\omega }`$ satisfy the half polarisation condition 5.13 and $`\stackrel{~}{v}>0.`$ It admits then one and only one analytic solution tending to zero at infinity.
To show that this result implies that we have a family of solutions of the Einstein - wave map evolution system which decays to solutions of the VTD equations, we need to verify that for a large enough $`t`$ we have $`\mathrm{\Phi }_t=_t\mathrm{\Phi },`$ $`\mathrm{\Phi }_a=_a\mathrm{\Phi }`$ and the like. To show that $`\mathrm{\Phi }_a=_a\mathrm{\Phi }`$ we use the equations 2.20 together with commutation of partial derivatives to show that:
$$_t(\varphi _a_a\varphi )=_a\varphi _t_a_t\varphi =0;$$
(5.29)
hence $`\varphi _a_a\varphi `$ is independent of $`t.`$ As $`t`$ tends to $`\mathrm{}`$ it tends to zero because
$$\varphi _a_a\varphi =e^{\epsilon _\varphi ^{}t}u_{\varphi _a}e^{\epsilon _\varphi t}(_au_\varphi \epsilon _\varphi u_\varphi ).$$
(5.30)
It must therefore always be zero. Analogous arguments can be used to show that $`\omega _a=_a\omega `$ and $`\sigma _c^{ab}=\stackrel{~}{}_c\sigma ^{ab}.`$
## 6 Constraints.
The solution of the evolution system satisfies the full Einstein equations so long as it satisfies also the Einstein constraints, that is
$$C_0:=\mathrm{\Sigma }_0^0\frac{1}{2}\{R(g)k.k+\tau ^2e^{2\lambda }_t\mathrm{\Phi }._t\mathrm{\Phi }\}=0$$
$$C_a:=e^\lambda \mathrm{\Sigma }_a^0\{_bk_a^b_a\tau +e^\lambda _t\mathrm{\Phi }._a\mathrm{\Phi }\}=0.$$
As usual we will rely on the Bianchi identities, here to construct a Fuchsian system satisfied by the constraints. Together with the wave equation satisfied by $`\mathrm{\Phi },`$ the Bianchi identities imply that
$${}_{}{}^{(3)}_{\alpha }^{}\mathrm{\Sigma }_\beta ^\alpha =0.\text{ }$$
(6.1)
Modulo the evolution equations $`{}_{}{}^{(3)}R_{a}^{b}\rho _a^b=0`$ that we have solved, with $`\rho _a^b\mathrm{\Phi }_a.\mathrm{\Phi }^b,`$ it holds that
$${}_{}{}^{(3)}R\rho =R_0^0\rho _0^0;$$
(6.2)
hence
$$\mathrm{\Sigma }_0^0R_0^0\rho _0^0\frac{1}{2}\delta _0^0(^{(3)}R\rho )=\frac{1}{2}\delta _0^0(^{(3)}R\rho )$$
(6.3)
and
$$\mathrm{\Sigma }_a^b=\frac{1}{2}\delta _a^b(^{(3)}R\rho )=\delta _a^b\mathrm{\Sigma }_0^0$$
(6.4)
We use these equations and the identities
$$\mathrm{\Sigma }_a^0e^\lambda C_a,\text{ }\mathrm{\Sigma }_0^aN^2\mathrm{\Sigma }^{a0}g^{ab}N^2\mathrm{\Sigma }_b^0e^\lambda g^{ab}C_b$$
(6.5)
together with the expressions for the Christoffel symbols of the metric $`{}_{}{}^{(3)}g.`$ We find that the equations 6.1 can be written in the form
$$_tC_02e^\lambda \tau C_0=g^{ab}_a(e^\lambda C_b)+g^{ab}e^\lambda _a\lambda C_a$$
(6.6)
and (after some simplifications and multiplying by $`e^\lambda )`$
$$_tC_ae^\lambda \tau C_a=e^\lambda _aC_0+2e^\lambda _a\lambda C_0.$$
(6.7)
Equivalently, we have
$$_t(e^{2\lambda }C_0)2\stackrel{~}{v}e^{2\lambda }C_02e^\lambda \tau e^{2\lambda }C_0=e^\lambda \sigma ^{ab}_a(e^\lambda C_b)+\sigma ^{ab}e^\lambda _a\lambda e^\lambda C_a$$
(6.8)
and
$$_t(e^\lambda C_a)\stackrel{~}{v}e^\lambda C_ae^{2\lambda }\tau C_a=_a(e^{2\lambda }C_0).$$
(6.9)
We see that $`e^{2\lambda }C_0`$ and $`e^\lambda C_a`$ satisfy a linear homogeneous system, which admits zero as a solution. This solution is the unique one tending to zero at infinity, so long as the system is Fuchsian.
###### Lemma 6.1
The system 6.8, 6.9 is Fuchsian, for a solution of the evolution system, if the VTD solution satisfies $`\widehat{C}_0=0`$ (i.e. $`\stackrel{~}{v}^2=\stackrel{~}{w}^2).`$
Proof. Since the coefficients of the equations 6.8, 6.9 are constructed from solutions of the evolution system we may use the expansions and estimates derived in previous sections. In particular we calculate
$$e^\lambda \tau \stackrel{~}{v}+e^{\epsilon _kt}u_{k,a}^a.$$
(6.10)
Equation 6.8 can therefore be written as the following equation of Fuchsian type:
$$_t(e^{2\lambda }C_0)4\stackrel{~}{v}e^{2\lambda }C_0e^{2\lambda }C_0=e^{\epsilon _kt}u_{k,a}^ae^{2\lambda }C_0+e^\lambda \sigma ^{ab}_a(e^\lambda C_b)+\sigma ^{ab}e^\lambda _a\lambda e^\lambda C_a.$$
(6.11)
Equation 6.9 is not a priori in Fuchsian form for the pair ($`e^\lambda C_a,e^{2\lambda }C_0)`$ in spite of the identity 6.10. However if we use the identity
$$e^{2\lambda }C_0\frac{1}{2}\{e^{2\lambda }R(g)e^{2\lambda }k.k+e^{2\lambda }\tau ^2_t\mathrm{\Phi }._t\mathrm{\Phi }\}$$
(6.12)
and the property $`\widehat{C}_0=0`$ together with the expression for $`R(g)`$ given in 5.9 we can show that there exists a number $`\mu >0`$ and a bounded function $`F(x,t)`$ such that we have
$$|e^{2\lambda }C_0|e^{\mu t}F(x,t)\text{ and }|_a(e^{2\lambda }C_0)|e^{\mu t}F(x,t).$$
(6.13)
It follows that 6.9 takes Fuchsian form.
###### Theorem 6.2
A solution of the evolution system satisfies the full Einstein wave map equations if and only if the half polarized asymptotic data satisfies the condition $`\stackrel{~}{w}=\stackrel{~}{v},`$ and also
$$\stackrel{~}{\mathrm{\Theta }}=1\text{ and }\stackrel{~}{v}e^{\stackrel{~}{\lambda }+2\stackrel{~}{\varphi }}=constant.$$
(6.14)
Proof. To complete the proof that $`C_0=C_a=0`$ it suffices to show that $`e^{2\lambda }C_0`$ and $`e^\lambda C_a`$ tend to zero at infinity. We have already checked that this is true for $`e^{2\lambda }C_0,`$ as long as $`\stackrel{~}{w}=\stackrel{~}{v};`$ i.e. $`\widehat{C}_0=0.`$
We now study the asymptotic behaviour of $`e^\lambda C_a.`$ If we denote by $`\delta u`$ the difference between a field $`u`$ and its VTD value, we calculate (recall that $`\lambda =\widehat{\lambda },\widehat{\sigma }=\stackrel{~}{\sigma })`$
$$e^\lambda (C_a\widehat{C}_a)e^\lambda \{(_b\stackrel{~}{}_b)k_a^b+\stackrel{~}{}_b\delta k_a^b_a\delta \tau \}+\delta (\mathrm{\Phi }_t.\mathrm{\Phi }_a)$$
(6.15)
with
$$\delta (\mathrm{\Phi }_t.\mathrm{\Phi }_a)2\varphi _t\delta \varphi _a+2\widehat{\varphi }_a\delta \varphi _t+\frac{1}{2}e^{2\varphi }\omega _t\delta (e^{2\varphi }\omega _a)+e^{2\widehat{\varphi }}\widehat{\omega }_a\delta (e^{2\varphi }\omega _t).$$
(6.16)
We see that, in the half polarized case, the Fuchsian expansions imply that $`e^\lambda (C_a\widehat{C}_a)`$ tends to zero as $`t`$ tends to infinity. Using the expressions for $`\widehat{k}_a^b`$ and $`\widehat{\lambda },`$ we see that $`e^{\widehat{\lambda }}\widehat{C}_a`$ reads:
$$e^{\widehat{\lambda }}\widehat{C}_a\frac{1}{2}e^{\widehat{\lambda }}_a(e^{\widehat{\lambda }}\stackrel{~}{v})\widehat{\mathrm{\Phi }}_t.\widehat{\mathrm{\Phi }}_a\frac{1}{2}(_a\stackrel{~}{v}\stackrel{~}{v}_a\stackrel{~}{\lambda }+\stackrel{~}{v}_a\stackrel{~}{v}t)\widehat{\mathrm{\Phi }}_t.\widehat{\mathrm{\Phi }}_a$$
Using the expressions of $`\widehat{\lambda },\widehat{\mathrm{\Phi }}_t,\widehat{\mathrm{\Phi }}_a`$ and the half polarization condition, we find after some computation that
$$\widehat{\mathrm{\Phi }}_t.\widehat{\mathrm{\Phi }}_a=\stackrel{~}{w}\{cos\theta _a\stackrel{~}{\varphi }+\frac{1}{2}\stackrel{~}{\mathrm{\Theta }}^1_a(\stackrel{~}{\mathrm{\Theta }}\stackrel{~}{w}t)\}.$$
(6.17)
Thus we find that the terms containing $`t`$ disappear from $`e^{\widehat{\lambda }}\widehat{C}_a`$ if $`\stackrel{~}{v}=\stackrel{~}{w}`$ and $`\stackrel{~}{\mathrm{\Theta }}=1.`$ It follows that $`e^{\widehat{\lambda }}\widehat{C}_a`$ tends to zero as $`t`$ tends to infinity (recall that $`cos\theta `$ tends then to $`1)`$ if and only if
$$\frac{1}{2}[_a\stackrel{~}{v}\stackrel{~}{v}_a\stackrel{~}{\lambda }]\stackrel{~}{v}_a\stackrel{~}{\varphi }=0,$$
a condition equivalent to the hypothesis 6.14 of the theorem.
###### Remark 6.3
In the half polarized case the VTD solution only satisfies asymptotically the VTD momentum constraint, and only after being multiplied $`e^{\widehat{\lambda }}.`$
Aknowledgements.
We are grateful to Vincent Moncrief for interesting discussions about this paper.
We thank the Kavli Institute for Theoretical Physics at Santa Barbara, LโInstitut des Hautes Etudes Scientifiques at Bures-sur-Yvette and the department of mathematics of the University of Washington for providing very pleasant and stimulating environments for our collaboration on this work. This work was partially supported by the NSF, under grants PHY-0099373 and PHY-0354659 at Oregon.
References.
\[BM\] B. K. Berger and V. Moncrief โNumerical evidence for an oscillatory singularity in generic $`U(1)`$ symmetric cosmologies on $`T^3\times R\mathrm{"}`$ Phys. Rev. D 58 064023-1-8 (1998).
\[M86\] V. Moncrief Reduction of Einstein equations for vacuum spacetimes with U(1) spacelike isometry group, Annals of Physics 167 (1986), 118-142
\[CB-M 96\] Y. Choquet-Bruhat and V. Moncrief Existence theorem for solutions of Einstein equations with 1 parameter spacelike isometry group, Proc. Symposia in Pure Math, 59, 1996, H. Brezis and I.E. Segal ed. 67-80.
\[CBIM\] Y. Choquet-Bruhat, J. Isenberg and V.Moncrief โTopologically general $`U(1)`$ symmetric Einsteinian spacetimes with AVTD behaviourโ Il Nuovo Cimento B, Vol. 119, issue no. 7-9, 2005.
\[IM\] J. Isenberg and V. Moncrief, โAsymptotic behavior in polarized and half-polarized $`U(1)`$ symmetric spacetimesโ, Class. Qtm. Grav. 19, 5361-5386 (2002).
\[KR\] S. Kichenassamy and A.D. Rendall, โAnalytical description of singularities in Gowdy spacetimesโ, Class. Qtm.Grav.15, 1339-1355 (1998).
\[AR\] L. Andersson and A. Rendall, โQuiescent cosmological singularitiesโ, Comm. Math. Phys. 218, 479-511 (2001).
\[DHRW\] T. Damour, M Henneaux, A. Rendall, and M. Weaver, โKasner-like behavior for subcritical Einstein-matter systemsโ, Ann. H. Poin. 3, 1049-1111 (2002).
\[IM92\] J. Isenberg and V, Moncrief, โAsymptotic behavior of the gravitational field and the nature of singularities in Gowdy spacetimesโ, Ann. Phys. 99, 84-122 (1992).
\[CBM\] Y. Choquet-Bruhat and V.Moncrief, โFuture complete $`U(1)`$ symmetric einsteinian spacetimesโ, Ann. Henri Poincare, 2, 1007-1064 (2001) See also โNon linear stability of einsteinian spacetimes with $`U(1)`$ isometry groupโ, gr-qc/0302021.
\[CB\] Y. Choquet-Bruhat, โFuture complete $`U(1)`$ symmetric einsteinian spacetimes, the unpolarized caseโ, in โ50 Years of the Cauchy Problemโ, eds. P. Chrusciel and H. Friedrich (2004).
\[K\] S. Kichenassamy, โNonlinear Wave Equationsโ (Dekker, NY) (1996).
\*Acadรฉmie des Sciences, 23 quai Conti 755270 Paris cedex 06, France
\**Department of Mathematics and Institute of Theoretical Science, University of Oregon, Eugene, OR 97403-5203, USA |
warning/0506/cond-mat0506746.html | ar5iv | text | # Universality Class of One-Dimensional Directed Sandpile Models
## Abstract
A general $`n`$-state directed โsandpileโ model is introduced. The stationary properties of the $`n`$-state model are derived for $`n<\mathrm{}`$, and analytical arguments based on a central limit theorem show that the model belongs to the universality class of the totally asymmetric Oslo model, with a crossover to uncorrelated branching process behavior for small system sizes. Hence, the central limit theorem allows us to identify the existence of a large universality class of one-dimensional directed sandpile models.
Sandpile models have attracted much analytical attention in recent years, largely due to their application to the development of self-organized criticality (SOC)Bak ; Dhar1989 ; Dhar ; Priezzhev2001 . Analytical solutions, scaling arguments and numerics have shown that many models share the same critical exponents and scaling functions, leading to the notion of universality classes for such systems, as in equilibrium systems Dhar1989 ; Paczuski1996 ; Mohanty2002 . In the following, we map an $`n`$-state directed model to a random walker problem and, using a central limit theorem for dependent random variables Brown1971 , derive the conditions for scaling and the associated critical exponents. The exponents are in exact agreement with those derived for the Totally Asymmetric Oslo Model (TAOM) Pruessner2003b ; Pruessner2004 , which is a special case of this general model. For small system sizes, we find that the model may exhibit different scaling, which corresponds to an uncorrelated branching process, with a crossover characterized by an $`n`$ dependent crossover length, $`\xi _n`$.
The model considered is an $`n`$-state directed โsandpileโ model. The system exists on a one-dimensional lattice with $`L`$ sites. Each site, $`i`$, is in one of $`n`$ states, $`z_i[0,n1]`$, which represents the number of particles on site $`i`$.
At the beginning of each time step, we add one particle to site $`i=1`$. This site may then topple a number of times, each toppling redistributing a particle to the next site, $`z_1z_11`$ and $`z_2z_2+1`$. When site $`i=2`$ receives a particle it may also undergo topplings, redistributing particles to site $`i=3`$, and so on, with particles being passed to sites of increasing $`i`$. Note that when site $`i=L`$ topples, the redistributed particle will leave the system. When all activity ceases, a new time step commences. The avalanche size, $`s`$, is defined as the total number of topplings during a single time step. The only restrictions on the toppling rules are: (i) A toppling may never cause $`z_i`$ to become negative. (ii) If $`z_i>n1`$, then site $`i`$ must topple. (iii) Each time site $`i`$ topples, it redistributes one particle to site $`i+1`$ only. (iv) Particles are conserved in the bulk, only leaving the system at the boundary site $`i=L`$ and entering when particles are added to the boundary site $`i=1`$ at the beginning of each time step. (v) The number of topplings a site undergoes is non-deterministic for at least one value of $`z`$. This final restriction disallows deterministic toppling rules, which lead to trivial dynamics.
For the following general discussion, there is no need to specify the toppling rules in any more detail. Later, when presenting numerical results, we will consider a specific implementation of a general class of probabilistic toppling rules that satisfies the โrestrictionsโ (i)-(v) above.
The quantity of interest in a sandpile which has reached a non-equilibrium steady state is the avalanche-size probability, $`P(s;L,1)`$, which is the probability of observing an avalanche of size $`s`$ in a system of size $`L`$ when one particle is added to site $`i=1`$. SOC is associated with a time-independent avalanche-size probability which obeys simple finite-size scaling
$$P(s;L,1)=as^\tau ๐ข\left(s/bL^D\right)\text{for }s,L1\text{,}$$
(1)
where $`a`$ and $`b`$ are non-universal constants, $`\tau `$ and $`D`$ are universal critical exponents, and $`๐ข`$ is a universal scaling function. The $`k`$th moment of the avalanche-size probability is
$`Q_{L,1}^{(k)}`$ $`={\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}s^kP(s;L,1)`$
$`{\displaystyle _1^{\mathrm{}}}๐sas^{k\tau }๐ข(s/bL^D)`$
$`=ab^{1+k\tau }G_k(L)L^{D(1+k\tau )},`$ (2)
where the sum has been approximated by an integral and $`G_k(L)=_{1/bL^D}^{\mathrm{}}๐uu^{k\tau }๐ข(u)`$. Provided that $`0<G_k(\mathrm{})<\mathrm{}`$, we have
$$Q_{L,1}^{(k)}=\mathrm{\Gamma }_kL^{\gamma _k}\text{for }L1\text{,}$$
(3)
where $`\gamma _k=D(1+k\tau )`$ is a universal exponent and $`\mathrm{\Gamma }_k=ab^{1+k\tau }G_k(L)`$ is a non-universal amplitude which is a constant for $`L1`$. Hence, the scaling of the moments with system size $`L`$ is a universal feature which is independent of particular details of the dynamics if Eq. (1) is valid, the approximation to the integral in Eq. (2) does not affect the scaling behavior of $`Q_{L,1}^{(k)}`$, and $`G_k(L)`$ approaches a non-zero constant for $`L\mathrm{}`$. In the following we shall show that under a precise set of conditions, Eq. (3) will hold with $`D=3/2`$ and $`\tau =4/3`$.
Using a simple extension to the Markov matrix methods used in Ref. Pruessner2004 , it can be shown that if the Markov matrix representing the evolution operator for the model is regular and $`n<\mathrm{}`$, then there is a unique stationary state. Note that we will only consider toppling rules for which the evolution operator is regular and so a unique stationary state exists. Following a similar calculation as in Pruessner2004 , we find that in this state, the number of particles, $`z_i`$, in each site, $`i`$, is an independent identically distributed random variable with probability $`p_z`$ tobepublished . Hence, we find that the probability of occurrence of a configuration $`\{z_i\}=\{z_1,z_2,\mathrm{},z_L\}`$ is
$$p_{\{z_i\}}=\underset{i=1}{\overset{L}{}}p_{z_i},$$
(4)
where the values of $`p_z`$ depend on the details of the toppling rules. This is known as a product state and has the property that there are no spatial correlations. However, we shall shortly argue that there exists a crossover length, $`\xi _n`$, such that for system sizes $`L\xi _n`$, there are temporal correlations which produce non-trivial behavior.
We define
$$Q_{L,m}^{(k)}\underset{s=0}{\overset{\mathrm{}}{}}s^kP(s;L,m)$$
(5)
as the $`k`$th moment of the avalanche-size probability, $`P(s;L,m)`$, for a system of size $`L`$ which has received $`m`$ particles at site $`i=1`$.
The first moment is easily derived from the fact that, in the stationary state, the average number of particles which leave the system through the open boundary must equal the number of particles added to the system. Each of the $`m`$ particles topples exactly $`L`$ times, and
$$Q_{L,m}^{(1)}=mL,$$
(6)
implying $`\gamma _1D(2\tau )=1`$.
To derive the scaling of higher moments, we introduce $`P(t,s;1,L,m)`$ as the joint probability that a system of size $`1+L`$ which has received $`m`$ particles at site $`i=1`$ undergoes $`t`$ topplings in the first site and $`s`$ in the remaining $`L`$ sites. We note that since the model is directed and the stationary state is a product state,
$$P(t,s;1,L,m)=P(t;1,m)P(s;L,t).$$
(7)
At the beginning of each time step we add one particle to site $`i=1`$ and it will topple $`s_1`$ times with probability $`P(s_1;1,1)`$. The second site therefore receives $`s_1`$ particles and as a result topples $`s_2`$ times with probability $`P(s_2;1,s_1)`$. The probability of site 2 toppling $`s_2`$ times, denoted $`\varphi _2(s_2)`$, is
$$\varphi _2(s_2)=\underset{s_1=1}{\overset{\mathrm{}}{}}P(s_2;1,s_1)P(s_1;1,1)$$
(8)
which follows from Eq. (7). If we define $`\varphi _i(x)`$ as the probability that site $`i`$ topples $`x`$ times, with $`\varphi _1(s_1)P(s_1;1,1)`$, then
$$\varphi _{i+1}(x)=\underset{y=1}{\overset{\mathrm{}}{}}P(x;1,y)\varphi _i(y).$$
(9)
This describes a discrete random walker on the interval $`[0,\mathrm{}]`$. Since activity stops when one of the sites topples zero times, there is an absorbing boundary at $`x=0`$. The probability of hopping from $`y`$ to $`x`$ in a single step is given by $`P(x;1,y)`$. If we denote a particular trajectory of a walker $`x(i)`$, $`i=0\mathrm{}L`$, then the corresponding avalanche size is
$$s=\underset{i=1}{\overset{L}{}}x(i)$$
(10)
with $`x(0)=1`$. This corresponds to the area under the first $`L`$ steps of a random walk with an absorbing boundary at $`x=0`$.
Of course, this is a correlated random walker and individual steps are not independent because the probability of hopping a certain distance varies depending on where the random walker is according to $`P(s;1,m)`$.
We now define
$$\stackrel{~}{Q}_{1,m}^{(2)}=\underset{s=0}{\overset{\mathrm{}}{}}(sm)^2P(s;1,m),$$
(11)
which is the width of the probability $`P(s;1,m)`$ around the mean value, $`m`$. Using a martingale theorem Brown1971 , we can show that iff there exists a number $`0<M<\mathrm{}`$ such that $`0<\stackrel{~}{Q}_{1,m}^{(2)}<M`$ for all $`m`$, then in the limit $`L\mathrm{}`$, the distribution $`\varphi _i(x)`$ will converge for large $`i`$ to that for an equivalent independent random walker tobepublished . Hence, we find that all moments scale,
$$Q_{L,1}^{(k)}L^{(3k1)/2}=L^{\frac{3}{2}(1+k\frac{4}{3})}$$
(12)
and we can read off the exponents $`D=3/2`$ and $`\tau =4/3`$.
Hence, we must find the conditions under which $`\stackrel{~}{Q}_{1,m}^{(2)}`$ is non-zero and finite. Note, again, that we are assuming the existence of a unique stationary state. Consider a site with $`z`$ particles which has received $`m`$ particles. After $`s`$ topplings have taken place it will contain $`z^{}=z+ms`$ particles. Since both $`z`$ and $`z^{}`$ must lie between $`0`$ and $`n1`$, $`P(s;1,m)`$ may only be non-zero for $`mn+1sm+n1`$. Hence $`\stackrel{~}{Q}_{1,m}^{(2)}(n1)^2`$, which is finite for $`n<\mathrm{}`$. In order to have $`\stackrel{~}{Q}_{1,m}^{(2)}=0`$, there must be an $`m`$ for which only an avalanche of size $`s=m`$ is allowed. However, this is only possible if the number of topplings a site undergoes on receiving a particle is fully deterministic, which are trivial dynamics. Hence, if $`n<\mathrm{}`$ and the toppling rule leads to non-trivial dynamics we have $`0<\stackrel{~}{Q}_{1,m}^{(2)}(n1)^2`$ and Eq. (12) follows.
The scaling of the TAOM will only be observed asymptotically for $`L\mathrm{}`$. However, we hypothesize the existence of an $`n`$ dependent crossover length, $`\xi _n`$, such that TAOM scaling is observed for $`L\xi _n`$, with different behavior for $`L\xi _n`$. To see what happens for system sizes $`L\xi _n`$, consider a system when a particle is added to site $`i=1`$. If the probability, $`p_z`$, that a site is occupied by $`z`$ particles has support for all $`z[0,n1]`$, then, for $`n1`$, it is likely that $`0zn1`$. If the number of particles on a site $`z`$ is neither close to $`0`$ nor $`n1`$, the propagating avalanches will not be sensitive to the medium within which it is propagating and the system will be temporally uncorrelated. However, as the avalanche propagates through the system, fluctuations in the number of topplings increase and as each subsequent site is less likely to have uncorrelated topplings they will start to feel the fact that $`n`$ is finite. Hence, for small system sizes, $`1L\xi _n`$, the avalanches are uncorrelated and will correspond to an uncorrelated branching process with exponents $`D=2`$ and $`\tau =3/2`$ Harris1963 . For large system sizes, $`L\xi _n`$, temporal correlations will emerge and the system will crossover to behavior of the TAOM with $`D=3/2`$ and $`\tau =4/3`$.
We support the above arguments with numerical data from the following โtypicalโ realization: A site $`i`$, $`0<z_i<n1`$, which receives a particle will topple once with probability $`1/2`$ and twice with probability $`1/4`$. If $`z_i=0`$, then it cannot topple twice and will topple once with probability $`1/2`$. If $`z_i=n1`$, it will topple once or twice, each with probability $`1/2`$. It can be shown that the support of $`p_z`$ extends over all possible states $`z[0,n1]`$. We shall compare numerical results from this model with exact results from the corresponding uncorrelated branching process, which we can calculate analytically Harris1963 .
Figure 2(a) displays measurements of the rescaled second moment, $`Q_{L,1}^{(2)}/L^{5/2}`$ vs. $`L`$. From the arguments above, we expect $`Q_{L,1}^{(2)}/L^{5/2}`$ to scale like $`L^{1/2}`$ for $`1L\xi _n`$ and approach a constant for $`L\xi _n`$, which is supported by the numerics. We do not attempt a data collapse because, although it is clear that the crossover length $`\xi _n`$ is a non-decreasing function of $`n`$, it has a non-universal functional dependence.
We also consider the moment ratios
$$g_k\frac{s^ks^{k2}}{s^2^{k1}}\frac{\mathrm{\Gamma }_k\mathrm{\Gamma }_1^{k2}}{\mathrm{\Gamma }_2^{k1}}.$$
(13)
For an avalanche-size probability of the form Eq. (1), $`g_k`$ will be universal constants, that is, they only depend on $`\tau `$ and $`๐ข`$. Figure 2(b) displays $`g_3`$ vs. system size, $`L`$. Since $`G_k(L)`$ is only constant for $`L1`$, the measured $`g_3`$ will only converge toward the universal constant for large $`L`$. These have values $`g_3=9/5`$ for the branching process and we measure $`g_31.29`$ for the TAOM. For $`1L\xi _n`$, the measured values follow the exact result for the uncorrelated branching process, with a crossover to the TAOM curve for $`L\xi _n`$.
We have shown that a general $`n`$-state directed sandpile model of self-organized criticality belongs to the same universality class as the totally asymmetric Oslo model, recently solved in Ref. Pruessner2004 . The precise conditions for this universality are that the evolution operator is regular, the avalanches are non-deterministic, and that $`n`$ is finite. We have argued that there is an $`n`$ dependent crossover length $`\xi _n`$, which separates uncorrelated branching process exponents, $`D=2`$ and $`\tau =3/2`$, for $`1L\xi _n`$ from TAOM exponents, $`D=3/2`$ and $`\tau =4/3`$ for $`L\xi _n`$. This crossover may be considered a consequence of temporal correlations emerging in the system, which moves it away from the uncorrelated branching process, associated with mean-field exponents. The conditions for a system to be in this universality class have been found using a central limit theorem for dependent random variables Brown1971 . This is the first time a technique of applying a central limit theorem to the discrete model has been used to explicitly and precisely identify a universality class of non-equilibrium self-organized critical systems and we expect much new research will follow along these lines.
The authors wish to thank G. Pruessner and D. Dhar for helpful comments on the manuscript, and B. Derrida for help in finding the area under a random walker. M.S. gratefully acknowledges the financial support of U.K. EPSRC through grant GR/P01625/01. |
warning/0506/hep-ph0506305.html | ar5iv | text | # The Interval Approach to Braneworld Gravity
## 1 Introduction
Models of braneworld gravity are so varied and popular that they are commonly denoted in a cryptic shorthand: ADD , RSI , RSII , LR , AED , UED , DGP , KR , etc. Most models are based upon an orbifold as the background geometry, usually $`๐^\mathrm{๐}/๐_\mathrm{๐}`$. The analysis of such models begins with making explicit orbifold projections on the equations of motion, and integrating the equations of motion through the orbifold fixed points to obtain junction conditions -. This is a convenient shortcut which produces correct results in simple analyses of simple systems.
However there are problems with this standard approach. The first is that general relativity (GR) is defined on manifolds, not orbifolds. An orbifold is a well-defined but singular limit of a smooth manifold. It is possible to treat a $`๐_\mathrm{๐}`$ orbifold as a limit of a manifold with boundaries, however a generally covariant action principle for manifolds with boundaries is usually only discussed in the case where metric fluctuations are restricted to vanish on the boundaries . In braneworld models we are specifically interested in the case where the metric fluctuations do not vanish on the boundaries. This problem is usually finessed by applying junction conditions , which does not address the status of general covariance in such a system, or resolve whether one can define an unambiguous action principle.
Another complication is that braneworld models often contain extra scalar degrees of freedom, coming from fluctuations of the higher dimensional metric. In many setups - it appears that these extra scalars are kinetic ghosts (i.e. they have kinetic terms with the wrong sign) or have a kinetic term with vanishing coefficient, leading to strong coupling behavior. To exhibit either kind of pathology explicitly requires computing the full gauge-fixed effective action of (at least) the linearized theory. This certainly requires a well-defined action principle as a starting point, and it requires an unambiguous understanding of the full general coordinate invariance of the model.
In this paper we provide a general set of definitions and methodologies for analyzing models of braneworld gravity. We begin with a number of familiar examples. In ยง2.1 we introduce basic concepts and notation of the interval picture using a simple 5d scalar field theory. In ยง2.2 we discuss 5d abelian gauge theory in a fixed braneworld background. In ยง3.1 and ยง3.2, we treat 5d gravity in a flat $`๐^\mathrm{๐}/๐_\mathrm{๐}`$ background. We contrast the usual orbifold techniques with the interval picture, where we never invoke $`๐_\mathrm{๐}`$ projections or junction conditions. To simplify the presentation we employ limits of the general results derived in ยง5. In ยง4 we do a similar analysis for the original Randall-Sundrum model, in the interval picture.
We then proceed to analyze a general $`AdS_5/AdS_4`$ setup with two branes, including brane kinetic terms for gravity. Here it is already not obvious from previous work how to count the physical scalars coming from the metric. Orbifold projections by their very nature are only implemented in a coordinate system where the branes are โstraightโ, i.e., located at fixed slices of the 5d coordinate $`y`$. We call such coordinate systems โstraight gaugesโ, and define them precisely. For setups with more than one brane, none of the standard gauge choices of gravity (axial, harmonic, de Donder, Gaussian normal) are straight gauges in a single coordinate patch. A coordinate transformation of the metric that violates the straight gauge condition has the appearance of a scalar metric perturbation, called a โbrane-bendingโ mode. In the orbifold approach one cannot distinguish between the following two possibilities:
* orbifold gravity does not respect the full general coordinate invariance of gravity on manifolds, and thus some brane-bending modes are physical;
* orbifold gravity does implement the full general coordinate invariance, and brane-bending modes are always pure gauge.
The interval picture shows that the second alternative is the correct one. We show that warped two-brane setups have at most a single 4d scalar mode (a radion) coming from the metric.
Other authors have already introduced some of the concepts employed in this paper -.
## 2 Orbifolds in field theory
### 2.1 Scalars on a 5d orbifold
Most of the literature on warped extra dimensions is based upon the idea of field theory on the simple orbifold $`๐^\mathrm{๐}/๐_\mathrm{๐}`$. For a nongravitational theory there is a simple unambiguous implementation of the orbifolding. Consider for example a real 5d scalar field $`\varphi (x^\mu ,y)`$. Compactifying the $`y`$ direction on a circle with radius $`L/\pi `$ implies that $`\varphi `$ should be decomposed into the appropriate Fourier modes:
$`\varphi (x,y)={\displaystyle \frac{a_0(x)}{\sqrt{2}}}+{\displaystyle \underset{n>0}{}}\left[a_n(x)\mathrm{cos}\left({\displaystyle \frac{\pi ny}{L}}\right)+b_n(x)\mathrm{sin}\left({\displaystyle \frac{\pi ny}{L}}\right)\right].`$ (1)
The $`๐_\mathrm{๐}`$ orbifolding around the point $`y=0`$ then amounts to projecting out all of the odd modes, i.e., setting all the 4d fields $`b_n(x)`$ to zero. It is also possible to define a different orbifolded theory in which all of the even modes are projected out. Note that the modes which are even around $`y=0`$ are also even around $`y=L`$. These are the two fixed points of the orbifold. By periodicity, the fixed point $`y=L`$ is identified with the point $`y=L`$. The orbifold can be regarded as extending from $`L`$ to $`L`$, with two fixed points but no boundaries.
This definition of a field theory orbifold is certainly not adequate for a theory which includes gravity. In particular the fixed points of an orbifold lead to ambiguities in the formulation of GR. This is especially true if one introduces delta function sources at the fixed points of the orbifold. It is not obvious in this case that there is a well-defined action principle, and the status of general coordinate invariance is murky.
To examine these issues, we first need a definition of the field theory orbifold at the level of the action, rather than as a projection on the equations of motion. We use essentially the same definition as . Consider again a real 5d scalar field. Including polynomial sources located at the fixed points, the $`๐^\mathrm{๐}/๐_\mathrm{๐}`$ field theory orbifold is defined by the following action (our metric signature is $`++++`$):
$`S`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle }d^4x({\displaystyle _\epsilon ^{L\epsilon }}dy+{\displaystyle _{L+\epsilon }^\epsilon }dy)\{{\displaystyle \frac{1}{2}}^M\varphi _M\varphi +V(\varphi )`$ (2)
$`+(\delta (y\epsilon )+\delta (y+\epsilon ))V_0(\varphi )+(\delta (yL+\epsilon )+\delta (y+L\epsilon ))V_L(\varphi )\}.`$
The action comes from integrating over two intervals: $`[L+\epsilon ,\epsilon ]`$ and $`[\epsilon ,L\epsilon ]`$. It is understood that we are imposing periodicity under $`yy+2L`$, as before. Thus the $`๐^\mathrm{๐}/๐_\mathrm{๐}`$ orbifold, which has two fixed points and no boundary, is here represented as a limit of a theory with two intervals and four boundary points. In this simple example there is a bulk potential $`V(\varphi )`$ and two โbraneโ sources $`V_0(\varphi )`$ and $`V_L(\varphi )`$. These brane sources have support only at the four boundary points; they are introduced symmetrically to reproduce the usual delta function brane sources in the limit, e.g.:
$`\underset{\epsilon 0}{lim}\left({\displaystyle _\epsilon ^{L\epsilon }}๐y+{\displaystyle _{L+\epsilon }^\epsilon }๐y\right)\left(\delta (y\epsilon )+\delta (y+\epsilon )\right)V_0(\varphi )`$
$`=\underset{\epsilon 0}{lim}\left({\displaystyle _{L+\epsilon }^{L\epsilon }}๐y\right){\displaystyle \frac{1}{2}}\left(\delta (y\epsilon )+\delta (y+\epsilon )\right)V_0(\varphi )={\displaystyle _L^L}๐y\delta (y)V_0(\varphi ).`$ (3)
It is important to note that we are assuming that the brane sources are continuous, e.g.
$`\underset{\epsilon 0}{lim}V_0(\varphi )|_\epsilon =\underset{\epsilon 0}{lim}V_0(\varphi )|_\epsilon V_0(\varphi )|_0.`$ (4)
In order to have a well-defined action principle, a field theory with boundaries requires the imposition of appropriate boundary conditions. To see how this works for our simple example, consider the full variation of the action, keeping surface terms:
$`\delta S`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle }d^4x\{({\displaystyle _\epsilon ^{L\epsilon }}dy+{\displaystyle _{L+\epsilon }^\epsilon }dy)(^M_M\varphi +{\displaystyle \frac{\delta V(\varphi )}{\delta \varphi }})\delta \varphi `$ (5)
$`+\left[\left(\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_L}{\delta \varphi }}\right)\delta \varphi \right]_{L\epsilon }+\left[\left(\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_0}{\delta \varphi }}\right)\delta \varphi \right]_\epsilon `$
$`+\left[(\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_0}{\delta \varphi }})\delta \varphi \right]_\epsilon +\left[(\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_L}{\delta \varphi }})\delta \varphi \right]_{L+\epsilon }\},`$
where prime denotes a derivative with respect to $`y`$.
The bulk equation of motion is
$`^M_M\varphi +{\displaystyle \frac{\delta V(\varphi )}{\delta \varphi }}=0.`$ (6)
To make the action stationary, this must be supplemented by boundary conditions at the four boundary points. One option is to impose Dirichlet boundary conditions, i.e., to require that $`\delta \varphi (x^\mu ,y)`$ vanishes at the boundaries. This is not usually what one wants for brane models, although it is the assumption used for general relativity with boundaries.
The other option is to supplement the bulk equations of motion by four โbrane-boundaryโ equations:
$`\left[\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_L}{\delta \varphi }}\right]_{L\epsilon }`$ $`=0;`$
$`\left[\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_0}{\delta \varphi }}\right]_\epsilon `$ $`=0;`$
$`\left[\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_0}{\delta \varphi }}\right]_\epsilon `$ $`=0;`$
$`\left[\varphi ^{}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta V_L}{\delta \varphi }}\right]_{L+\epsilon }`$ $`=0.`$
In the limit $`\epsilon 0`$ this is equivalent to:
$`\varphi ^{}|_{0^+}=\varphi ^{}|_0^{};`$
$`\varphi ^{}|_L^{}=\varphi ^{}|_{L^+};`$
$`2\varphi ^{}|_{0^+}+{\displaystyle \frac{\delta V_0}{\delta \varphi }}|_0=0;`$ (8)
$`2\varphi ^{}|_L^{}+{\displaystyle \frac{\delta V_L}{\delta \varphi }}|_L=0.`$
The brane-boundary conditions (2.1) are invariant under interchanging the two intervals combined with $`yy`$. From this it is clear that we can always restrict our attention to solving for $`\varphi `$ in the interval $`0<y<L`$, imposing the second two boundary equations of (2.1). The solution in the interval $`L<y<0`$ then follows by applying the first two relations of (2.1). In this paper we will always be content to display our solutions on $`0<y<L`$.
Note that if we remove the brane sources, we get simple Neumann boundary conditions at the boundaries. Then in the limit $`\epsilon 0`$, the brane-boundary equations are precisely equivalent to the usual orbifold projection.
Following earlier work , we will use the name โinterval pictureโ to refer to this approach to defining field theory orbifolds at the level of the action.
### 2.2 Abelian gauge theory on a warped orbifold
A more ambitious example is to consider a 5d abelian gauge theory, with brane kinetic terms, in a warped orbifold background with two branes. This setup was analyzed in the conventional orbifold picture in . The interval picture action is given by:
$`S`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle }d^4x({\displaystyle _\epsilon ^{L\epsilon }}dy+{\displaystyle _{L+\epsilon }^\epsilon }dy)\sqrt{G}{\displaystyle \frac{1}{8g_5^2}}\{G^{MP}G^{NQ}F_{MN}F_{PQ}`$
$`+2r_U\{\delta (y\epsilon )+\delta (y+\epsilon )\}G^{\mu \rho }G^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma }`$
$`+2r_I\{\delta (yL+\epsilon )+\delta (y+L\epsilon )\}G^{\mu \rho }G^{\nu \sigma }F_{\mu \nu }F_{\rho \sigma }\}.`$
Here we have introduced a fixed background metric which is warped: $`G_{\mu \nu }=a^2(y)\eta _{\mu \nu }`$, $`G_{44}=1`$, $`G_{\mu 4}=0`$. The function $`a(y)`$ is the warp factor. For e.g. an $`AdS_5`$ background as in Randall-Sundrum (RS), we would have:
$`a(y)=\{\begin{array}{cc}\mathrm{e}^{ky}\hfill & L<y<0,\hfill \\ \mathrm{e}^{ky}\hfill & 0<y<L,\hfill \end{array}`$ (12)
where $`k`$ is the inverse $`AdS_5`$ radius of curvature.
To be clear, letโs pull out all the warp factors explicitly, and for the rest of this example we raise and lower indices with $`\eta _{\mu \nu }`$. Then:
$`F^{MN}F_{MN}`$ $`=`$ $`{\displaystyle \frac{2}{a^4}}\left(_\mu A_\nu ^\mu A^\nu _\mu A_\nu ^\nu A^\mu \right)`$ (13)
$`+{\displaystyle \frac{2}{a^2}}\left(_\mu A_4^\mu A_42^\mu A_4A_\mu ^{}+A_{}^{\mu }{}_{}{}^{}A_\mu ^{}\right),`$
where prime denotes a derivative with respect to $`y`$. The bulk equations of motion (EOM) are:
$`^2A_\mu _\mu ^\nu A_\nu \left(a^2_\mu A_4\right)^{}+\left(a^2A_\mu ^{}\right)^{}=0;`$ (14)
$`a^2^2A_4a^2^\mu A_\mu ^{}=0.`$ (15)
The two brane-boundary equations are:
$`a^2_\mu A_4|_{0^+}+a^2A_\mu ^{}|_{0^+}+r_U\left[^2A_\mu _\mu ^\nu A_\nu \right]_{y=0}`$ $`=`$ $`0;`$
$`a^2_\mu A_4|_L^{}a^2A_\mu ^{}|_L^{}+r_I\left[^2A_\mu _\mu ^\nu A_\nu \right]_{y=L}`$ $`=`$ $`0.`$ (16)
The 5d abelian gauge transformations are generated by $`\mathrm{\Lambda }(x,y)`$:
$`A_\mu `$ $``$ $`A_\mu +_\mu \mathrm{\Lambda };`$
$`A_4`$ $``$ $`A_4+\mathrm{\Lambda }^{}.`$ (17)
We want to determine the physical degrees of freedom of this theory in the interval picture. To begin, we do a partial gauge-fixing by choosing
$`\mathrm{\Lambda }^{(\mathrm{I})}(x,y)={\displaystyle ^y}A_4๐y+{\displaystyle ^y}F(y)\psi (x)๐y.`$ (18)
With this partial gauge-fixing we have
$`A_4(x,y)=F(y)\psi (x).`$ (19)
The function $`F(y)`$ is fixed but arbitrary; different choices of $`F(y)`$ correspond to different gauges. The 4d field $`\psi (x)`$, on the other hand, appears at this point to be a 4d scalar degree of freedom.
The bulk equation (15) becomes:
$`F^2\psi =^\mu A_\mu ^{}.`$ (20)
The 5d field $`A_M(x,y)`$ has a different 4d tensor decomposition depending upon whether or not it is a zero mode of the operator $`^2`$, i.e., whether it is a massless mode in the 4d sense. Thus we need to solve separately for the massless and massive modes.
#### 2.2.1 massless modes
When $`A_M(x,y)`$ is a zero mode of $`^2`$, we can write:
$`A_\mu (x,y)=A_\mu ^T(x,y)+_\mu \varphi (x,y)+A_\mu ^L(x,y),`$ (21)
where $`\varphi (x,y)`$ is pure gauge, $`A_\mu ^L(x,y)`$ is the 4d longitudinal mode, and $`A_\mu ^T(x,y)`$ are the two remaining transverse modes which are not pure gauge. In addition, we are only looking at the part of $`\psi (x)`$ which satisfies $`^\mu _\mu \psi =0`$.
The bulk equation (20) reduces to:
$`^\mu A_\mu ^L{}_{}{}^{}(x,y)=0^\mu A_\mu ^L(x,y)=\rho (x),`$ (22)
where $`\rho (x)`$ is an arbitrary function. Defining
$`\chi _\mu ^{(0)}{}_{}{}^{}(x,y)=A_\mu ^{}F(y)_\mu \psi (x),`$ (23)
the remaining bulk equation (14) gives:
$`(a^2\chi _\mu ^{(0)}{}_{}{}^{})^{}=_\mu \rho ,`$ (24)
while the brane-boundary equations become:
$`\left[a^2\chi _\mu ^{(0)}{}_{}{}^{}r_U_\mu \rho \right]_0`$ $`=`$ $`0;`$
$`\left[a^2\chi _\mu ^{(0)}{}_{}{}^{}+r_I_\mu \rho \right]_L`$ $`=`$ $`0,`$ (25)
where we are employing a shorthand notation $`(0,L)`$ to distinguish the two independent brane-boundary conditions.
Provided that $`r_U+r_I+L0`$, the only simultaneous solution of (24-25) is
$`\chi _\mu ^{(0)}{}_{}{}^{}=0;A_\mu ^L=0.`$ (26)
This in turn implies:
$`A_\mu ^T(x,y)`$ $`=`$ $`A_\mu ^T(x);`$
$`\varphi (x,y)`$ $`=`$ $`(y)\psi (x)+\varphi (x),`$ (27)
where $`(y)`$ is defined by
$`^{}(y)=F(y),`$ (28)
with the integration constant set to zero.
To count massless degrees of freedom, we perform the gauge transformation defined by
$`\mathrm{\Lambda }^{(\mathrm{II})}(x,y)=(y)\psi (x)\varphi (x).`$ (29)
This takes us to an axial gauge, which is also the unitary gauge for this model:
$`A_\mu (x,y)`$ $`=`$ $`A_\mu ^T(x);`$
$`A_4(x,y)`$ $`=`$ $`0.`$ (30)
There are no extra massless scalar modes, as expected.
#### 2.2.2 massive modes
In this case we decompose
$`A_\mu =A_\mu ^T+_\mu \varphi ,`$ (31)
where $`A_\mu ^T`$ is transverse. Then (20) becomes:
$`F^2\psi =^2\varphi ^{}.`$ (32)
However we already gauge-fixed $`\psi (x)`$ (massless and massive modes) to zero by the transformation (29). Since also we are looking only at massive modes of $`\varphi `$ we can remove the $`^2`$ and conclude:
$`\varphi (x,y)=\varphi (x).`$ (33)
Note $`\varphi (x)`$ just represents the residual 4d gauge freedom that preserves the axial gauge. Thus we can gauge-fix it to zero.
So far we have:
$`A_\mu (x,y)`$ $`=`$ $`A_\mu ^T(x,y);`$
$`A_4(x,y)`$ $`=`$ $`0.`$ (34)
The massive KK modes have three physical polarizations (and no residual gauge freedom), as appropriate for a massive vector.
Plug this into the bulk equation (14):
$`^2A_\mu ^T+(a^2A_\mu ^T{}_{}{}^{})^{}=0.`$ (35)
The brane-boundary equations become:
$`\left[a^2A_\mu ^T{}_{}{}^{}+r_U^2A_\mu ^T\right]_0=0;`$ (36)
$`[a^2A_\mu ^T{}_{}{}^{}+r_I^2A_\mu ^T]_L=0.`$ (37)
We introduce a Kaluza-Klein (KK) decomposition for the massive transverse modes $`A_\mu ^T(x,y)`$:
$`A_\mu ^T(x,y)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}A_\mu ^{(n)}(x)\chi ^{(n)}(y).`$ (38)
We can take the $`A_\mu ^{(n)}(x)`$ to be on-shell in the 4d sense, so $`^2`$ $``$ $`p^2`$ $``$ $`m_n^2`$. The bulk equation of motion becomes:
$`(a^2\chi ^{(n)}{}_{}{}^{}(y))^{}=m_n^2\chi ^{(n)}(y).`$ (39)
We can turn the above into Besselโs equation by making the substitutions:
$`\chi ^{(n)}={\displaystyle \frac{1}{a(y)}}f^{(n)};z_n={\displaystyle \frac{m_n}{ka(y)}},`$ (40)
where now we are going to restrict to the RS case, so $`a^{}/a=k`$, $`a^{\prime \prime }/a=k^2`$. This produces
$`\left(z_n^2{\displaystyle \frac{d^2}{dz_n^2}}+z_n{\displaystyle \frac{d}{dz_n}}+(z_n^21)\right)f^{(n)}=0.`$ (41)
The solutions are:
$`\chi ^{(n)}={\displaystyle \frac{1}{a}}N_n\left(J_1(z_n)+bY_1(z_n)\right),`$ (42)
where $`b`$ is a constant and the $`N_n`$ are normalization constants.
The brane-boundary equations now become:
$`\chi ^{(n)}{}_{}{}^{}|_{0}^{}={\displaystyle \frac{r_Um_n^2}{a^2}}\chi ^{(n)}|_0;\chi ^{(n)}{}_{}{}^{}|_{L}^{}={\displaystyle \frac{r_Im_n^2}{a^2}}\chi ^{(n)}|_L.`$ (43)
Using
$`\chi ^{(n)}{}_{}{}^{}={\displaystyle \frac{m_n}{a}}N_n(J_0(z_n)+bY_0(z_n)),`$ (44)
we determine the constant $`b`$ and the eigenvalues $`m_n`$:
$`b=\left[{\displaystyle \frac{J_0(z_n)+\frac{r_Um_n}{a}J_1(z_n)}{Y_0(z_n)+\frac{r_Um_n}{a}Y_1(z_n)}}\right]_0=\left[{\displaystyle \frac{J_0(z_n)\frac{r_Um_n}{a}J_1(z_n)}{Y_0(z_n)\frac{r_Um_n}{a}Y_1(z_n)}}\right]_L.`$ (45)
These results are identical to those of , computed in the orbifold picture.
## 3 Gravity on a flat orbifold
### 3.1 Orbifold picture
Consider 5d gravity on an $`๐^\mathrm{๐}/๐_\mathrm{๐}`$ orbifold in the simplest case where there are no brane sources and there is no bulk cosmological constant. This would be, e.g., the gravity background for the simplest model of Universal Extra Dimensions . Letโs find the physical degrees of freedom coming from the 5d metric, using the conventional orbifold language. The background metric is flat:
$`G_{\mu \nu }^\mathrm{๐}=\eta _{\mu \nu };G_{\mu 4}^\mathrm{๐}=0;G_{44}^\mathrm{๐}=1.`$ (46)
Including linearized metric fluctuations, we write:
$`G_{MN}=G_{MN}^\mathrm{๐}+h_{MN}=\left(\begin{array}{cc}\eta _{\mu \nu }& 0\\ 0& 1\end{array}\right)+\left(\begin{array}{cc}h_{\mu \nu }& h_{\mu 4}\\ h_{4\nu }& h_{44}\end{array}\right).`$ (47)
Plugging this into the standard source-free 5d Einstein equation gives the following bulk equations of motion:
$`0`$ $`=`$ $`_P_\mu h_\nu ^P+_P_\nu h_\mu ^P^P_Ph_{\mu \nu }_\mu _\nu h_M^M`$ (48)
$`\eta _{\mu \nu }\left(^M^Nh_{MN}^P_Ph_M^M\right);`$
$`0`$ $`=`$ $`_P_\mu h_4^P+_Ph_{\mu }^{P}{}_{}{}^{}^P_Ph_{\mu 4}_\mu h_{M}^{M}{}_{}{}^{};`$ (49)
$`0`$ $`=`$ $`2_Ph_{4}^{P}{}_{}{}^{}^P_Ph_{44}h_{M}^{M}{}_{}{}^{\prime \prime }^M^Nh_{MN}+^P_Ph_M^M,`$ (50)
where as always a prime indicates derivative with respect to $`y`$.
Because of the periodicity in $`y`$, all of the metric fluctuations are expanded in a tower of KK modes, which are just sines and cosines. In the orbifold picture, we impose the $`๐_\mathrm{๐}`$ symmetry explicitly, by projecting out the sine modes for $`h_{\mu \nu }`$ and $`h_{44}`$, as well as the cosine modes for $`h_{\mu 4}`$:
$`h_{\mu \nu }(x,y)`$ $`=`$ $`{\displaystyle \frac{h_{\mu \nu }^{(0)}(x)}{\sqrt{2}}}+{\displaystyle \underset{n>0}{\overset{\mathrm{}}{}}}h_{\mu \nu }^{(n)}(x)\mathrm{cos}\left({\displaystyle \frac{n\pi y}{L}}\right);`$
$`h_{\mu 4}(x,y)`$ $`=`$ $`{\displaystyle \underset{n>0}{\overset{\mathrm{}}{}}}h_{\mu 4}^{(n)}(x)\mathrm{sin}\left({\displaystyle \frac{n\pi y}{L}}\right);`$ (51)
$`h_{44}(x,y)`$ $`=`$ $`{\displaystyle \frac{h_{44}^{(0)}(x)}{\sqrt{2}}}+{\displaystyle \underset{n>0}{\overset{\mathrm{}}{}}}h_{44}^{(n)}(x)\mathrm{cos}\left({\displaystyle \frac{n\pi y}{L}}\right).`$
Note there are ten zero modes from $`h_{\mu \nu }`$ and one more from $`h_{44}`$.
Linearized general coordinate transformations (GCTs) $`x^Mx^M+\xi ^M`$ give:
$`h_{\mu \nu }`$ $``$ $`h_{\mu \nu }_\mu \xi _\nu _\nu \xi _\mu ,`$ (52)
$`h_{\mu 4}`$ $``$ $`h_{\mu 4}\xi _\mu {}_{}{}^{}_\mu \xi ^4,`$ (53)
$`h_{44}`$ $``$ $`h_{44}2\xi ^4{}_{}{}^{},`$ (54)
where the gauge parameters $`\xi ^M(x,y)`$ are also expanded in KK modes and subjected to $`๐_\mathrm{๐}`$ projections:
$`\xi ^\mu (x,y)`$ $`=`$ $`{\displaystyle \frac{\xi ^{\mu (0)}(x)}{\sqrt{2}}}+{\displaystyle \underset{n>0}{\overset{\mathrm{}}{}}}\xi ^{\mu (n)}(x)\mathrm{cos}\left({\displaystyle \frac{n\pi y}{L}}\right);`$
$`\xi ^4(x,y)`$ $`=`$ $`{\displaystyle \underset{n>0}{\overset{\mathrm{}}{}}}\xi ^{4(n)}(x)\mathrm{sin}\left({\displaystyle \frac{n\pi y}{L}}\right).`$ (55)
We can go to a convenient coordinate system by choosing
$`{\displaystyle \frac{2n\pi }{L}}\xi ^{4(n)}`$ $`=`$ $`h_{44}^{(n)},`$ (56)
$`{\displaystyle \frac{n\pi }{L}}\xi _\mu ^{(n)}`$ $`=`$ $`h_{\mu 4}^{(n)}+{\displaystyle \frac{L}{2n\pi }}_\mu h_{44}^{(n)},`$ (57)
so that all $`h_{\mu 4}^{(n)}`$ and $`h_{44}^{(n)}`$ with $`n>0`$ are gauged away, leaving only
$`h_{\mu \nu }(x,y)`$ $`=`$ $`h_{\mu \nu }^{(0)}(x)+{\displaystyle \underset{n>0}{\overset{\mathrm{}}{}}}h_{\mu \nu }^{(n)}(x)\mathrm{cos}\left({\displaystyle \frac{n\pi y}{L}}\right);`$ (58)
$`h_{44}(x,y)`$ $`=`$ $`h_{44}^{(0)}(x).`$ (59)
The residual gauge freedom is generated by $`\xi ^M(x,y)`$ that satisfy
$`\xi ^4{}_{}{}^{}=0,\xi _\mu ^{}=_\mu \xi ^4.`$ (60)
In addition, we began with a coordinate system in which the branes are straight, i.e., they are located at fixed slices of $`y`$. In the orbifold picture it is not obvious whether we are allowed to deviate from this โstraight gaugeโ. In the literature it is usually assumed (implicitly) that one should not deviate from straight gauges. This assumption implies that $`\xi ^4`$ should vanish at the brane locations, reducing (60) to:
$`\xi ^4=0,\xi _\mu ^{}=0.`$ (61)
Thus the remaining gauge freedom is just the 4d GCTs generated by $`\xi ^{\mu (0)}(x)`$. The equations of motion for the gauge-fixed degrees of freedom can be decoupled by a standard analysis . From (58) we read off the physical degrees of freedom (DOF):
* a massless graviton $`h_{\mu \nu }^{(0)}(x)`$ with two on-shell degrees of freedom,
* a massless radion $`h_{44}^{(0)}(x)`$,
* a Kaluza-Klein tower of massive gravitons $`h_{\mu \nu }^{(n)}(x)`$ with 5 DOF each.
### 3.2 Interval picture
In the interval picture we begin with an action:
$`S={\displaystyle d^4x\left(_{0^+}^L^{}๐y+_{L^+}^0^{}๐y\right)\sqrt{G}\mathrm{\hspace{0.17em}2}M^3R}+4M^3{\displaystyle _{}}K,`$ (62)
where $`R`$ is the Ricci scalar, $`K`$ is extrinsic curvature, and $`M`$ is the 5d Planck mass. We are using coordinates in which the branes (i.e. the boundaries) are straight, that is, they are located at fixed slices of $`y`$.
The second term in the action is the usual Gibbons-Hawking modification of GR for the case of manifolds with boundaries . The addition of this term ensures that the bulk EOM is the usual Einstein equation, for metric variations which vanish on the boundaries. Since we need an action principle for metric variations which do not necessarily vanish at the boundaries, we must supplement the bulk Einstein equation by appropriate brane-boundary equations.
We want to compare the orbifold results of the previous section with what we obtain using the interval picture. Keep in mind that in the interval picture we do not impose any $`๐_\mathrm{๐}`$ projections, on either the metric fluctuations or on the generators of general coordinate transformations. We will only quote results since a general derivation is given in ยง5. There we will also generalize to the case of non-straight gauges.
As in our previous example of the 5d gauge field, we begin by performing a partial gauge fixing. The GCT with $`\xi ^{(\mathrm{I})}{}_{}{}^{\mu }=0`$ and
$`\xi ^{(\mathrm{I})}{}_{}{}^{4}={\displaystyle \frac{1}{2}}{\displaystyle ^y}h_{44}dy{\displaystyle \frac{1}{2}}{\displaystyle ^y}F(y)\psi (x)dy,`$ (63)
with $`F(y)`$ a fixed but arbitrary function of $`y`$, transforms an arbitrary $`h_{44}`$ into
$`h_{44}=F(y)\psi (x).`$ (64)
Since we want to be in a straight gauge, we must require that $`\xi ^{(\mathrm{I})}^4`$ vanishes at the locations of the branes. On the interval $`0<y<L`$, this fixes the $`y`$-independent part of (63) to be
$`\xi ^{(\mathrm{I})}{}_{}{}^{4}={\displaystyle \frac{1}{2}}{\displaystyle _0^y}h_{44}dy{\displaystyle \frac{1}{2}}{\displaystyle _0^y}F(y)\psi (x)dy,`$ (65)
and fixes a relation between the radion field $`\psi (x)`$, $`F(y)`$, and the original metric fluctuation $`h_{44}(x,y)`$:
$`\psi (x)={\displaystyle \frac{_0^Lh_{44}๐y}{_0^LF(y)๐y}}.`$ (66)
From (66) we see that $`F(y)`$, though arbitrary, must be nonzero. More precisely, the straight gauge condition requires:
$`{\displaystyle _0^L}F(y)๐y0.`$ (67)
Thus the analog of axial gauge is not a straight gauge.
Next we can perform an additional partial gauge-fixing to eliminate $`h_{\mu 4}`$. Choose $`\xi ^{(\mathrm{II})}{}_{}{}^{4}=0`$ and
$`\xi ^{(\mathrm{II})}{}_{}{}^{\mu }={\displaystyle ^y}h^{\mu 4}dy,`$ (68)
which leaves $`h_{44}`$ unaffected and gives
$`h_{\mu 4}=0.`$ (69)
The remaining gauge freedom is just the 4d general coordinate transformation generated by
$`\xi ^4=0,\xi ^\mu =\xi ^\mu (x).`$ (70)
Note that the coordinate transformation generated by
$`\xi ^4ฯต(x),\xi ^\mu =y_\mu ฯต(x),`$ (71)
respects the gauge conditions (64) and (69) but does not keep us in a straight gauge. Treated as a scalar metric perturbation, $`ฯต(x)`$ is the putative brane-bending mode.
To identify the physical DOF, we examine the bulk equations of motion obtained from (48-50):
$`0`$ $`=`$ $`_\rho _\mu h_\nu ^\rho +_\rho _\nu h_\mu ^\rho ^2h_{\mu \nu }_\mu _\nu \stackrel{~}{h}`$
$`\eta _{\mu \nu }(_\rho _\sigma h^{\rho \sigma }^2\stackrel{~}{h})h_{\mu \nu }^{\prime \prime }+\eta _{\mu \nu }\stackrel{~}{h}^{\prime \prime }F_\mu _\nu \psi +\eta _{\mu \nu }F^2\psi ;`$
$`0`$ $`=`$ $`_\nu h_{\mu }^{\nu }{}_{}{}^{}_\mu \stackrel{~}{h}^{};`$ (72)
$`0`$ $`=`$ $`_\mu _\nu h^{\mu \nu }+^2\stackrel{~}{h};`$
$`0`$ $`=`$ $`\stackrel{~}{h}^{\prime \prime }+F^2\psi ,`$
where $`\stackrel{~}{h}=\eta ^{\mu \nu }h_{\mu \nu }`$, and the fourth equation is an auxiliary relation obtained from twice the third equation subtracted from the trace of the first. From the general formula (243) that we will derive in ยง5, the brane-boundary equations are
$`0=\left[h_{\mu \nu }^{}\eta _{\mu \nu }\stackrel{~}{h}^{}\right]_{y=0,L},`$ (73)
As in the gauge field example of the previous section, we will need to solve these equations separately for the cases where the metric perturbations are massless or massive in the 4d sense. To be completely explicit, we will Fourier transform to a 4d momentum space representation of the metric perturbations $`\overline{h}_{\mu \nu }(p,y)`$. The bulk and brane-boundary equations become:
$`0`$ $`=`$ $`p_\mu p_\rho \overline{h}_\nu ^\rho p_\nu p_\rho \overline{h}_\mu ^\rho +p^2\overline{h}_{\mu \nu }+p_\mu p_\nu \overline{h}`$ (74)
$`+\eta _{\mu \nu }(p_\rho p_\sigma \overline{h}^{\rho \sigma }p^2\overline{h})\overline{h}_{\mu \nu }^{\prime \prime }+\eta _{\mu \nu }\overline{h}^{\prime \prime }+Fp_\mu p_\nu \overline{\psi }\eta _{\mu \nu }Fp^2\overline{\psi };`$
$`0`$ $`=`$ $`p_\nu \overline{h}_\mu ^\nu {}_{}{}^{}p_\mu \overline{h}^{};`$ (75)
$`0`$ $`=`$ $`p_\mu p_\nu \overline{h}^{\mu \nu }p^2\overline{h};`$ (76)
$`0`$ $`=`$ $`\overline{h}^{\prime \prime }Fp^2\overline{\psi },`$ (77)
$`0`$ $`=`$ $`\left[\overline{h}_{\mu \nu }^{}\eta _{\mu \nu }\overline{h}^{}\right]_{y=0,L}.`$ (78)
#### 3.2.1 $`๐ฉ^\mathrm{๐}\mathrm{๐}`$
As discussed in the Appendix, for $`p^20`$ the tensor $`\overline{h}_{\mu \nu }(p,y)`$ can be decomposed as
$`\overline{h}_{\mu \nu }(p,y)=\overline{b}_{\mu \nu }(p,y)+ip_\mu \overline{V}_\nu (p,y)+ip_\nu \overline{V}_\mu (p,y)p_\mu p_\nu \overline{\varphi }_1(p,y)+\eta _{\mu \nu }\overline{\varphi }_2(p,y),`$ (79)
where $`\overline{b}_{\mu \nu }(p,y)`$ is traceless transverse, $`\overline{V}_\mu (p,y)`$ is transverse and pure gauge, $`\overline{\varphi }_1(p,y)`$ is a pure gauge scalar, and $`\overline{\varphi }_2(p,y)`$ is another scalar. The bulk equation (76) immediately gives the constraint
$`\overline{\varphi }_2=0.`$ (80)
Then (75) gives
$`ip^2\overline{V}_\mu ^{}=0.`$ (81)
Since $`\overline{V}_\mu `$ is both pure gauge and $`y`$-independent, we can gauge it away using the transverse modes of our residual gauge freedom (70).
Integrating (77) twice in $`y`$ gives
$`\overline{\varphi }_1=\overline{f}_2(p)+y\overline{f}_1(p)๐\overline{\psi },`$ (82)
where $`\overline{f}_1(p)`$, $`\overline{f}_2(p)`$ are integration โconstantsโ, and $`๐(y)`$ is defined by $`๐^{\prime \prime }(y)=F(y)`$, with no integration constants. We can remove $`\overline{f}_2(p)`$ using the longitudinal mode of the residual gauge freedom (70).
The trace of the brane-boundary equation (78) gives
$`\left[3p^2\overline{\varphi }_1^{}\right]_{y=0,L}=\left[3p^2(\overline{f}_1\overline{\psi })\right]_{y=0,L}=0,`$ (83)
with $`^{}(y)=F(y)`$. Due to (67), $`(0)(L)`$, and thus the only solution of (83) is
$`\overline{f}_1(p)=0,\overline{\psi }(p)=0.`$ (84)
Now only $`\overline{b}_{\mu \nu }`$ is left, and (74) and (78) determine its mass spectrum. By going on-shell, i.e., $`p^2=m^2`$:
$`m^2\overline{b}_{\mu \nu }+\overline{b}_{\mu \nu }^{\prime \prime }=0,\left[\overline{b}_{\mu \nu }^{}\right]_{y=0,L}=0`$
$`\overline{b}_{\mu \nu }(p,y)=\overline{B}_{\mu \nu }(p)\mathrm{cos}{\displaystyle \frac{n\pi }{L}}y,n=1,2,\mathrm{},`$ (85)
which agrees with the results obtained for the massive sector in the orbifold approach.
#### 3.2.2 $`๐ฉ^\mathrm{๐}=\mathrm{๐}`$
The massless modes of $`\overline{h}_{\mu \nu }`$ have the more complicated tensor decomposition given by (313):
$`\overline{h}_{\mu \nu }(p,y)`$ $`=`$ $`\overline{\beta }_{\mu \nu }(p,y)+ip_\mu \overline{v}_\nu (p,y)+ip_\nu \overline{v}_\mu (p,y)p_\mu p_\nu \overline{\phi }_1(p,y)`$ (86)
$`+ip_\mu \overline{n}_\nu (p,y)+ip_\nu \overline{n}_\mu (p,y)+\overline{c}_{\mu \nu }(p,y)+\eta _{\mu \nu }\overline{\phi }_2(p,y).`$
Here $`\overline{v}_\mu (p,y)`$ is transverse and pure gauge (2 DOF), $`\overline{\phi }_1(p,y)`$ is also pure gauge (1 DOF), and $`\overline{n}_\mu (p,y)`$ is pure gauge but not transverse (1 DOF). Also, $`\overline{c}_{\mu \nu }(p,y)`$ is traceless but not transverse (3 DOF), $`\overline{\phi }_2(p,y)`$ is a scalar (1 DOF), and $`\overline{\beta }_{\mu \nu }(p,y)`$ are the remaining traceless transverse components (2 DOF).
The bulk equation (76) gives the constraint:
$`p_\mu p_\nu \overline{c}^{\mu \nu }=0,`$ (87)
which, because $`p_\nu \overline{c}^{\mu \nu }0`$, implies
$`\overline{c}^{\mu \nu }=0.`$ (88)
Then, (75) becomes
$`p_\mu (ip_\nu \overline{n}^\nu {}_{}{}^{}+3\overline{\phi }_2^{})=0ip_\nu \overline{n}^\nu {}_{}{}^{}+3\overline{\phi }_2^{}=0,`$ (89)
while (77) gives
$`ip_\nu \overline{n}^\nu {}_{}{}^{\prime \prime }+2\overline{\phi }_2^{\prime \prime }=0ip_\nu \overline{n}^\nu {}_{}{}^{}+2\overline{\phi }_2^{}=\overline{f}_1(p).`$ (90)
Solving for $`p_\nu \overline{n}^\nu ^{}`$ and $`\overline{\phi }_2^{}`$, we get
$`ip_\nu \overline{n}^\nu {}_{}{}^{}=3\overline{f}_1,\overline{\phi }_2^{}=\overline{f}_1.`$ (91)
Taking the trace of the brane-boundary equation (78):
$`0=\left[\overline{h}^{}\right]_{y=0,L}=\left[2ip_\nu \overline{n}^\nu {}_{}{}^{}+4\overline{\phi }_2^{}\right]_{y=0,L}=2\overline{f}_1.`$ (92)
Then
$`\overline{\phi }_2=\overline{f}_2(p),`$ (93)
and due to $`p_\mu \overline{n}^\mu 0`$,
$`\overline{n}^\mu {}_{}{}^{}=0.`$ (94)
Since $`\overline{n}^\mu `$ is pure gauge and $`y`$-independent, it can be eliminated by the longitudinal part of the residual gauge freedom.
Finally, (74) gives
$`\overline{\beta }_{\mu \nu }^{\prime \prime }+ip_\mu \overline{v}_\nu ^{\prime \prime }+ip_\nu \overline{v}_\mu ^{\prime \prime }p_\mu p_\nu (\overline{\phi }_1^{\prime \prime }+2\overline{f}_2+F\overline{\psi })=0.`$ (95)
Now we contract (95) with $`\overline{n}^\mu \overline{n}^\nu `$. In the Appendix, we show that $`\overline{n}^\mu \overline{\beta }_{\mu \nu }=0`$, $`\overline{n}^\mu \overline{v}_\mu =0`$, and $`p_\mu \overline{n}^\mu 0`$; these are 4d tensor relations which are unchanged if we replace $`\overline{\beta }_{\mu \nu }`$ by $`\overline{\beta }_{\mu \nu }^{\prime \prime }`$ or $`\overline{v}_\mu `$ by $`\overline{v}_\mu ^{\prime \prime }`$. Thus we get
$`\overline{\phi }_1^{\prime \prime }+2\overline{f}_2+F\overline{\psi }=0,`$ (96)
and
$`\overline{\beta }_{\mu \nu }^{\prime \prime }+ip_\mu \overline{v}_\nu ^{\prime \prime }+ip_\nu \overline{v}_\mu ^{\prime \prime }=0.`$ (97)
For convenience, letโs define $`\overline{t}_{\mu \nu }=\overline{\beta }_{\mu \nu }+ip_\mu \overline{v}_\nu +ip_\nu \overline{v}_\mu `$. Then (78) gives
$`0=\left[\overline{t}_{\mu \nu }{}_{}{}^{}p_\mu p_\nu \overline{\phi }_1^{}\right]_{y=0,L},`$ (98)
thus
$`0`$ $`=`$ $`[\overline{t}_{\mu \nu }{}_{}{}^{}]_{y=0,L},`$ (99)
$`0`$ $`=`$ $`\left[\overline{\phi }_1^{}\right]_{y=0,L}.`$ (100)
From (97) and (99), we see that $`\overline{t}_{\mu \nu }`$ is $`y`$-independent. Then since $`\overline{v}_\nu `$ is pure gauge, we can gauge it away using two transverse components of the residual gauge freedom.
Next, from (96) we get
$`\overline{\phi }_1(p,y)=\overline{f}_4(p)+y\overline{f}_3(p)y^2\overline{f}_2(p)๐(y)\overline{\psi }.`$ (101)
Since $`\overline{\phi }_1`$ is pure gauge, we can eliminate $`\overline{f}_4`$ by the remaining transverse component of the residual gauge freedom. Using (100), $`\overline{\phi }_1`$ and $`\overline{f}_2`$ can be written in terms of $`\overline{\psi }`$:
$`\overline{f}_2`$ $`=`$ $`{\displaystyle \frac{(L)(0)}{2L}}\overline{\psi },`$ (102)
$`\overline{\phi }_1`$ $`=`$ $`\left\{{\displaystyle \frac{(L)(0)}{2L}}y^2(๐(y)(0)y)\right\}\overline{\psi }.`$ (103)
Recall that $`F(y)`$ is an arbitrary function satisfying (67). For example we can choose $`F(y)=1`$, in which case the above reduces to $`f_2(x)=\psi (x)/2`$, $`\phi _1(x)=0`$.
In short, the physical degrees of freedom of the massless sector consist of a massless graviton $`\beta _{\mu \nu }(x)`$ with two on-shell degrees of freedom, together with the massless radion $`\psi (x)`$. This agrees with the results of the orbifold approach.
## 4 The Randall-Sundrum model in the interval picture
Letโs repeat the same exercise for the case of the RSI background . We want to reproduce the well-known results from the orbifold approach using the interval picture analysis. Here we have a nonzero bulk cosmological constant and brane tensions, which are tuned to give a warped background solution with flat 4d slices. The interval picture action is
$`{\displaystyle \frac{S}{2M^3}}`$ $`=`$ $`{\displaystyle d^4x\left(_{0^+}^L^{}๐y+_{L^+}^0^{}๐y\right)\sqrt{G}\left(R+12k^2\right)}`$ (104)
$`+{\displaystyle \underset{i}{}}{\displaystyle _{y=y_i}}d^4x\sqrt{g}\mathrm{\hspace{0.33em}12}k\theta _i+2{\displaystyle _{}}K,`$
where $`\theta _1=\theta _2=1`$, $`y_1=0`$, $`y_2=L`$, and we have already inserted the tuned values for the two brane tensions. The background solution is
$`G_{MN}^\mathrm{๐}=\left(\begin{array}{cc}g_{\mu \nu }^\mathrm{๐}& 0\\ 0& 1\end{array}\right),`$ (107)
where $`g_{\mu \nu }^\mathrm{๐}=a^2(y)\eta _{\mu \nu }`$ with
$`a(y)=\{\begin{array}{cc}\mathrm{e}^{ky}\hfill & L<y<0,\hfill \\ \mathrm{e}^{ky}\hfill & 0<y<L.\hfill \end{array}`$ (110)
As in all our examples we will restrict our attention to the interval $`0<y<L`$. Indices are raised and lowered with the warped background metric $`G_{MN}^\mathrm{๐}`$. Following the procedure of the previous section, we gauge fix to
$`h_{\mu 4}=0,h_{44}=F(y)\psi (x),`$ (111)
with the residual gauge freedom generated by
$`\xi ^4=0,\xi ^\mu =\xi ^\mu (x).`$ (112)
Bulk and brane-boundary equations of motion are obtained from (239-243):
$`0`$ $`=`$ $`_\rho _\mu h_\nu ^\rho +_\rho _\nu h_\mu ^\rho ^2h_{\mu \nu }_\mu _\nu \stackrel{~}{h}g_{\mu \nu }^\mathrm{๐}(_\rho _\sigma h^{\rho \sigma }^2\stackrel{~}{h})`$
$`h_{\mu \nu }^{\prime \prime }+g_{\mu \nu }^\mathrm{๐}\stackrel{~}{h}^{\prime \prime }4kg_{\mu \nu }^\mathrm{๐}\stackrel{~}{h}^{}+4k^2h_{\mu \nu }`$
$`F_\mu _\nu \psi +g_{\mu \nu }^\mathrm{๐}F^2\psi +3kg_{\mu \nu }^\mathrm{๐}F^{}\psi 12k^2g_{\mu \nu }^\mathrm{๐}F\psi ,`$
$`0`$ $`=`$ $`_\nu h_\mu ^\nu {}_{}{}^{}_\mu \stackrel{~}{h}^{}3kF_\mu \psi ,`$
$`0`$ $`=`$ $`_\mu _\nu h^{\mu \nu }+^2\stackrel{~}{h}3k\stackrel{~}{h}^{}12k^2F\psi ,`$ (113)
$`0`$ $`=`$ $`{\displaystyle \frac{(a^2\stackrel{~}{h}^{})^{}}{a^2}}+F^2\psi +4kF^{}\psi 8k^2F\psi ,`$
$`0`$ $`=`$ $`[h_{\mu \nu }{}_{}{}^{}g_{\mu \nu }^\mathrm{๐}\stackrel{~}{h}^{}+2kh_{\mu \nu }3kg_{\mu \nu }^\mathrm{๐}F\psi ]_{y=y_i},`$
with $`^2=g^\mathrm{๐}{}_{}{}^{\mu \nu }_{\mu }^{}_\nu `$, $`\stackrel{~}{h}=g^\mathrm{๐}{}_{}{}^{\mu \nu }h_{\mu \nu }^{}`$.
Although the spacetime is warped, we can still perform a 4d Fourier analysis on the flat 4d slices. Using $`p^2`$ to denote $`a^2g^\mathrm{๐}{}_{}{}^{\mu \nu }p_{\mu }^{}p_\nu `$ we write:
$`0`$ $`=`$ $`p_\rho p_\mu \overline{h}_\nu ^\rho p_\rho p_\nu \overline{h}_\mu ^\rho +a^2p^2\overline{h}_{\mu \nu }+p_\mu p_\nu \overline{h}+g_{\mu \nu }^\mathrm{๐}(p_\rho p_\sigma \overline{h}^{\rho \sigma }a^2p^2\overline{h})`$
$`\overline{h}_{\mu \nu }^{\prime \prime }+g_{\mu \nu }^\mathrm{๐}\overline{h}^{\prime \prime }4kg_{\mu \nu }^\mathrm{๐}\overline{h}^{}+4k^2\overline{h}_{\mu \nu }`$
$`+Fp_\mu p_\nu \overline{\psi }g_{\mu \nu }^\mathrm{๐}Fa^2p^2\overline{\psi }+3kg_{\mu \nu }^\mathrm{๐}F^{}\overline{\psi }12k^2g_{\mu \nu }^\mathrm{๐}F\overline{\psi },`$
$`0`$ $`=`$ $`p_\nu \overline{h}_\mu ^\nu {}_{}{}^{}p_\mu \overline{h}^{}3kFp_\mu \overline{\psi },`$ (115)
$`0`$ $`=`$ $`p_\mu p_\nu \overline{h}^{\mu \nu }a^2p^2\overline{h}3k\overline{h}^{}12k^2F\overline{\psi },`$ (116)
$`0`$ $`=`$ $`{\displaystyle \frac{(a^2\overline{h}^{})^{}}{a^2}}Fa^2p^2\overline{\psi }+4kF^{}\overline{\psi }8k^2F\overline{\psi },`$ (117)
$`0`$ $`=`$ $`[\overline{h}_{\mu \nu }{}_{}{}^{}g_{\mu \nu }^\mathrm{๐}\overline{h}^{}+2k\overline{h}_{\mu \nu }3kg_{\mu \nu }^\mathrm{๐}F\overline{\psi }]_{y=y_i}.`$ (118)
### 4.1 $`๐ฉ^\mathrm{๐}\mathrm{๐}`$
We use the tensor decomposition from the Appendix:
$`\overline{h}_{\mu \nu }=\overline{b}_{\mu \nu }+ip_\mu \overline{V}_\nu +ip_\nu \overline{V}_\mu a^2p_\mu p_\nu \overline{\varphi }_1+g_{\mu \nu }^\mathrm{๐}\overline{\varphi }_2,`$ (119)
where we put $`a^2`$ in front of $`\overline{\varphi }_1`$ for convenience. Integrating (117), we get
$`\overline{h}^{}=p^2\overline{\varphi }_1^{}(p,y)+4\overline{\varphi }_2^{}(p,y)=e^{2ky}\overline{f}_1(p)+e^{2ky}(y)p^2\overline{\psi }(p)4kF(y)\overline{\psi }(p),`$ (120)
with $`^{}(y)=F(y)`$. Then (116) and (115) are written as
$`p^2\overline{\varphi }_2(p,y)`$ $`=`$ $`k\left(\overline{f}_1(p)+p^2\overline{\psi }(p)\right),`$ (121)
$`\left(e^{2ky}p^2\overline{V}_\mu (p,y)\right)^{}3ip_\mu \overline{\varphi }_2^{}(p,y)`$ $`=`$ $`3ikFp_\mu \overline{\psi }(p).`$ (122)
We can solve (120-122) for $`\overline{\varphi }_1`$, $`\overline{\varphi }_2`$ and $`\overline{V}_\mu `$, to get
$`\overline{\varphi }_1`$ $`=`$ $`\overline{f}_2(p){\displaystyle \frac{1}{p^2}}\left({\displaystyle \frac{4k}{p^2}}+{\displaystyle \frac{e^{2ky}}{2k}}\right)\overline{f}_1๐\overline{\psi },`$ (123)
$`\overline{\varphi }_2`$ $`=`$ $`k\left({\displaystyle \frac{1}{p^2}}\overline{f}_1+\overline{\psi }\right),`$ (124)
$`0`$ $`=`$ $`p^2\left(e^{2ky}\overline{V}_\mu \right)^{},`$ (125)
where we have defined $`๐^{}(y)=e^{2ky}(y)`$.
Equation (125) fixes the $`y`$-dependence of $`\overline{V}_\mu `$ to be $`e^{2ky}`$, which allows us to eliminate $`\overline{V}_\mu `$ by the transverse part of the residual gauge freedom. Similarly, $`\overline{f}_2(p)`$ is removed by the longitudinal part of the residual gauge freedom. Then, $`\overline{h}_{\mu \nu }`$ becomes
$`\overline{h}_{\mu \nu }=\overline{b}_{\mu \nu }+e^{2ky}p_\mu p_\nu \left\{{\displaystyle \frac{1}{p^2}}\left({\displaystyle \frac{4k}{p^2}}+{\displaystyle \frac{e^{2ky}}{2k}}\right)\overline{f}_1+๐\overline{\psi }\right\}kg_{\mu \nu }^\mathrm{๐}\left({\displaystyle \frac{1}{p^2}}\overline{f}_1+\overline{\psi }\right).`$ (126)
Plugging this into the Fourier-transformed version of the bulk $`\mu \nu `$-EOM and the boundary EOM, we have
$`0`$ $`=`$ $`e^{2ky}p^2\overline{b}_{\mu \nu }\overline{b}_{\mu \nu }^{\prime \prime }+4k^2\overline{b}_{\mu \nu },`$ (127)
$`0`$ $`=`$ $`\left[\overline{b}_{\mu \nu }^{}+2k\overline{b}_{\mu \nu }+(p_\mu p_\nu \eta _{\mu \nu }p^2)\left({\displaystyle \frac{1}{p^2}}\overline{f}_1+\overline{\psi }\right)\right]_{y=y_i}.`$ (128)
Contracting (128) with $`g^\mathrm{๐}^{\mu \nu }`$ gives
$`{\displaystyle \frac{1}{p^2}}\overline{f}_1+(0)\overline{\psi }=0,{\displaystyle \frac{1}{p^2}}\overline{f}_1+(L)\overline{\psi }=0.`$ (129)
Since $`(0)(L)`$, this implies $`\overline{f}_1=\overline{\psi }=0`$.
Going on-shell, we substitute $`m^2`$ for $`p^2`$:
$`0`$ $`=`$ $`\overline{b}_{\mu \nu }^{\prime \prime }+\left(m^2e^{2ky}4k^2\right)\overline{b}_{\mu \nu },`$ (130)
$`0`$ $`=`$ $`\left[\overline{b}_{\mu \nu }^{}+2k\overline{b}_{\mu \nu }\right]_{y=y_i}.`$ (131)
The solution of (130) is
$`\overline{b}_{\mu \nu }(p,y)=\overline{A}_{\mu \nu }(p)J_2\left({\displaystyle \frac{m}{k}}e^{ky}\right)+\overline{B}_{\mu \nu }(p)Y_2\left({\displaystyle \frac{m}{k}}e^{ky}\right),`$ (132)
where $`J_n(Y_n)`$ is the Bessel function of the first(second) kind. Equation (131) provides boundary conditions:
$`0`$ $`=`$ $`m\left\{\overline{A}_{\mu \nu }(p)J_1\left({\displaystyle \frac{m}{k}}\right)+\overline{B}_{\mu \nu }(p)Y_1\left({\displaystyle \frac{m}{k}}\right)\right\},`$ (133)
$`0`$ $`=`$ $`me^{kL}\left\{\overline{A}_{\mu \nu }(p)J_1\left({\displaystyle \frac{m}{k}}e^{kL}\right)+\overline{B}_{\mu \nu }(p)Y_1\left({\displaystyle \frac{m}{k}}e^{kL}\right)\right\}.`$ (134)
These can have a non-trivial solution only when
$`J_1\left({\displaystyle \frac{m}{k}}\right)Y_1\left({\displaystyle \frac{m}{k}}e^{kL}\right)Y_1\left({\displaystyle \frac{m}{k}}\right)J_1\left({\displaystyle \frac{m}{k}}e^{kL}\right)=0,`$ (135)
which determines the discrete spectrum of massive graviton modes. Then, (132) becomes
$`\overline{b}_{\mu \nu }(p,y)=\overline{B}_{\mu \nu }(p)\left\{Y_1\left({\displaystyle \frac{m}{k}}\right)J_2\left({\displaystyle \frac{m}{k}}e^{ky}\right)J_1\left({\displaystyle \frac{m}{k}}\right)Y_2\left({\displaystyle \frac{m}{k}}e^{ky}\right)\right\},`$ (136)
up to an overall normalization. Note $`\overline{b}_{\mu \nu }(p,y)`$ has five polarizations from the transverse-traceless $`\overline{B}_{\mu \nu }`$. Thus the physical content of the massive sector is a Kaluza-Klein tower of massive gravitons coming from $`b_{\mu \nu }(x,y)`$, in agreement with the Randall-Sundrum result.
### 4.2 $`๐ฉ^\mathrm{๐}=\mathrm{๐}`$
We use the same massless tensor decomposition as in the flat orbifold case:
$`\overline{h}_{\mu \nu }=\overline{t}_{\mu \nu }a^2p_\mu p_\nu \overline{\phi }_1+ip_\mu \overline{n}_\nu +ip_\nu \overline{n}_\mu +\overline{c}_{\mu \nu }+g_{\mu \nu }^\mathrm{๐}\overline{\phi }_2,`$ (137)
where $`\overline{t}_{\mu \nu }=\overline{\beta }_{\mu \nu }+ip_\mu \overline{v}_\nu +ip_\nu \overline{v}_\mu `$. Equation (117) gives
$`2ip_\mu \overline{n}^\mu {}_{}{}^{}+4\overline{\phi }_2^{}=e^{2ky}\overline{f}_1(p)4kF(y)\overline{\psi }(p),`$ (138)
and (116) and (115) become
$`p_\mu p_\nu \overline{c}^{\mu \nu }`$ $`=`$ $`3ke^{2ky}\overline{f}_1,`$ (139)
$`p_\nu \overline{c}_\mu ^\nu ^{}`$ $`=`$ $`p_\mu (e^{2ky}\overline{f}_1kF\overline{\psi }ip_\nu \overline{n}^\nu {}_{}{}^{}\overline{\phi }_2^{}).`$ (140)
Now we contract (140) with $`\overline{n}^\mu `$. In the Appendix, we show that $`\overline{n}^\mu \overline{c}_{\mu \nu }=0`$ and $`p_\mu \overline{n}^\mu 0`$; these are 4d tensor relations which are unchanged if replace $`\overline{n}^\mu `$ by $`\overline{n}^\mu ^{}`$ or $`\overline{c}_{\mu \nu }`$ by $`\overline{c}_{\mu \nu }^{}`$. Thus we get
$`ip_\nu \overline{n}^\nu {}_{}{}^{}+\overline{\phi }_2^{}=e^{2ky}\overline{f}_1kF\overline{\psi }.`$ (141)
Solving this and (138) for $`p_\nu \overline{n}^\nu ^{}`$ and $`\overline{\phi }_2^{}`$,
$`ip_\nu \overline{n}^\nu ^{}`$ $`=`$ $`{\displaystyle \frac{3}{2}}e^{2ky}\overline{f}_1,`$ (142)
$`\overline{\phi }_2^{}`$ $`=`$ $`{\displaystyle \frac{e^{2ky}}{2}}\overline{f}_1kF\overline{\psi }.`$ (143)
Since the trace of (118) gives
$`\left[3\overline{h}^{}12kF\overline{\psi }\right]_{y=y_i}=\left[3e^{2ky}\overline{f}_1\right]_{y=y_i}=0\overline{f}_1(p)=0,`$ (144)
then (139) dictates
$`\overline{c}_{\mu \nu }=0,`$ (145)
and from (142), we see $`\overline{n}_\mu `$ can be gauged away by the longitudinal part of the residual gauge freedom.
Equation (143) is integrated to give
$`\overline{\phi }_2=\overline{f}_2(p)k\overline{\psi }.`$ (146)
Then, contracting (4),
$`0`$ $`=`$ $`\overline{t}_{\mu \nu }{}_{}{}^{\prime \prime }+4k^2\overline{t}_{\mu \nu }+p_\mu p_\nu \{e^{2ky}\overline{\phi }_1^{\prime \prime }4ke^{2ky}\overline{\phi }_1^{}+2\overline{f}_2+(F2k)\overline{\psi }\},`$ (147)
with $`\overline{n}^\mu \overline{n}^\nu `$, we get
$`0`$ $`=`$ $`\overline{t}_{\mu \nu }{}_{}{}^{\prime \prime }+4k^2\overline{t}_{\mu \nu },`$ (148)
$`0`$ $`=`$ $`e^{2ky}\overline{\phi }_1^{\prime \prime }4ke^{2ky}\overline{\phi }_1^{}+2\overline{f}_2+\left(F2k\right)\overline{\psi }.`$ (149)
Similarly for (118);
$`0`$ $`=`$ $`\left[\overline{t}_{\mu \nu }^{}+2k\overline{t}_{\mu \nu }\right]_{y=y_i},`$ (150)
$`0`$ $`=`$ $`\left[e^{2ky}\overline{\phi }_1^{}\right]_{y=y_i}.`$ (151)
Using (148) and (150), we obtain
$`\overline{t}_{\mu \nu }(p,y)=\overline{B}_{\mu \nu }(p)e^{2ky}.`$ (152)
This means that $`\overline{v}_\mu `$ has the correct $`y`$-dependence to be gauged away by two transverse components of the residual gauge freedom.
Finally, solving (149), we get
$`\overline{\phi }_1=\overline{f}_4(p)+e^{4ky}\overline{f}_3(p)+{\displaystyle \frac{e^{2ky}}{2k^2}}\overline{f}_2(p)๐\overline{\psi }(p),`$ (153)
where we can gauge away $`\overline{f}_4`$ by the remaining transverse component of the residual gauge freedom. Using (151), $`\overline{f}_2`$ and $`\overline{f}_3`$ can be written in terms of $`\overline{\psi }`$:
$`\overline{f}_2=k{\displaystyle \frac{e^{2kL}(0)(L)}{e^{2kL}1}}\overline{\psi },\overline{f}_3={\displaystyle \frac{(L)(0)}{4k(e^{2kL}1)}}\overline{\psi },`$ (154)
and then
$`\overline{\phi }_1=\left\{{\displaystyle \frac{(L)(0)}{4k(e^{2kL}1)}}e^{4ky}+{\displaystyle \frac{e^{2kL}(0)(L)}{2k(e^{2kL}1)}}e^{2ky}๐\right\}\overline{\psi }.`$ (155)
Thus all the surviving scalars are linearly dependent on $`\overline{\psi }`$. Since $`F(y)`$ is an arbitrary function satisfying (67), we can simplify the above expressions. For example, choosing $`F(y)=1/a^2`$, the above reduces to $`f_2(x)=0`$, $`\phi _1(x)=0`$, and $`h_{\mu \nu }=a^2(y)B_{\mu \nu }(x)(1/2)\eta _{\mu \nu }\psi (x)`$.
We see that the physical content of the massless sector consists of a massless graviton $`B_{\mu \nu }(x)`$ with two on-shell degrees of freedom, and a massless radion $`\psi (x)`$. This agrees with the standard results .
## 5 Gravity in a general warped background
We are interested in warped background solutions which are generalizations of the original two brane setup of Randall and Sundrum . We have a 5d spacetime, $``$, which extends to infinity along the usual (1+3) dimensions (denoted by $`x^\mu `$) and has an extra spatial dimension (denoted by $`y`$) compactified on a circle with circumference $`2L`$. There are two branes, which are nonintersecting codimension one hypersurfaces described by $`\mathrm{\Phi }_1(x,y)=0`$ and $`\mathrm{\Phi }_2(x,y)=0`$. The branes divide the 5d spacetime $``$ into two pieces: $`_1`$, which extends from $`\mathrm{\Phi }_1(x,y)=0^+`$ to $`\mathrm{\Phi }_2(x,y)=0^{}`$, and $`_2`$, which extends from $`\mathrm{\Phi }_1(x,y)=0^{}`$ to $`\mathrm{\Phi }_2(x,y)=0^+`$. The branes have tension, which may be positive or negative, and the brane actions have kinetic terms for gravity, which in a complete model would be induced by radiative corrections involving brane matter -.
The bulk part of the action will be written with a bulk metric
$$G_{MN}=\left(\begin{array}{cc}g_{\mu \nu }& G_{\mu 4}\\ G_{4\nu }& G_{44}\end{array}\right),$$
whereas brane parts are written in terms of the induced metric
$`g_{\alpha \beta }^{(i)}=\left[{\displaystyle \frac{x^M}{x^{(i)\alpha }}}{\displaystyle \frac{x^N}{x^{(i)\beta }}}G_{MN}\right]_{\mathrm{\Phi }_i=0},`$
with $`x^{(i)\alpha }`$ a coordinate on the boundary hypersurface $`\mathrm{\Phi }_i=0`$, i.e., a โbrane coordinateโ. Since the superscript <sup>(i)</sup> on any entity always implies that it is evaluated on the $`\mathrm{\Phi }_i=0`$ hypersurface, we will omit $`[]_{\mathrm{\Phi }_i=0}`$ hereafter unless there is room for confusion. The inverse of the bulk and induced metrics satisfy the relation :
$`[G^{MN}]_{\mathrm{\Phi }_i=0}=N^{(i)M}N^{(i)N}+g^{(i)}{}_{}{}^{\alpha \beta }{\displaystyle \frac{x^M}{x^{(i)\alpha }}}{\displaystyle \frac{x^N}{x^{(i)\beta }}},`$ (156)
where $`N^{(i)M}`$ is the unit vector outward-normal to $`\mathrm{\Phi }_i=0`$, which can be written as
$`N_M^{(i)}={\displaystyle \frac{\theta _i_M\mathrm{\Phi }_i}{\sqrt{G^{PQ}_P\mathrm{\Phi }_i_Q\mathrm{\Phi }_i}}}.`$ (157)
Our convention for โoutwardโ is that $`\theta _1`$ is chosen to be $`1`$, while $`\theta _2`$ is $`+1`$.
The brane coordinate system together with $`N_M^{(i)}`$ naturally induces a bulk coordinate system on the brane, which we will call the โboundary normal coordinatesโ (BNCs) denoted by $`x^{(i)}^{\overline{M}}`$: on the $`\mathrm{\Phi }_i=0`$ hypersurface, we have $`x^{(i)\alpha }`$-coordinates and $`N_M^{(i)}`$ defining the directions orthogonal to them. Then at every point on the $`\mathrm{\Phi }_i=0`$ hypersurface, we choose the $`x^{(i)\overline{\alpha }}`$ to be in the directions of the $`x^{(i)\alpha }`$โs, and $`\overline{y}^{(i)}`$ to be in the direction of $`N_M^{(i)}`$. One of the useful features of this BNC is that since the $`\overline{y}^{(i)}`$-coordinate is orthogonal the to $`x^{(i)\overline{\alpha }}`$-ones,
$`G_{\overline{\alpha }\overline{4}}^{(i)}=0,`$ (158)
which, in turn, implies that $`\overline{\mu }`$-indices are raised and lowered by $`g_{\overline{\alpha }\overline{\beta }}^{(i)}`$ only, and the $`\overline{4}`$-index by $`G_{\overline{4}\overline{4}}^{(i)}`$ only. Note that the BNC is not necessarily the same as Gaussian normal coordinates on the brane, since we donโt require $`G_{\overline{4}\overline{4}}^{(i)}=1`$. Also
$`g_{\alpha \beta }^{(i)}={\displaystyle \frac{x^{(i)\overline{M}}}{x^{(i)\alpha }}}{\displaystyle \frac{x^{(i)\overline{N}}}{x^{(i)\beta }}}G_{\overline{M}\overline{N}}^{(i)}=\delta _{\overline{\alpha }}^{\overline{M}}\delta _{\overline{\beta }}^{\overline{N}}G_{\overline{M}\overline{N}}^{(i)}=g_{\overline{\alpha }\overline{\beta }}^{(i)},`$ (159)
i.e., the $`\overline{\alpha }\overline{\beta }`$-components of the bulk metric are the same as the induced metric. By construction we have
$`N_{\overline{M}}^{(i)}=(0,0,0,0,\theta _i\sqrt{G_{\overline{4}\overline{4}}^{(i)}}),N^{(i)\overline{M}}=(0,0,0,0,\theta _i\sqrt{G^{(i)\overline{4}\overline{4}}}).`$ (160)
To summarize, we have three types of coordinate system in this section: Roman indices denote bulk coordinates, barred Roman indices with superscript <sup>(i)</sup> denote BNCs on the $`i`$-th brane, and Greek indices with superscript <sup>(i)</sup> denote brane coordinates on the $`i`$-th brane. One exception to these rules is $`N_M^{(i)}`$: even though $`N_\alpha ^{(i)}`$ has a Greek index and a superscript <sup>(i)</sup>, it denotes part of a bulk vector.
### 5.1 Derivation of the equations of motion
The most general interval picture action for 5d braneworld gravity, up to second order in derivatives, is
$`S`$ $`=`$ $`\left({\displaystyle __1}d^5x+{\displaystyle __2}d^5x\right)\sqrt{G}\left(2M^3R\mathrm{\Lambda }\right)`$ (161)
$`+2M^3{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}}d^4x^{(i)}\sqrt{g^{(i)}}(\lambda _i\stackrel{~}{}^{(i)}U_i)+4M^3{\displaystyle _{_1+_2}}K.`$
$`R`$ is a Ricci scalar constructed from $`G_{MN}`$, while $`\stackrel{~}{}^{(i)}`$ is a 4d Ricci scalar made of only $`g_{\mu \nu }^{(i)}`$, with $`\stackrel{~}{}`$ indicating that it is a 4d quantity. The brane tensions $`V_i`$ have been rescaled: $`U_i=V_i/2M^3`$, as have the coefficients of the brane kinetic terms: $`\lambda _i=M_i^2/M^3`$.
$`K_{\alpha \beta }`$ is the extrinsic curvature, defined on the boundary hypersurface :
$`K_{\alpha \beta }^{(i)}=_MN_N^{(i)}{\displaystyle \frac{x^M}{x^{(i)\alpha }}}{\displaystyle \frac{x^N}{x^{(i)\beta }}}.`$ (162)
So
$`K^{(i)}=g^{(i)}{}_{}{}^{\alpha \beta }K_{\alpha \beta }^{(i)}=(G^{MN}N^{(i)M}N^{(i)N})_MN_N^{(i)}=G^{MN}_MN_N^{(i)},`$ (163)
where the last equality is because $`N^{(i)M}N_M^{(i)}=1`$, implying $`N^{(i)N}_MN_N^{(i)}=0`$.
The Gibbons-Hawking (GH) extrinsic curvature term in (161) is essential for a gravity analysis in spaces with nontrivial boundary. It ensures that, in the absence of boundary/brane sources, the EOM reduce to the usual Einstein equations for variations of the metric which vanish on the boundary. However in brane setups such as we are considering, the variations of the metric do not vanish on the boundary. As a result, the GH term will make a nontrivial contribution to the boundary part of the EOM.
Letโs find the equations of motion for (161). Replacing $`G_{MN}`$ by $`G_{MN}+\delta G_{MN}`$ and expanding up to first order in $`\delta G_{MN}`$, the first term of (161) gives
$`\sqrt{G}\left(R{\displaystyle \frac{\mathrm{\Lambda }}{2M^3}}\right)\sqrt{G}\left\{\delta R+\delta G^{MN}R_{MN}+\left(R{\displaystyle \frac{\mathrm{\Lambda }}{2M^3}}\right){\displaystyle \frac{\delta G}{2}}\right\},`$ (164)
where
$`\delta G`$ $`=`$ $`G^{MN}\delta G_{MN}=G_{MN}\delta G^{MN},`$ (165)
$`\delta R`$ $`=`$ $`_M(_N\delta G^{MN}G_{PQ}^M\delta G^{PQ}).`$ (166)
The last two terms of (164) give the bulk part of the variation, which contains the Einstein tensor:
$`{\displaystyle \frac{\delta S}{2M^3}}|_{\mathrm{bulk}}={\displaystyle d^5x\sqrt{G}\left\{R_{MN}\frac{G_{MN}}{2}\left(R\frac{\mathrm{\Lambda }}{2M^3}\right)\right\}\delta G^{MN}}.`$ (167)
Next, from the brane part of (161) we get
$`\sqrt{g^{(i)}}\left\{\lambda _i\delta \stackrel{~}{}^{(i)}+\lambda _i\delta g^{(i)\alpha \beta }\stackrel{~}{}_{\alpha \beta }^{(i)}{\displaystyle \frac{g_{\alpha \beta }^{(i)}}{2}}\left(\lambda _i\stackrel{~}{}^{(i)}U_i\right)\delta g^{(i)\alpha \beta }\right\},`$ (168)
where
$`\delta \stackrel{~}{}^{(i)}=\stackrel{~}{}_\alpha ^{(i)}(\stackrel{~}{}_\beta ^{(i)}\delta g^{(i)\alpha \beta }g_{\gamma \delta }^{(i)}\stackrel{~}{}^{(i)\alpha }\delta g^{(i)\gamma \delta }),`$ (169)
with $`\stackrel{~}{}^{(i)}`$ a covariant derivative with respect to $`g_{\alpha \beta }^{(i)}`$. Since our bent branes extend to infinity along $`x^{(i)\mu }`$-directions, we can drop 4d total derivatives, and the brane part of the variation is
$`{\displaystyle \frac{\delta S}{2M^3}}|_{\mathrm{brane}}={\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}}d^4x^{(i)}\sqrt{g^{(i)}}\left(\lambda _i\stackrel{~}{}_{\alpha \beta }^{(i)}{\displaystyle \frac{g_{\alpha \beta }^{(i)}}{2}}(\lambda _i\stackrel{~}{}^{(i)}U_i)\right)\delta g^{(i)\alpha \beta }.`$ (170)
The $`\delta R`$ term in (164) and the last term of (161) produce the boundary part of the variation: applying the Gauss theorem in the curved spacetime, the $`\delta R`$ term gives
$`{\displaystyle \frac{\delta S}{2M^3}}|_{\delta R}={\displaystyle d^5x\sqrt{G}\delta R}`$
$`=2{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\sqrt{g^{(i)}}N_M^{(i)}\left(_N\delta G^{MN}G_{PQ}^M\delta G^{PQ}\right),`$ (171)
where the factor of 2 is because we have used the symmetry discussed in ยง2.1 to write four boundary contributions in terms of two, and
$`{\displaystyle _{\mathrm{\Phi }_1=0}^{(\mathrm{bdy})}}={\displaystyle _{\mathrm{\Phi }_1=0^+}},{\displaystyle _{\mathrm{\Phi }_2=0}^{(\mathrm{bdy})}}={\displaystyle _{\mathrm{\Phi }_2=0^{}}}.`$
The $`K`$-term is a bit more complicated; we get
$`{\displaystyle \frac{\delta S}{2M^3}}|_K=4{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\delta \left(\sqrt{g^{(i)}}G^{MN}_MN_N^{(i)}\right)`$
$`=2{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\sqrt{g^{(i)}}\{2N_M^{(i)}_N\delta G^{MN}(G_{MN}+N_M^{(i)}N_N^{(i)})N^{(i)P}_P\delta G^{MN}`$
$`+(2_MN_N^{(i)}_P(N^{(i)P}N_M^{(i)}N_N^{(i)})(G_{MN}N_M^{(i)}N_N^{(i)})_PN^{(i)P})\delta G^{MN}\}.`$ (172)
Note that we are varying $`g_{\alpha \beta }^{(i)}`$ and $`N_M^{(i)}`$ as well because $`\delta G^{MN}`$ does not vanish on the boundary. Then, combining (5.1) and (5.1) gives
$`{\displaystyle \frac{\delta S}{2M^3}}|_{\mathrm{bdy}}`$ $`=`$ $`2{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\sqrt{g^{(i)}}\{N_M^{(i)}(_N\delta G^{MN}N_N^{(i)}N^{(i)P}_P\delta G^{MN})`$ (173)
$`+(2_MN_N^{(i)}2N^{(i)P}N_M^{(i)}_PN_N^{(i)}G_{MN}_PN^{(i)P})\delta G^{MN}\}.`$
By introducing the projection operator, $`P^{(i)}`$, onto the $`i`$-th hypersurface, defined by
$`P_{MN}^{(i)}G_{MN}N_M^{(i)}N_N^{(i)},`$ (174)
(173) can be further simplified into
$`{\displaystyle \frac{\delta S}{2M^3}}|_{\mathrm{bdy}}`$ $`=`$ $`2{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\sqrt{g^{(i)}}\{P^{(i)}{}_{M}{}^{P}_{P}^{}(N_N^{(i)}\delta G^{MN})`$ (175)
$`+(P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}G_{MN}_PN^{(i)P})\delta G^{MN}\}.`$
Using the identity
$`P_Q^{(i)P}_PP_M^{(i)Q}=N_M^{(i)}_PN^{(i)P},`$ (176)
we can rewrite (175) as follows:
$`{\displaystyle \frac{\delta S}{2M^3}}|_{\mathrm{bdy}}`$ $`=`$ $`2{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\sqrt{g^{(i)}}\{P^{(i)}{}_{Q}{}^{P}_{P}^{}(P_M^{(i)Q}N_N^{(i)}\delta G^{MN})`$ (177)
$`+(P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}P_{MN}^{(i)}_PN^{(i)P})\delta G^{MN}\}.`$
At this point (177) does not seem to give us an EOM because of the first term of the integrand, which contains a derivative of $`\delta G^{MN}`$. However $`P_M^{(i)P}_P`$ is the tangential covariant derivative along the boundary hypersurface and $`P_M^{(i)Q}N_N^{(i)}\delta G^{MN}`$ is a vector tangential to the hypersurface. Thus the first term in (177) is a total tangential divergence, which is equivalent to a 4d total divergence, and can be dropped.
Now the complete variation of the action is
$`{\displaystyle \frac{\delta S}{2M^3}}`$ $`=`$ $`{\displaystyle d^5x\sqrt{G}\left\{R_{MN}\frac{G_{MN}}{2}\left(R\frac{\mathrm{\Lambda }}{2M^3}\right)\right\}\delta G^{MN}}`$
$`+{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}}d^4x^{(i)}\sqrt{g^{(i)}}\left(\lambda _i\stackrel{~}{}_{\alpha \beta }^{(i)}{\displaystyle \frac{g_{\alpha \beta }^{(i)}}{2}}(\lambda _i\stackrel{~}{}^{(i)}U_i)\right)\delta g^{(i)\alpha \beta }`$
$`+{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}^{(\mathrm{bdy})}}d^4x^{(i)}\sqrt{g^{(i)}}(2P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}2P_{MN}^{(i)}_PN^{(i)P})\delta G^{MN}.`$
From the arguments presented in ยง2.1, we can drop the distinction between brane and boundary contributions, obtaining
$`{\displaystyle \frac{\delta S}{2M^3}}`$ $`=`$ $`{\displaystyle d^5x\sqrt{G}\left\{R_{MN}\frac{G_{MN}}{2}\left(R\frac{\mathrm{\Lambda }}{2M^3}\right)\right\}\delta G^{MN}}`$
$`+{\displaystyle \underset{i}{}}{\displaystyle _{\mathrm{\Phi }_i=0}}d^4x^{(i)}\sqrt{g^{(i)}}\{(\lambda _i\stackrel{~}{}_{\alpha \beta }^{(i)}{\displaystyle \frac{g_{\alpha \beta }^{(i)}}{2}}(\lambda _i\stackrel{~}{}^{(i)}U_i))e^{(i)}{}_{M}{}^{\overline{\alpha }}e_{}^{(i)}_N^{\overline{\beta }}`$
$`+2P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}2P_{MN}^{(i)}_PN^{(i)P}\}\delta G^{MN},`$
where
$`e_M^{(i)\overline{M}}={\displaystyle \frac{x^{(i)\overline{M}}}{x^M}}`$ (180)
transforms $`M`$-indices into $`\overline{M}`$-ones. Thus we have the bulk equations of motion:
$`R_{MN}{\displaystyle \frac{G_{MN}}{2}}\left(R{\displaystyle \frac{\mathrm{\Lambda }}{2M^3}}\right)=0,`$ (181)
supplemented by the brane-boundary equations:
$`[(\lambda _i\stackrel{~}{}_{\alpha \beta }^{(i)}{\displaystyle \frac{g_{\alpha \beta }^{(i)}}{2}}(\lambda _i\stackrel{~}{}^{(i)}U_i))e^{(i)}{}_{M}{}^{\overline{\alpha }}e_{}^{(i)}_N^{\overline{\beta }}`$ (182)
$`+2P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}2P_{MN}^{(i)}_PN^{(i)P}]_{\mathrm{\Phi }_i=0}=0.`$
Equations (181-182) are completely general and can be applied to arbitrary boundaries. It is completely covariant under general coordinate transformations, including those that bend the branes.
### 5.2 Straight gauges
It is extremely convenient to work in straight gauges. These are defined as follows:
> straight gauge: a choice of 5d bulk coordinate system such that:
>
> * both of the branes are described by straight slices $`y=y_i`$;
> * $`G_{\mu 4}^{(i)}=[G_{\mu 4}]_{y=y_i}=0`$ for $`i=1,2`$.
From (157), straight slices at $`y=y_i`$ implies that
$`N_\mu ^{(i)}=0,N_4^{(i)}={\displaystyle \frac{\theta _i}{\sqrt{G^{(i)44}}}}.`$ (183)
Then using
$`e_{\overline{\alpha }}^{(i)M}N_M^{(i)}=e_{\overline{\alpha }}^{(i)\overline{M}}N_{\overline{M}}^{(i)}=N_{\overline{\alpha }}^{(i)}=0,`$ (184)
and (183), we get
$`e_{\overline{\alpha }}^{(i)4}=0.`$ (185)
Note however that $`G_{44}^{(i)}=1/G^{(i)44}`$ is still arbitrary.
Thus we see that an equivalent (and perhaps more intuitive) definition of straight gauges is:
> straight gauge: a choice of 5d bulk coordinate system such that $`N^{(i)\mu }=0`$ and $`N_\mu ^{(i)}=0`$.
Yet another equivalent definition is that a straight gauge is any choice of 5d bulk coordinates such that the bulk coordinates are BNCโs at the locations of both branes.
A natural question is whether it is always possible to impose a straight gauge, starting from an arbitrary bulk coordinate system. We can prove this, without loss of generality, by starting from a bulk coordinate system where the first brane is at $`y=0`$ with $`N_\mu ^{(1)}=0`$ and $`[G_{\mu 4}]_{y=0}=0`$, while the second brane is bent:
$`\mathrm{\Phi }_1=y,\mathrm{\Phi }_2=yL\rho (x),`$ (186)
and $`N_\mu ^{(2)}`$, $`[G_{\mu 4}]_{\mathrm{\Phi }_2=0}`$ do not necessarily vanish.
To get to a straight gauge, we first perform a GCT defined by
$`\stackrel{ห}{x}^\mu =x^\mu ,\stackrel{ห}{y}=y{\displaystyle \frac{\rho (x)}{L+\rho (x)}}y,`$ (187)
under which
$`\stackrel{ห}{\mathrm{\Phi }}_1=\stackrel{ห}{y},\stackrel{ห}{\mathrm{\Phi }}_2=\stackrel{ห}{y}L,`$ (188)
but
$`[\stackrel{ห}{G}_{\mu 4}]_{\stackrel{ห}{y}=y_i}`$ $`=`$ $`\left[{\displaystyle \frac{x^M}{\stackrel{ห}{x}^\mu }}{\displaystyle \frac{x^N}{\stackrel{ห}{y}}}G_{MN}\right]_{\stackrel{ห}{y}=y_i}=\left[{\displaystyle \frac{y}{\stackrel{ห}{y}}}\left(G_{\mu 4}+{\displaystyle \frac{y}{\stackrel{ห}{x}^\mu }}G_{44}\right)\right]_{\stackrel{ห}{y}=y_i}`$ (189)
$`=`$ $`{\displaystyle \frac{L+\rho (x)}{L}}\left[G_{\mu 4}+{\displaystyle \frac{\stackrel{ห}{y}}{L}}{\displaystyle \frac{\rho }{\stackrel{ห}{x}^\mu }}G_{44}\right]_{\stackrel{ห}{y}=y_i}.`$
That is, the first condition in our definition of a straight gauge is satisfied but $`[\stackrel{ห}{G}_{\mu 4}]_{\stackrel{ห}{y}=L}`$ is still non-vanishing. Now we perform a second GCT such that
$`\widehat{y}=\stackrel{ห}{y},\widehat{x}^\mu =f^\mu (\stackrel{ห}{x},\stackrel{ห}{y});`$ (190)
both branes are still described by $`\widehat{y}=y_i`$ and
$`[\widehat{G}_{\mu 4}]_{\stackrel{ห}{y}=y_i}`$ $`=`$ $`\left[{\displaystyle \frac{\stackrel{ห}{x}^\alpha }{\widehat{x}^\mu }}\left({\displaystyle \frac{\stackrel{ห}{x}^\beta }{\widehat{y}}}\stackrel{ห}{g}_{\alpha \beta }+\stackrel{ห}{G}_{\alpha 4}\right)\right]_{\widehat{y}=y_i}.`$ (191)
(191) does not necessarily vanish at $`\widehat{y}=0,L`$ for arbitrary $`\stackrel{ห}{G}_{MN}`$. But for any fixed $`\stackrel{ห}{G}_{MN}`$, the quantity inside the parentheses can be set to be zero by choosing, for example,
$`\widehat{x}^\alpha =\stackrel{ห}{x}^\alpha +{\displaystyle ๐\widehat{y}\stackrel{ห}{g}^{\alpha \beta }\stackrel{ห}{G}_{\beta 4}}.`$ (192)
Therefore it is always possible to find a bulk coordinate system satisfying straight gauge conditions.
The general brane-boundary equations (182) simplify quite a bit in a straight gauge. To see this, contract the tensor equations (182) with $`e_{\overline{M}}^{(i)M}e_{\overline{N}}^{(i)N}`$:
$`[(\lambda _i\stackrel{~}{}_{\alpha \beta }^{(i)}{\displaystyle \frac{g_{\alpha \beta }^{(i)}}{2}}(\lambda _i\stackrel{~}{}^{(i)}U_i))\delta _{\overline{M}}^\alpha \delta _{\overline{N}}^\beta `$
$`+e_{\overline{M}}^{(i)M}e_{\overline{N}}^{(i)N}(P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}+P^{(i)}{}_{N}{}^{P}_{P}^{}N_M^{(i)}2P_{MN}^{(i)}_PN^{(i)P})]_{\mathrm{\Phi }_i=0}=0.`$ (193)
These brane-boundary equations break up into three tensor equations each. The $`\overline{4}\overline{4}`$ equation is:
$`\left[e_{\overline{4}}^{(i)M}e_{\overline{4}}^{(i)N}(P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}+P^{(i)}{}_{N}{}^{P}_{P}^{}N_M^{(i)}2P_{MN}^{(i)}_PN^{(i)P})\right]_{\mathrm{\Phi }_i=0}=0.`$ (194)
This is trivially satisfied, since $`e_{\overline{4}}^{(i)M}`$ is parallel to $`N^{(i)M}`$,<sup>1</sup><sup>1</sup>1For any 5-vector $`T_M`$ tangential to $`\mathrm{\Phi }_i=0`$-hypersurface, $`e_{\overline{4}}^{(i)M}T_M=e_{\overline{4}}^{(i)\overline{M}}T_{\overline{M}}=T_{\overline{4}}=0.`$ and $`N^{(i)M}`$ contracted with a projection operator $`P_M^{(i)P}`$ vanishes.
The $`\overline{\mu }\overline{4}`$ part is:
$`\left[e_{\overline{\mu }}^{(i)M}e_{\overline{4}}^{(i)N}(P^{(i)}{}_{M}{}^{P}_{P}^{}N_N^{(i)}+P^{(i)}{}_{N}{}^{P}_{P}^{}N_M^{(i)}2P_{MN}^{(i)}_PN^{(i)P})\right]_{\mathrm{\Phi }_i=0}=0.`$ (195)
The second and third terms vanish for the same reason as above, leaving only the first term, which is proportional to $`N^{(i)N}_PN_N^{(i)}=0`$.
So only the $`\overline{\mu }\overline{\nu }`$ brane-boundary equation has any content. It can be simplified using (184):
$`[(\lambda _i\stackrel{~}{}_{\overline{\alpha }\overline{\beta }}^{(i)}{\displaystyle \frac{g_{\overline{\alpha }\overline{\beta }}^{(i)}}{2}}(\lambda _i\stackrel{~}{}^{(i)}U_i))`$
$`+e_{\overline{\alpha }}^{(i)M}e_{\overline{\beta }}^{(i)N}(_MN_N^{(i)}+_NN_M^{(i)})2g_{\overline{\alpha }\overline{\beta }}_PN^{(i)P}]_{\mathrm{\Phi }_i=0}=0.`$ (196)
Now we impose a straight gauge. Then because of (185), we can always choose the BNCs such that
$`e_{\overline{\alpha }}^{(i)M}=\delta _{\overline{\alpha }}^M.`$ (197)
Furthermore, in the straight gauge the $`y`$-direction of the bulk coordinate system on the $`i`$-th brane is parallel to $`N_M^{(i)}`$ which is in the $`\overline{y}^{(i)}`$-direction, and thus we can take
$`e_{\overline{4}}^{(i)M}=\delta _{\overline{4}}^M.`$ (198)
That is, $`e^{(i)}=1`$ and we need not distinguish between $`x^{(i)\overline{M}}`$-system and $`[x^M]_{\mathrm{\Phi }_i=0}`$-one; one bulk coordinate patch can describe the whole spacetime including the boundary while keeping a straight gauge, which justifies dropping bars on indices in (5.2).
Due to (183) and the second condition in our definition of a straight gauge, we get
$`_\alpha N_\beta ^{(i)}+_\beta N_\alpha ^{(i)}=2\mathrm{\Gamma }_{\alpha \beta }^4N_4^{(i)}=\theta _i\sqrt{G^{44}}g_{\alpha \beta }^{}.`$ (199)
Similarly:
$`2g_{\alpha \beta }_PN^{(i)P}=2\theta _ig_{\alpha \beta }(\sqrt{G^{44}}^{}+\sqrt{G^{44}}\mathrm{\Gamma }_{P4}^P)=\theta _i\sqrt{G^{44}}g_{\alpha \beta }g^{\rho \sigma }g_{\rho \sigma }^{}.`$ (200)
Putting together (5.2-200) we get the full EOM in an arbitrary straight gauge:
$`\mathrm{bulk}:R_{MN}{\displaystyle \frac{1}{2}}G_{MN}(R{\displaystyle \frac{\mathrm{\Lambda }}{2M^3}})=0,`$ (201)
$`\mathrm{brane}\mathrm{boundary}:\left[\lambda _i\stackrel{~}{}_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }(\lambda _i\stackrel{~}{}U_i)+\theta _i\sqrt{G^{44}}(g_{\mu \nu }^{}g_{\mu \nu }g_{\rho \sigma }^{}g^{\rho \sigma })\right]_{y=y_i}=0.`$ (202)
Recall that $`\theta _1=\theta _2=1`$, and that strictly speaking the terms multiplying $`\theta _i`$ are evaluated at $`y=0^+`$, $`L^{}`$, not at $`y=0`$, $`L`$.
### 5.3 Background solutions
For a linearized analysis, we write
$`G_{MN}=G_{MN}^\mathrm{๐}+h_{MN}=\left(\begin{array}{cc}g_{\mu \nu }^\mathrm{๐}& 0\\ 0& 1\end{array}\right)+\left(\begin{array}{cc}h_{\mu \nu }& h_{\mu 4}\\ h_{4\nu }& h_{44}\end{array}\right).`$ (203)
We can solve (201-202) with a straight gauge ansatz for a general warped $`AdS_4`$ background metric:
$`g_{\mu \nu }^\mathrm{๐}={\displaystyle \frac{a(y)^2}{(1\frac{H^2x^2}{4})^2}}\eta _{\mu \nu },`$ (204)
with $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ and $`x^2=\eta _{\mu \nu }x^\mu x^\nu `$. This corresponds to a warped geometry where each slice is $`AdS_4`$ (or 4d Minkowski space in the limit $`H^20`$). <sup>2</sup><sup>2</sup>2It is also possible to obtain factorizable backgrounds with $`dS_4`$ slices, but we will not consider these solutions here. See .
It is easy to show that
$`\mathrm{\Gamma }^\mathrm{๐}_{\nu \lambda }^\mu `$ $`=`$ $`{\displaystyle \frac{H^2}{2(1\frac{H^2x^2}{4})}}(\delta _\lambda ^\mu \eta _{\nu \sigma }x^\sigma +\delta _\nu ^\mu \eta _{\lambda \sigma }x^\sigma \eta _{\nu \lambda }x^\mu ),`$ (205)
$`\mathrm{\Gamma }^\mathrm{๐}_{\nu 4}^\mu `$ $`=`$ $`{\displaystyle \frac{a^{}}{a}}\delta _\nu ^\mu ,`$ (206)
$`\mathrm{\Gamma }^\mathrm{๐}_{\mu \nu }^4`$ $`=`$ $`{\displaystyle \frac{a^{}}{a}}g_{\mu \nu }^\mathrm{๐},`$ (207)
while all the other components of $`\mathrm{\Gamma }^\mathrm{๐}_{NP}^M`$ vanish.
Also we find
$`R_{\mu \nu }^\mathrm{๐}`$ $`=`$ $`{\displaystyle \frac{3H^2+3a^2+aa^{\prime \prime }}{a^2}}g_{\mu \nu }^\mathrm{๐},`$ (208)
$`R_{44}^\mathrm{๐}`$ $`=`$ $`{\displaystyle \frac{4a^{\prime \prime }}{a}},`$ (209)
$`R^\mathrm{๐}`$ $`=`$ $`{\displaystyle \frac{4(3H^2+3a^2+2aa^{\prime \prime })}{a^2}},`$ (210)
$`\stackrel{~}{}_{\mu \nu }^\mathrm{๐}|_{y=y_i}`$ $`=`$ $`\left[{\displaystyle \frac{3H^2}{a^2}}g_{\mu \nu }^\mathrm{๐}\right]_{y=y_i},`$ (211)
$`\stackrel{~}{}^\mathrm{๐}|_{y=y_i}`$ $`=`$ $`\left[{\displaystyle \frac{12H^2}{a^2}}\right]_{y=y_i},`$ (212)
and $`R_{\mu 4}^\mathrm{๐}=0`$. Then, with $`G_{MN}`$ replaced by $`G_{MN}^\mathrm{๐}`$, (201) gives
$`\left(H^2+a_{}^{}{}_{}{}^{2}2k^2a^2+aa^{\prime \prime }\right)g_{\mu \nu }^\mathrm{๐}`$ $`=`$ $`0,`$ (213)
$`H^2+a_{}^{}{}_{}{}^{2}k^2a^2`$ $`=`$ $`0,`$ (214)
where $`k^2=\mathrm{\Lambda }/24M^3`$. We will restrict our consideration to models with a negative bulk cosmological constant. From (202) we get
$`\left[\left({\displaystyle \frac{U_i}{6}}+{\displaystyle \frac{\lambda _iH^2}{a^2}}\theta _i{\displaystyle \frac{2a^{}}{a}}\right)g_{\mu \nu }^\mathrm{๐}\right]_{y=y_i}=0.`$ (215)
The general solution of (214) with normalization $`a(0)=1`$ has the form:
$`a(y)={\displaystyle \frac{\mathrm{cosh}k(yy_0)}{\mathrm{cosh}ky_0}},`$ $`0<y<L,`$ (216)
where
$`\mathrm{cosh}ky_0={\displaystyle \frac{k}{H}}.`$ (217)
With this solution, (213) is automatically satisfied. (215) gives boundary conditions at $`y=0`$ and $`L`$:
$`y=0`$ $`:`$ $`2kT_0={\displaystyle \frac{U_0}{6}}+{\displaystyle \frac{\lambda _0H^2}{a(0)^2}},`$ (218)
$`y=L`$ $`:`$ $`2kT_L={\displaystyle \frac{U_L}{6}}+{\displaystyle \frac{\lambda _LH^2}{a(L)^2}},`$ (219)
where $`T_0=\mathrm{tanh}ky_0`$ and $`T_L=\mathrm{tanh}k(Ly_0)`$.
For convenience we define $`v_i=k\lambda _i=kM_i^2/M^3`$ and $`w_i=U_i/k=V_i/(2kM^3)`$ and solve (218) and (219) for $`T_0`$ and $`T_L`$ respectively to get
$`T_i`$ $`=`$ $`{\displaystyle \frac{U_i}{12k}}+{\displaystyle \frac{\lambda _i}{2k}}{\displaystyle \frac{H^2\mathrm{cosh}^2ky_0}{\mathrm{cosh}^2k(y_iy_0)}}={\displaystyle \frac{w_i}{12}}+{\displaystyle \frac{v_i}{2}}(1T_i^2)`$ (220)
$`T_i^\pm ={\displaystyle \frac{1}{v_i}}\left(1\pm \sqrt{1+{\displaystyle \frac{1}{6}}w_iv_i+v_i^2}\right).`$
Given any input values for the brane tensions $`V_i`$ and brane Planck constants $`M_i`$, we can find a background solution by solving for the 4d curvature parameter $`H`$ and the brane separation $`L`$. Equivalently, we can specify $`w_0`$, $`w_L`$, $`v_0`$ and $`v_L`$ as inputs and solve for $`T_0`$ and $`T_L`$ using (220). For example, if $`w_0=w_L=12`$, $`v_0>1`$ and $`v_L<1`$, then $`H=0`$ (i.e. the branes are flat), the value of $`L`$ is undetermined, and $`T_0=T_L=1`$. This special case becomes the original Randall-Sundrum model when we take $`v_0,v_L0`$.
Recall that we are only considering the case where the 4d curvature is $`AdS`$-like, i.e. the bulk space is approximately $`AdS_5/AdS_4`$. This means that $`H^2>0`$, and the $`T_i`$ are real and satisfy $`|T_i|<1`$. Choices of input parameters which do not satisfy these conditions do not give $`AdS_5/AdS_4`$ solutions. Solving $`1<T_i^+<1`$, we get
$`(v_i0w_i6v_i{\displaystyle \frac{6}{v_i}}`$
$`((v_i1w_i12)(v_i<112w_i12)))`$
$`(v_i<0w_i6v_i{\displaystyle \frac{6}{v_i}}`$
$`((v_i<1w_i12)(v_i112w_i12))).`$ (221)
The results for $`T_i^{}`$ are similar. Note that there are solutions for both positive and negative brane tensions, and for both positive and negative brane Planck constants.
### 5.4 Gauge fixing
Having determined the general background solution, we have to deal with the metric fluctuations, $`h_{MN}`$, as given in (203). We will perform a complete gauge-fixing, starting with the straight gauge implied by the background solution. All indices will be raised and lowered using the background metric $`G_{MN}^\mathrm{๐}`$, but to reduce clutter we will omit the superscript <sup>0</sup> on $`g_{\mu \nu }`$.
Under a linearized 5d general coordinate transformation $`x^Mx^M+\xi ^M`$ the metric fluctuations transform as follows:
$`h_{\mu \nu }`$ $``$ $`h_{\mu \nu }g_{\mu \nu }{\displaystyle \frac{2a^{}}{a}}\xi ^4\stackrel{~}{}_\mu \xi _\nu \stackrel{~}{}_\nu \xi _\mu ,`$ (222)
$`h_{\mu 4}`$ $``$ $`h_{\mu 4}g_{\mu \nu }\xi ^\nu {}_{}{}^{}_\mu \xi ^4,`$ (223)
$`h_{44}`$ $``$ $`h_{44}2\xi ^4{}_{}{}^{}.`$ (224)
We start with a partial gauge-fixing to exhibit the radion, letting $`\xi ^{(\mathrm{I})}{}_{}{}^{\mu }=0`$ and
$`\xi ^{(\mathrm{I})}{}_{}{}^{4}={\displaystyle \frac{1}{2}}{\displaystyle ^y}h_{44}dy{\displaystyle \frac{1}{2}}{\displaystyle ^y}F(y)\psi (x)dy,`$ (225)
with $`F(y)`$ a fixed but arbitrary function of $`y`$. This transforms an arbitrary $`h_{44}`$ into
$`h_{44}=F(y)\psi (x).`$ (226)
Since we want to be in a straight gauge, we must require that $`\xi ^{(\mathrm{I})}^4`$ vanishes at the locations of the branes. On the interval $`0<y<L`$, this fixes the $`y`$-independent part of (225):
$`\xi ^{(\mathrm{I})}{}_{}{}^{4}={\displaystyle \frac{1}{2}}{\displaystyle _0^y}h_{44}dy{\displaystyle \frac{1}{2}}{\displaystyle _0^y}F(y)\psi (x)dy,`$ (227)
and fixes a relation between the radion field $`\psi (x)`$, $`F(y)`$ and the original metric fluctuation $`h_{44}(x,y)`$:
$`\psi (x)={\displaystyle \frac{_0^Lh_{44}๐y}{_0^LF(y)๐y}}.`$ (228)
From (228) we see that $`F(y)`$, though arbitrary, must be nonzero. More precisely, the straight gauge condition requires:
$`{\displaystyle _0^L}F(y)๐y0.`$ (229)
Note that for a general metric fluctuation $`h_{\mu 4}(x,y)`$, we are not yet in a straight gauge since $`G_{\mu 4}^{(i)}0`$. So our next step is to fix to a straight gauge, by a partial gauge-fixing which eliminates $`h_{\mu 4}(x,y)`$ altogether. Choose $`\xi ^{(\mathrm{II})}{}_{}{}^{4}=0`$ and
$`\xi ^{(\mathrm{II})}{}_{}{}^{\mu }={\displaystyle ^y}h^{\mu 4}dy.`$ (230)
Then $`h_{\mu \nu }`$ is still arbitrary, $`h_{44}`$ is unaffected, and
$`h_{\mu 4}=0.`$ (231)
Given the straight gauge conditions and the gauge choices (226) and (231), the residual gauge freedom is generated by
$`\xi ^4=0,\xi ^\mu =\xi ^\mu (x).`$ (232)
Note that what actually appears in the general coordinate transformation for $`h_{\mu \nu }`$ is $`\stackrel{~}{}_\mu \xi _\nu (x)+\stackrel{~}{}_\nu \xi _\mu (x)`$, which picks up a nontrivial $`y`$ dependence, $`a^2(y)`$, from lowering the vector index.
The general coordinate transformation generated by
$`\xi ^4=\xi ^4(x)ฯต(x),\xi ^\mu ={\displaystyle \frac{a^2}{H^2}}{\displaystyle \frac{a^{}}{a}}\stackrel{~}{}^\mu ฯต(x),`$ (233)
respects (226) and (231) but takes us out of the straight gauge. The scalar $`ฯต(x)`$ is the putative brane-bending mode. Since the equations of motion are covariant, even under a brane-bending transformation generated by $`ฯต(x)`$, this mode is pure gauge.
The full linearized bulk equations of motion are given by:
$`\mu \nu \mathrm{part}`$ $`:`$ $`_P_\mu h_\nu ^P+_P_\nu h_\mu ^P^2h_{\mu \nu }_\mu _\nu h`$ (234)
$`g_{\mu \nu }(_M_Nh^{MN}^2h)4k^2g_{\mu \nu }h+8k^2h_{\mu \nu }=0,`$
$`\mu 4\mathrm{part}`$ $`:`$ $`_P_\mu h_4^P+_P_4h_\mu ^P^2h_{\mu 4}_\mu _4h=0,`$ (235)
$`44\mathrm{part}`$ $`:`$ $`2_P_4h_4^P^2h_{44}_4_4h`$ (236)
$`_M_Nh^{MN}+^2h4k^2h+8k^2h_{44}=0,`$
where $`h=G^{MN}h_{MN}`$. In our background the above EOM can be expanded using the following identities, which hold for any 5-vector $`T^M`$:
$`_\mu T^\nu `$ $`=`$ $`\stackrel{~}{}_\mu T^\nu +{\displaystyle \frac{a^{}}{a}}\delta _\mu ^\nu T^4,`$ (237)
$`_\mu T^4`$ $`=`$ $`\stackrel{~}{}_\mu T^4{\displaystyle \frac{a^{}}{a}}T_\mu .`$ (238)
Using these and our partial gauge-fixings, (226) and (231), we obtain
$`0`$ $`=`$ $`\stackrel{~}{}_\rho \stackrel{~}{}_\mu h_\nu ^\rho +\stackrel{~}{}_\rho \stackrel{~}{}_\nu h_\mu ^\rho \stackrel{~}{}^2h_{\mu \nu }\stackrel{~}{}_\mu \stackrel{~}{}_\nu \stackrel{~}{h}g_{\mu \nu }(\stackrel{~}{}_\rho \stackrel{~}{}_\sigma h^{\rho \sigma }\stackrel{~}{}^2\stackrel{~}{h})`$ (239)
$`h_{\mu \nu }^{\prime \prime }+g_{\mu \nu }\stackrel{~}{h}^{\prime \prime }+{\displaystyle \frac{4a^{}}{a}}g_{\mu \nu }\stackrel{~}{h}^{}+{\displaystyle \frac{8H^2+4a^{}^2}{a^2}}h_{\mu \nu }{\displaystyle \frac{3H^2}{a^2}}g_{\mu \nu }\stackrel{~}{h}`$
$`F\stackrel{~}{}_\mu \stackrel{~}{}_\nu \psi +g_{\mu \nu }F\stackrel{~}{}^2\psi {\displaystyle \frac{3a^{}}{a}}g_{\mu \nu }F^{}\psi {\displaystyle \frac{6H^2+12a^{}^2}{a^2}}g_{\mu \nu }F\psi ,`$
$`0`$ $`=`$ $`(\stackrel{~}{}_\nu h_\mu ^\nu )^{}_\mu \stackrel{~}{h}^{}+{\displaystyle \frac{3a^{}}{a}}F_\mu \psi ,`$ (240)
$`0`$ $`=`$ $`\stackrel{~}{}_\mu \stackrel{~}{}_\nu h^{\mu \nu }+\stackrel{~}{}^2\stackrel{~}{h}+{\displaystyle \frac{3a^{}}{a}}\stackrel{~}{h}^{}{\displaystyle \frac{3H^2}{a^2}}\stackrel{~}{h}{\displaystyle \frac{12a^{}^2}{a^2}}F\psi ,`$ (241)
with $`\stackrel{~}{h}=g^{\mu \nu }h_{\mu \nu }`$. Also twice (241) subtracted from the trace of (239) gives the auxiliary EOM:
$`0={\displaystyle \frac{(a^2\stackrel{~}{h}^{})^{}}{a^2}}+F\stackrel{~}{}^2\psi {\displaystyle \frac{4a^{}}{a}}F^{}\psi 8k^2F\psi .`$ (242)
By a similar procedure the brane-boundary equations become
$`0`$ $`=`$ $`[\theta _i(h_{\mu \nu }{}_{}{}^{}g_{\mu \nu }\stackrel{~}{h}^{})+({\displaystyle \frac{3\lambda _iH^2}{a^2}}2kT_i)h_{\mu \nu }{\displaystyle \frac{3\lambda _iH^2}{2a^2}}g_{\mu \nu }\stackrel{~}{h}+3kT_ig_{\mu \nu }F\psi `$ (243)
$`+{\displaystyle \frac{\lambda _i}{2}}(\stackrel{~}{}_\rho \stackrel{~}{}_\mu h_\nu ^\rho +\stackrel{~}{}_\rho \stackrel{~}{}_\nu h_\mu ^\rho \stackrel{~}{}^2h_{\mu \nu }\stackrel{~}{}_\mu \stackrel{~}{}_\nu \stackrel{~}{h})`$
$`{\displaystyle \frac{\lambda _i}{2}}g_{\mu \nu }(\stackrel{~}{}_\rho \stackrel{~}{}_\sigma h^{\rho \sigma }\stackrel{~}{}^2\stackrel{~}{h})]_{y=y_i}.`$
#### 5.4.1 โmassiveโ case
We can generalize (305) of Appendix (and shuffle $`\varphi _1`$ and $`\varphi _2`$) to get
$`h_{\mu \nu }=b_{\mu \nu }+\stackrel{~}{}_\mu V_\nu +\stackrel{~}{}_\nu V_\mu +a^2\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }\stackrel{~}{}^2\right)\varphi _1+g_{\mu \nu }\varphi _2,`$ (244)
with
$`\stackrel{~}{}^\mu b_{\mu \nu }=0,\stackrel{~}{b}=0,`$ (245)
$`\stackrel{~}{}^\mu V_\mu =0.`$ (246)
After $`y`$-integration, (242) gives the first equation for $`\varphi _1`$ and $`\varphi _2`$:
$`\stackrel{~}{h}^{}=4\varphi _2^{}={\displaystyle \frac{f_1(x)}{a^2}}๐_4\psi (x)+{\displaystyle \frac{4a^{}}{a}}F\psi (x),`$ (247)
where $`๐_n=\stackrel{~}{}^2\frac{nH^2}{a^2}`$, $`^{}(y)=F(y)`$ and a new field, $`f_1(x)`$, is introduced as an integration โconstantโ for $`\stackrel{~}{h}^{}`$. Of course, there should have been other generic integration constants arising from integrating $`F(y)`$ and $`a(y)`$. But all of them can be absorbed into $``$ and $`f_1`$.
Then (241) and (240) become
$`{\displaystyle \frac{a^2}{4}}๐_4\stackrel{~}{}^2\varphi _1=๐_4\left(\varphi _2{\displaystyle \frac{a^{}}{a}}\psi \right)+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{f_1}{a^2}},`$ (248)
$`(๐_3V_\mu )^{}=3\stackrel{~}{}_\mu \left(\varphi _2^{}{\displaystyle \frac{a^{}}{a}}F\psi {\displaystyle \frac{a^2}{4}}๐_4\varphi _1^{}\right).`$ (249)
As in the flat or RSI case, (244) breaks down when $`\stackrel{~}{}_\mu V_\nu +\stackrel{~}{}_\nu V_\mu `$ and $`\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu \frac{1}{4}g_{\mu \nu }\stackrel{~}{}^2\right)\varphi _1`$ become transverse, i.e. when
$`\stackrel{~}{}^\nu (\stackrel{~}{}_\mu V_\nu +\stackrel{~}{}_\nu V_\mu )=๐_3V_\mu =0,`$ (250)
$`\stackrel{~}{}^\nu \left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }\stackrel{~}{}^2\right)\varphi _1={\displaystyle \frac{3}{4}}\stackrel{~}{}_\mu ๐_4\varphi _1=0.`$ (251)
So our โmassiveโ tensor decomposition (244) is valid for modes such that:
1. scalar modes are not annihilated by $`๐_4`$, and
2. vector modes are not annihilated by $`๐_3`$.
Due to the condition i, we can safely rewrite $`f_1`$ as
$`f_1(x)=a^2๐_4\sigma (x).`$ (252)
We remove $`๐_4`$ from (248) to get
$`{\displaystyle \frac{a^2}{4}}\stackrel{~}{}^2\varphi _1=\varphi _2+{\displaystyle \frac{a^{}}{a}}(\sigma \psi ).`$ (253)
Taking a $`y`$-derivative of it and using (247),
$`\stackrel{~}{}^2(a^2\varphi _1^{}\sigma +\psi )=0.`$ (254)
In $`AdS_4`$, the eigenvalue of $`\stackrel{~}{}^2`$ acting on a scalar is bounded below by $`4H^2/a^2`$, that is, $`\stackrel{~}{}^2`$ cannot kill a scalar . Then (254) gives
$`a^2\varphi _1^{}=\sigma \psi .`$ (255)
Plugging this and (247) into (249), we get
$`(๐_3V_\mu )^{}=0.`$ (256)
Noting condition ii, we see that the $`y`$-dependence of $`V_\mu `$ should be $`a^2`$. This allows us to eliminate $`V_\mu `$ using the transverse part of the residual gauge freedom.
Now (239) boils down to
$`0`$ $`=`$ $`\stackrel{~}{}^2b_{\mu \nu }b_{\mu \nu }^{\prime \prime }+{\displaystyle \frac{4a^2}{a^2}}b_{\mu \nu }`$ (257)
$`+(g_{\mu \nu }\stackrel{~}{}^2\stackrel{~}{}_\mu \stackrel{~}{}_\nu )(a^2\varphi _1^{\prime \prime }+4aa^{}\varphi _1^{}{\displaystyle \frac{a^2}{2}}\stackrel{~}{}^2\varphi _1+2\varphi _2+F\psi )`$
$`+g_{\mu \nu }({\displaystyle \frac{3}{4}}a^2\stackrel{~}{}^2\varphi _1^{\prime \prime }3aa^{}\stackrel{~}{}^2\varphi _1^{}+{\displaystyle \frac{3H^2}{2}}\stackrel{~}{}^2\varphi _1`$
$`+3\varphi _2^{\prime \prime }+{\displaystyle \frac{12a^{}}{a}}\varphi _2^{}{\displaystyle \frac{6H^2}{a^2}}\varphi _2{\displaystyle \frac{3a^{}}{a}}F^{}\psi {\displaystyle \frac{6H^2+12a^{}^2}{a^2}}F\psi )`$
$`=`$ $`\stackrel{~}{}^2b_{\mu \nu }b_{\mu \nu }^{\prime \prime }+{\displaystyle \frac{4a^2}{a^2}}b_{\mu \nu },`$
and (243) becomes
$`0`$ $`=`$ $`[\theta _ib_{\mu \nu }^{}({\displaystyle \frac{\lambda _iH^2}{a^2}}+2kT_i)b_{\mu \nu }{\displaystyle \frac{\lambda _i}{2}}\stackrel{~}{}^2b_{\mu \nu }`$ (258)
$`+(g_{\mu \nu }\stackrel{~}{}^2\stackrel{~}{}_\mu \stackrel{~}{}_\nu )(\theta _ia^2\varphi _1^{}{\displaystyle \frac{\lambda _i}{4}}a^2\stackrel{~}{}^2\varphi _1+\lambda _i\varphi _2)`$
$`+g_{\mu \nu }({\displaystyle \frac{3\theta _i}{4}}a^2\stackrel{~}{}^2\varphi _1^{}+{\displaystyle \frac{3\lambda _iH^2}{4}}\stackrel{~}{}^2\varphi _1{\displaystyle \frac{3\lambda _iH^2}{a^2}}\varphi _23\theta _i\varphi _2^{}+3kT_iF\psi )]_{y=y_i}`$
$`=`$ $`[\theta _ib_{\mu \nu }^{}({\displaystyle \frac{\lambda _iH^2}{a^2}}+2kT_i)b_{\mu \nu }{\displaystyle \frac{\lambda _i}{2}}\stackrel{~}{}^2b_{\mu \nu }`$
$`+\theta _i(1+k\lambda _iT_i)(\stackrel{~}{}_\mu \stackrel{~}{}_\nu g_{\mu \nu }\stackrel{~}{}^2+{\displaystyle \frac{3H^2}{a^2}}g_{\mu \nu })(\sigma \psi )]_{y=y_i}.`$
The trace of (258) is
$`๐_4\left(\sigma (0)\psi \right)=0,๐_4\left(\sigma (L)\psi \right)=0,`$ (259)
from which, considering condition i and $`(0)(L)`$, it follows that
$`\sigma =0,\psi =0.`$ (260)
Then, from (247), (253) and (255):
$`\varphi _2^{}=0,{\displaystyle \frac{a^2}{4}}\stackrel{~}{}^2\varphi _1=\varphi _2,\varphi _1=f_2(x),`$ (261)
and $`h_{\mu \nu }`$ is
$`h_{\mu \nu }=b_{\mu \nu }+a^2\stackrel{~}{}_\mu \stackrel{~}{}_\nu f_2.`$ (262)
Since $`a^2\stackrel{~}{}_\mu \stackrel{~}{}_\nu f_2`$ has the correct $`y`$-dependence and form, it is removed by the longitudinal component of the residual gauge freedom, leaving only $`b_{\mu \nu }`$.
To get the spectrum of $`b_{\mu \nu }`$, first we solve (257). Using the EOM for a transverse-traceless spin-2 field of mass $`m0`$ in an $`AdS_4`$ background :
$`\stackrel{~}{}^2b_{\mu \nu }+{\displaystyle \frac{2H^2m^2}{a^2}}b_{\mu \nu }=0,`$ (263)
and substituting $`z=\mathrm{tanh}k(yy_0)`$, it becomes
$`(1z^2){\displaystyle \frac{d^2b_{\mu \nu }}{dz^2}}2z{\displaystyle \frac{db_{\mu \nu }}{dz}}+\left(2+{\displaystyle \frac{m^2}{H^2}}{\displaystyle \frac{4}{1z^2}}\right)b_{\mu \nu }=0,`$ (264)
and its solution is
$`b_{\mu \nu }=A_{\mu \nu }P(l,2,z)+B_{\mu \nu }Q(l,2,z),`$ (265)
where $`P`$ and $`Q`$ are associated Legendre functions of the 1st and 2nd kind respectively and $`l=\frac{1}{2}(1+\sqrt{9+\frac{4m^2}{H^2}})`$.
With (260), (258) gives the boundary conditions:
$`\left[2k(1z^2){\displaystyle \frac{db_{\mu \nu }}{dz}}+\left(4kT_0+\lambda _0k^2{\displaystyle \frac{m^2}{H^2}}(1z^2)\right)b_{\mu \nu }\right]_{y=0}`$ $`=`$ $`0,`$ (266)
$`\left[2k(1z^2){\displaystyle \frac{db_{\mu \nu }}{dz}}+\left(4kT_L+\lambda _Lk^2{\displaystyle \frac{m^2}{H^2}}(1z^2)\right)b_{\mu \nu }\right]_{y=L}`$ $`=`$ $`0.`$ (267)
Plugging (265) into (266) and (267), we get
$`\left(\begin{array}{cc}a_0& b_0\\ a_L& b_L\end{array}\right)\left(\begin{array}{c}A_{\mu \nu }\\ B_{\mu \nu }\end{array}\right)=\left(\begin{array}{c}0\\ 0\end{array}\right),`$ (268)
where
$`a_0`$ $`=`$ $`(k\lambda _0q(1T_0^2)+(3+\sqrt{9+4q})T_0)P({\displaystyle \frac{1}{2}}(1+\sqrt{9+4q}),2,T_0)`$
$`+(3+\sqrt{9+4q})P({\displaystyle \frac{1}{2}}(3+\sqrt{9+4q}),2,T_0),`$
$`b_0`$ $`=`$ $`(k\lambda _0q(1T_0^2)+(3+\sqrt{9+4q})T_0)Q({\displaystyle \frac{1}{2}}(1+\sqrt{9+4q}),2,T_0)`$ (269)
$`+(3+\sqrt{9+4q})Q({\displaystyle \frac{1}{2}}(3+\sqrt{9+4q}),2,T_0),`$
$`a_L`$ $`=`$ $`(k\lambda _Lq(1T_L^2)+(3+\sqrt{9+4q})T_L)P({\displaystyle \frac{1}{2}}(1+\sqrt{9+4q}),2,T_L)`$
$`(3+\sqrt{9+4q})P({\displaystyle \frac{1}{2}}(3+\sqrt{9+4q}),2,T_L),`$
$`b_L`$ $`=`$ $`(k\lambda _Lq(1T_L^2)+(3+\sqrt{9+4q})T_L)Q({\displaystyle \frac{1}{2}}(1+\sqrt{9+4q}),2,T_L)`$
$`(3+\sqrt{9+4q})Q({\displaystyle \frac{1}{2}}(3+\sqrt{9+4q}),2,T_L),`$
with $`q=\frac{m^2}{H^2}`$. The condition
$`0`$ $`=`$ $`a_0b_La_Lb_0,`$ (270)
determines the mass spectrum of the massive graviton. Up to an overall normalization (265) can now be written as
$`b_{\mu \nu }=B_{\mu \nu }\{b_0P(l,2,z)a_0Q(l,2,z)\}.`$ (271)
It seems that (270) can be solved by $`q=2`$, which implies the emergence of a tachyon. But actually this is just an artifact of (270): when $`m^2/H^2=2`$, i.e., $`l=0`$, $`P(0,2,z)`$ identically vanishes, and we need another independent solution. Solving (264) with $`m^2/H^2=2`$, we get
$`b_{\mu \nu }^{(q=2)}=A_{\mu \nu }^{(q=2)}{\displaystyle \frac{1+z^2}{1z^2}}+B_{\mu \nu }^{(q=2)}{\displaystyle \frac{z}{1z^2}},`$ (272)
which can satisfy (266-267) only by $`A_{\mu \nu }^{(q=2)}=B_{\mu \nu }^{(q=2)}=0`$. That is, $`b_{\mu \nu }^{(q=2)}=0`$ and therefore there is no tachyon.
The final result is that the physical degrees of freedom in the massive sector consist of a Kaluza-Klein tower of massive gravitons from $`b_{\mu \nu }(x,y)`$, with 5 DOF each.
#### 5.4.2 โmasslessโ case
For modes which do not satisfy the โmassiveโ conditions i and ii, we should use the curved space version of (311):
$`h_{\mu \nu }`$ $`=`$ $`\beta _{\mu \nu }+\stackrel{~}{}_\mu v_\nu +\stackrel{~}{}_\nu v_\mu +a^2\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }\stackrel{~}{}^2\right)\phi _1`$ (273)
$`+\stackrel{~}{}_\mu n_\nu +\stackrel{~}{}_\nu n_\mu +c_{\mu \nu }+g_{\mu \nu }\phi _2.`$
In this decomposition, vector and scalar modes are annihilated by $`๐_3`$ and $`๐_4`$, respectively, while tensor modes (see (263)) are annihilated by $`๐_2`$.
Equation (242) gives
$`2\stackrel{~}{}_\mu n^\mu {}_{}{}^{}+4\phi _2^{}={\displaystyle \frac{f_1}{a^2}}+{\displaystyle \frac{4a^{}}{a}}F\psi ,`$ (274)
so (241) and (240) become
$`\stackrel{~}{}_\mu \stackrel{~}{}_\nu c^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{3a^{}}{a}}{\displaystyle \frac{f_1}{a^2}},`$ (275)
$`\stackrel{~}{}_\nu c_\mu ^\nu ^{}`$ $`=`$ $`_\mu ({\displaystyle \frac{f_1}{a^2}}+{\displaystyle \frac{a^{}}{a}}F\psi \stackrel{~}{}_\nu n^\nu {}_{}{}^{}\phi _2^{}).`$ (276)
Since
$`\stackrel{~}{}^2\stackrel{~}{}^\nu c_{\mu \nu }`$ $`=`$ $`\stackrel{~}{}^\nu \stackrel{~}{}^2c_{\mu \nu }+{\displaystyle \frac{5H^2}{a^2}}\stackrel{~}{}^\nu c_{\mu \nu }={\displaystyle \frac{3H^2}{a^2}}\stackrel{~}{}^\nu c_{\mu \nu },`$ (277)
$`\stackrel{~}{}^2\stackrel{~}{}_\mu \stackrel{~}{}_\nu n^\nu `$ $`=`$ $`\stackrel{~}{}_\mu \stackrel{~}{}_\nu \stackrel{~}{}^2n^\nu ={\displaystyle \frac{3H^2}{a^2}}\stackrel{~}{}_\mu \stackrel{~}{}_\nu n^\nu ,`$ (278)
$`\stackrel{~}{}^2\stackrel{~}{}_\mu (\mathrm{scalar})`$ $`=`$ $`\left(\stackrel{~}{}_\mu \stackrel{~}{}^2{\displaystyle \frac{3H^2}{a^2}}\stackrel{~}{}_\mu \right)(\mathrm{scalar})={\displaystyle \frac{H^2}{a^2}}\stackrel{~}{}_\mu (\mathrm{scalar}),`$ (279)
acting $`๐_3`$ on (276) reduces it into
$`0={\displaystyle \frac{H^2}{a^2}}_\mu \left({\displaystyle \frac{f_1}{a^2}}+{\displaystyle \frac{a^{}}{a}}F\psi \phi _2^{}\right),`$ (280)
or
$`\phi _2^{}={\displaystyle \frac{f_1}{a^2}}+{\displaystyle \frac{a^{}}{a}}F\psi .`$ (281)
Then (274) gives
$`\stackrel{~}{}_\mu n^\mu {}_{}{}^{}={\displaystyle \frac{3f_1}{2a^2}}.`$ (282)
The trace of the brane-boundary equation (243) is
$`0`$ $`=`$ $`\left[3\theta _i\stackrel{~}{h}^{}{\displaystyle \frac{3\lambda _iH^2}{a^2}}\stackrel{~}{h}+12kT_iF\psi \lambda _i(\stackrel{~}{}_\rho \stackrel{~}{}_\sigma h^{\rho \sigma }\stackrel{~}{}^2\stackrel{~}{h})\right]_{y=y_i}`$ (283)
$`=`$ $`3\theta _i(1+k\lambda _iT_i){\displaystyle \frac{f_1(x)}{a^2}}|_{y=y_i},`$
i.e.,
$`f_1(x)=0.`$ (284)
Since $`\stackrel{~}{}^\mu c_{\mu \nu }0`$ and $`\stackrel{~}{}^\mu n_\mu 0`$, (275) and (282) give
$`c_{\mu \nu }=0,n^\mu {}_{}{}^{}=0.`$ (285)
Now that $`n^\mu `$ has the same $`y`$-dependence as $`\xi ^\mu `$, it is gauged away. Also from (281) we get
$`\phi _2=f_2(x)+\left({\displaystyle \frac{a^{}}{a}}H^2๐\right)\psi ,`$ (286)
where $`๐^{}(y)=/a^2`$.
With $`t_{\mu \nu }=\beta _{\mu \nu }+\stackrel{~}{}_\mu v_\nu +\stackrel{~}{}_\nu v_\mu `$, (239) and (243) become
$`0`$ $`=`$ $`t_{\mu \nu }{}_{}{}^{\prime \prime }+{\displaystyle \frac{2H^2+4a^{}^2}{a^2}}t_{\mu \nu }`$
$`\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }\right)\left\{a^2\phi _1^{\prime \prime }+4aa^{}\phi _1^{}2H^2\phi _1+2f_2+\left(F+{\displaystyle \frac{2a^{}}{a}}2H^2๐\right)\psi \right\},`$
$`0`$ $`=`$ $`[\theta _it_{\mu \nu }{}_{}{}^{}2kT_it_{\mu \nu }`$
$`+(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu })\{\theta _ia^2\phi _1^{}+\lambda _iH^2\phi _1\lambda _if_2(x)\lambda _i({\displaystyle \frac{a^{}}{a}}H^2๐)\psi \}]_{y=y_i}.`$
By construction $`๐_2`$ kills $`\beta _{\mu \nu }`$, and since
$`\stackrel{~}{}^2\stackrel{~}{}_{(\mu }v_{\nu )}=\stackrel{~}{}_{(\mu }๐_3v_{\nu )}{\displaystyle \frac{2H^2}{a^2}}\stackrel{~}{}_{(\mu }v_{\nu )}={\displaystyle \frac{2H^2}{a^2}}\stackrel{~}{}_{(\mu }v_{\nu )},`$ (289)
$`\stackrel{~}{}_\mu v_\nu +\stackrel{~}{}_\nu v_\mu `$ is also annihilated by $`๐_2`$. But, acting on a scalar,
$`๐_2\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }\right)=\left({\displaystyle \frac{4H^2}{a^2}}+{\displaystyle \frac{2H^2}{a^2}}\right)\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }\right).`$ (290)
Therefore, by applying $`๐_2`$ to (5.4.2) and (5.4.2) we separate scalar parts from the remaining. The separated scalar parts have the form
$`\stackrel{~}{}_\mu \stackrel{~}{}_\nu (\mathrm{scalars})={\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }(\mathrm{scalars}),`$ (291)
which can only be solved by $`(\mathrm{scalars})=0`$. Thus (5.4.2) gives
$`0`$ $`=`$ $`t_{\mu \nu }{}_{}{}^{\prime \prime }+{\displaystyle \frac{2H^2+4a^{}^2}{a^2}}t_{\mu \nu },`$ (292)
$`0`$ $`=`$ $`a^2\phi _1^{\prime \prime }+4aa^{}\phi _1^{}2H^2\phi _1+2f_2+\left(F+{\displaystyle \frac{2a^{}}{a}}2H^2๐\right)\psi ,`$ (293)
while from (5.4.2) we get
$`0`$ $`=`$ $`\left[\theta _it_{\mu \nu }{}_{}{}^{}2kT_it_{\mu \nu }\right]_{y=y_i},`$ (294)
$`0`$ $`=`$ $`[\theta _ia^2\phi _1^{}+\lambda _iH^2\phi _1\lambda _if_2(x)\lambda _i({\displaystyle \frac{a^{}}{a}}H^2๐)\psi \}]_{y=y_i}.`$ (295)
Introducing $`z=\mathrm{tanh}k(yy_0)`$, the most general solution of (292) is
$`t_{\mu \nu }(x,y)=A_{\mu \nu }(x){\displaystyle \frac{z\frac{z^3}{3}}{1z^2}}+B_{\mu \nu }(x){\displaystyle \frac{1}{1z^2}}.`$ (296)
The boundary conditions provided by (294) requires $`A_{\mu \nu }=0`$. Thus,
$`t_{\mu \nu }=B_{\mu \nu }(x){\displaystyle \frac{1}{1z^2}}.`$ (297)
Since $`1/(1z^2)=\mathrm{cosh}^2k(yy_0)=a^2(y)\mathrm{cosh}^2ky_0`$, (297) is up to overall normalization
$`t_{\mu \nu }(x,y)=a^2(y)B_{\mu \nu }(x).`$ (298)
Then $`v_\mu `$ has the correct $`y`$-dependence to be gauged away, leaving only $`\beta _{\mu \nu }`$.
(293) has a general solution
$`\phi _1(x,y)={\displaystyle \frac{f_2(x)}{H^2}}+(1z)^2C(x)+zD(x)๐(y)\psi (x),`$ (299)
and $`h_{\mu \nu }`$ becomes
$`h_{\mu \nu }`$ $`=`$ $`\beta _{\mu \nu }+a^2\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }\right){\displaystyle \frac{f_2(x)}{H^2}}+g_{\mu \nu }\left\{f_2(x)+\left({\displaystyle \frac{a^{}}{a}}H^2๐\right)\psi \right\}`$ (300)
$`+a^2\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }\right)\left((1z)^2C(x)+zD(x)๐(y)\psi (x)\right)`$
$`=`$ $`\beta _{\mu \nu }+a^2\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{f_2}{H^2}}\left(a^2๐\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{a^{}}{a}}g_{\mu \nu }\right)\psi `$
$`+a^2\left(\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{H^2}{a^2}}g_{\mu \nu }\right)\left((1z)^2C(x)+zD(x)\right).`$
Then we can see that $`f_2`$ can be gauged away.
Now (295) gives
$`\alpha _0C+\beta _0D`$ $`=`$ $`{\displaystyle \frac{k}{H^2}}\beta _0(0)\psi ,`$
$`\alpha _LC+\beta _LD`$ $`=`$ $`{\displaystyle \frac{k}{H^2}}\beta _L(L)\psi ,`$
where
$`\alpha _0=2(1+T_0)+k\lambda _0(1+T_0)^2,\beta _0=1+k\lambda _0T_0,`$
$`\alpha _L=2(1T_L)k\lambda _L(1T_L)^2,\beta _L=1+k\lambda _LT_L.`$
$`C(x)`$ and $`D(x)`$ can be solved;
$`C`$ $`=`$ $`{\displaystyle \frac{k}{H^2}}{\displaystyle \frac{\beta _0\beta _L}{\alpha _0\beta _L\alpha _L\beta _0}}\left((L)(0)\right)\psi ,`$ (301)
$`D`$ $`=`$ $`{\displaystyle \frac{k}{H^2}}{\displaystyle \frac{\alpha _0\beta _L(L)\alpha _L\beta _0(0)}{\alpha _0\beta _L\alpha _L\beta _0}}\psi .`$ (302)
We can use the gauge freedom of $`F(y)`$ to simplify $`C`$ and $`D`$. For example, choosing
$`k(y)={\displaystyle \frac{y}{L}}\left({\displaystyle \frac{\alpha _0}{\beta _0}}{\displaystyle \frac{\alpha _L}{\beta _L}}\right)+{\displaystyle \frac{\alpha _0}{\beta _0}},`$ (303)
gives
$`C={\displaystyle \frac{1}{H^2}}\psi ,D=0.`$ (304)
All the scalars are written in terms of $`\psi (x)`$.
In summary, the physical degrees of freedom of the massless sector consist of a massless graviton from $`\beta _{\mu \nu }(x)`$ with two on-shell degrees of freedom, and a massless radion $`\psi (x)`$.
## 6 Conclusion
In this paper we have developed a detailed methodology for analyzing models of braneworld gravity. We have used the interval picture, in which braneworld gravity has a well-defined action principle. The key result is equation (5.1), which gives the full variation of the braneworld gravity action with respect to an arbitrary metric variation. From this, we obtain the usual bulk Einstein equations, supplemented by additional constraints which we call โbrane-boundaryโ equations.
The brane-boundary equations are generally covariant, even for coordinate transformations that change the boundary. An immediate consequence of our result is that there are no physical โbrane-bendingโ modes of the 5d metric in braneworld gravity, as one would expect if general covariance were partially broken. This is important since scalar modes can lead to strong coupling behavior and kinetic ghosts. In the general class of models considered in this paper, the radion and the KK gravitons are the only possible sources of such pathologies.
We have introduced the concept of straight gauges, and showed how it is always possible to reach a straight gauge starting from an arbitrary bulk coordinate system. Then we showed how the analysis of linearized metric fluctuations and their equations of motion simplify in a straight gauge. The equations of motion for metric fluctuations of higher dimensional gravity have previously been analyzed in axial, harmonic, de Donder, or Gaussian normal gauges. However, for braneworld setups with more than one brane, none of these gauge choices corresponds to a straight gauge in a single coordinate patch.
In ยง3, ยง4, and ยง5, we have explicitly gauge-fixed and solved the equations of motions for setups with two branes, and 5d backgrounds that are flat, warped Randall-Sundrum, or general warped $`AdS_5/AdS_4`$. In all three cases we define a family of straight gauges. The straight gauges are parametrized by a single function $`F(y)`$, that obeys the condition (229) but is otherwise arbitrary.
The greatest practical importance of our work is in applications to more complicated models and to more subtle issues. Since we start with a well-defined 5d generally covariant action, and gauge-fix it explicitly to an effective 4d action, there can be no arguments about the counting of physical degrees of freedom, the identification of kinetic ghosts, or the onset of strong coupling behavior (to the extent that such behavior can be accessed starting from a linearized theory). We intend to exploit these advantages in future work.
### Acknowledgments
We thank Ruoyu Bao, Seungyeop Lee, Eduardo Pontรณn, Josรฉ Santiago, and Robert Wald for useful discussions. This research was supported by the U.S. Department of Energy Grants DE-AC02-76CHO3000 and DE-FG02-90ER40560.
## Tensor decomposition
A massive symmetric tensor field $`T_{\mu \nu }`$ in flat 4d spacetime has the decomposition
$`T_{\mu \nu }=b_{\mu \nu }+_\mu V_\nu +_\nu V_\mu +_\mu _\nu \varphi _1+\eta _{\mu \nu }\varphi _2,`$ (305)
where
$`b\eta ^{\mu \nu }b_{\mu \nu }`$ $`=`$ $`0,^\mu b_{\mu \nu }=0,`$ (306)
$`^\mu V_\mu `$ $`=`$ $`0.`$ (307)
(306) provides $`4+1`$ conditions, and then $`b_{\mu \nu }`$ has only $`105=5`$ DOF. Similarly, $`V_\mu `$ has $`41=3`$ DOF due to 1 condition imposed by (307). Obviously, $`\varphi _1`$ and $`\varphi _2`$ have one DOF each.
When the 4d fields are massless, i.e., $`^2(\mathrm{fields})=0`$, both $`_\mu V_\nu +_\nu V_\mu `$ and $`_\mu _\nu \varphi _1`$ become transverse-traceless, and a simple decomposition like (305) breaks down. Letโs derive the correct decomposition in the massless case.
First of all, we know that 4 out of 10 DOF of $`T_{\mu \nu }`$ should be expressed in a pure gauge form, $`_\mu V_\nu +_\nu V_\mu `$. The vector $`V_\mu `$ should have three transverse components, one of which can be written as the gradient of a massless scalar $`_\mu \phi _1`$. Let the transverse vector $`v_\mu `$ denote the other two transverse DOF, and $`n_\mu `$ denote the longitudinal component, so:
$`V_\mu =v_\mu +_\mu \phi _1+n_\mu ,`$ (308)
with $`_\mu v^\mu =0`$ and $`_\mu n^\mu 0`$.
The DOF of any symmetric tensor can be divided into the following:
1. $`_{\mu \nu }`$: transverse-traceless,
2. $`๐_{\mu \nu }`$: traceless but not transverse,
3. $`๐_{\mu \nu }`$: trace piece, which we can take to be proportional to $`\eta _{\mu \nu }`$.
Obviously $``$ has $`105=5`$ DOF. We have already exhibited 3 of them; $`_\mu v_\nu +_\nu v_\mu `$ and $`_\mu _\nu \phi _1`$. Therefore,
$`_{\mu \nu }=\beta _{\mu \nu }+_\mu v_\nu +_\nu v_\mu +_\mu _\nu \phi _1,`$ (309)
where $`\beta _{\mu \nu }`$ is a traceless-transverse tensor with 2 DOF.
The sum of the DOF of $``$ and $`๐`$ is 9, so $`๐`$ has 4 DOF. One of these is the pure gauge DOF $`n_\mu `$; we can write:
$`๐_{\mu \nu }=c_{\mu \nu }+_\mu n_\nu +_\nu n_\mu {\displaystyle \frac{1}{2}}\eta _{\mu \nu }_\rho n^\rho ,`$ (310)
where $`c_{\mu \nu }`$ is a traceless but not transverse tensor with 3 DOF.
Collecting the pieces, we get the decomposition of a massless tensor:
$`T_{\mu \nu }=\beta _{\mu \nu }+_\mu v_\nu +_\nu v_\mu +_\mu _\nu \phi _1+c_{\mu \nu }+_\mu n_\nu +_\nu n_\mu +\eta _{\mu \nu }\phi _2.`$ (311)
Letโs look at the massless decomposition in momentum space, i.e., consider the decomposition of $`\overline{T}_{\mu \nu }(p)=d^4xT_{\mu \nu }e^{ipx}`$, with $`p^2=0`$. When $`p^\mu `$ is null, it is not possible to find three vectors which are mutually orthogonal and transverse to $`p^\mu `$. Instead, we introduce the following explicit basis:
$`ฯต_\mu ^{(1)}:\mathrm{parallel}\mathrm{to}p_\mu (\mathrm{helicity}+1),ฯต^{(1)\mu }ฯต_\mu ^{(1)}=0;`$
$`ฯต_\mu ^{(2)}:\mathrm{antiparallel}\mathrm{to}p_\mu (\mathrm{helicity}1),ฯต^{(2)\mu }ฯต_\mu ^{(2)}=0,ฯต^{(1)\mu }ฯต_\mu ^{(2)}0;`$ (312)
$`ฯต_\mu ^{(j)}(j=3,\mathrm{\hspace{0.33em}4}):ฯต^{(1)\mu }ฯต_\mu ^{(j)}=ฯต^{(2)\mu }ฯต_\mu ^{(j)}=0,ฯต^{(j)\mu }ฯต_\mu ^{(k)}=\delta _{jk},`$
from which we can build bases for second rank symmetric tensors:
$`\epsilon _{\mu \nu }^{(1)}=ฯต_{(\mu }^{(3)}ฯต_{\nu )}^{(4)},\epsilon _{\mu \nu }^{(2)}=ฯต_\mu ^{(3)}ฯต_\nu ^{(3)}ฯต_\mu ^{(4)}ฯต_\nu ^{(4)},`$
$`\epsilon _{\mu \nu }^{(3)}=ฯต_{(\mu }^{(1)}ฯต_{\nu )}^{(3)},\epsilon _{\mu \nu }^{(4)}=ฯต_{(\mu }^{(1)}ฯต_{\nu )}^{(4)},`$
$`\epsilon _{\mu \nu }^{(5)}=ฯต_\mu ^{(1)}ฯต_\nu ^{(1)},`$
$`\epsilon _{\mu \nu }^{(6)}=ฯต_{(\mu }^{(1)}ฯต_{\nu )}^{(2)},`$
$`\epsilon _{\mu \nu }^{(7)}=ฯต_\mu ^{(2)}ฯต_\nu ^{(2)},\epsilon _{\mu \nu }^{(8)}=ฯต_{(\mu }^{(2)}ฯต_{\nu )}^{(3)},\epsilon _{\mu \nu }^{(9)}=ฯต_{(\mu }^{(2)}ฯต_{\nu )}^{(4)},`$
$`\epsilon _{\mu \nu }^{(10)}=ฯต_{(\mu }^{(1)}ฯต_{\nu )}^{(2)}+ฯต_\mu ^{(3)}ฯต_\nu ^{(3)}+ฯต_\mu ^{(4)}ฯต_\nu ^{(4)}.`$
With (Tensor decomposition), we can read off characteristics of each basis component:
1. $`\epsilon ^{(12)}`$: traceless-transverse, and transverse to $`ฯต_\mu ^{(2)}`$,
2. $`\epsilon ^{(35)}`$: traceless-transverse,
3. $`\epsilon ^{(6)}`$: neither traceless nor transverse,
4. $`\epsilon ^{(79)}`$: traceless but not transverse, and transverse to $`ฯต_\mu ^{(2)}`$,
5. $`\epsilon ^{(10)}\eta _{\mu \nu }`$.
Then using the fact that $`ฯต_\mu ^{(1)}`$ is parallel to $`p_\mu `$, we can get the decomposition
$`\overline{T}_{\mu \nu }=\overline{\beta }_{\mu \nu }+ip_\mu \overline{v}_\nu +ip_\nu \overline{v}_\mu p_\mu p_\nu \overline{\phi }_1+ip_\mu \overline{n}_\nu +ip_\nu \overline{n}_\mu +\overline{c}_{\mu \nu }+\eta _{\mu \nu }\overline{\phi }_2,`$ (313)
where $`\overline{n}_\mu `$ is proportional to $`ฯต_\mu ^{(2)}`$, and
1. $`\overline{\beta }_{\mu \nu }(\epsilon ^{(12)})`$: traceless-transverse, and transverse to $`\overline{n}_\mu `$, 2 DOF;
2. $`p_\mu \overline{v}_\nu +p_\nu \overline{v}_\mu (\epsilon ^{(34)})`$: traceless-transverse, where $`\overline{v}_\mu `$ is transverse to $`p^\mu `$ and to $`\overline{n}_\mu `$, 2 DOF;
3. $`p_\mu p_\nu \overline{\phi }_1(\epsilon ^{(5)})`$: traceless-transverse, 1 DOF;
4. $`p_\mu \overline{n}_\nu +p_\nu \overline{n}_\mu (\epsilon ^{(6)})`$: 1 DOF. $`p_\mu \overline{n}^\mu 0`$;
5. $`\overline{c}_{\mu \nu }(\epsilon ^{(79)})`$: traceless but not transverse, and transverse to $`\overline{n}_\mu `$, 3 DOF;
6. $`\eta _{\mu \nu }\overline{\phi }_2(\epsilon ^{(10)})`$: 1 DOF. |
warning/0506/quant-ph0506224.html | ar5iv | text | # State space structure and entanglement of rotationally invariant spin systems
## I Introduction
Entanglement is a basic feature of composite quantum systems connected to the tensor product structure of the underlying Hilbert space of states. A mixed state of a bipartite quantum system described by some density matrix $`\rho `$ is said to be entangled or inseparable if $`\rho `$ cannot be written as a convex linear combination of product states. Otherwise it is called classically correlated or separable WERNER . The properties of entangled states are responsible for many of the fascinating and curious aspects of the quantum world and lie at the core of many proposed applications in quantum information processing ECKERT ; ALBER ; NIELSEN .
The general characterization and quantification of entanglement in mixed quantum states is a highly non-trivial problem. It is even very difficult in general to formulate simple operational criteria which allow a unique identification of all separable states of a given composite system. There do exist, however, many necessary separability criteria PERES ; HORODECKI96a ; HORODECKI99 ; CERF ; KEMPE ; RUDOLPH03 ; CHEN ; TERHAL ; LEWENSTEIN00 . A simple and, in fact, very strong criterion is the Peres-Horodecki criterion PERES ; HORODECKI96a which states that a necessary condition for a given density matrix $`\rho `$ to be separable is that it has a positive partial transposition (PPT states). It is known that this criterion is necessary and sufficient for certain low-dimensional systems, while it is only necessary in higher dimensions HORODECKI96a .
The analysis of the entanglement structure is greatly facilitated through the introduction of symmetries, i. e., if one restricts to those states of the composite system which are invariant under certain groups of symmetry transformations. Important examples in this context are the manifolds of the Werner states WERNER , of the isotropic states RAINS ; HORODECKI99 and of the orthogonal states VOLLBRECHT . Here, we investigate entanglement under the symmetry group $`\mathrm{SO}(3)`$ of proper three-dimensional rotations of the coordinate axes. More precisely, we consider the problem of mixed state entanglement of systems which are composed of two particles with spins $`j_1`$ and $`j_2`$, and which are invariant under product representations of the group $`\mathrm{SO}(3)`$ or, equivalently, of the covering group $`\mathrm{SU}(2)`$. A basic tool of our analysis is the work of Vollbrecht and Werner VOLLBRECHT which provides a general scheme for the treatment of entanglement under given symmetry groups.
Mixed $`\mathrm{SO}(3)`$-invariant states of composite systems arise, for example, from the interaction of open systems TheWork with isotropic environments GORINI . Their analysis is of great importance and leads to many applications. As examples we mention investigations on the connection between quantum phase transitions and the behaviour of entanglement measures (see, e. g., OSTERLOH ; OSBORNE ), the analysis of entanglement of $`\mathrm{SU}(2)`$-invariant multiphoton states generated by parametric down-conversion DURKIN , and studies of the entanglement of formation CAVES . The technique of this paper could also be relevant for the characterization of quantum correlations in Fermionic or Bosonic systems developed recently SCHLIEMANN01 ; BRUSS .
The Hilbert space of a system which is composed of two particles with spins $`j_1`$ and $`j_2`$ is given by the tensor product $`^{N_1}^{N_2}`$, where $`N_1=2j_1+1`$ and $`N_2=2j_2+1`$ are the dimensions of the local spin spaces. We call such a system an $`N_1N_2`$ system. Throughout the paper we will assume that $`j_1j_2`$, i. e., $`N_1N_2`$.
According to the Peres-Horodecki criterion PERES ; HORODECKI96a the cases of $`22`$ and $`23`$ systems are trivial: It is known that in these cases the PPT criterion is necessary and sufficient for all states, i. e., even for states which are not invariant under rotations. Schliemann SCHLIEMANN1 has shown recently that the PPT criterion is also necessary and sufficient for $`\mathrm{SO}(3)`$-invariant $`2N_2`$ systems with arbitrary $`N_2`$. The case of $`33`$ systems has been treated by Vollbrecht and Werner VOLLBRECHT , who proved that the PPT criterion is again necessary and sufficient for separability. For $`44`$ systems a qualitatively new situation arises: It has been demonstrated in NtensorN that the PPT criterion is not sufficient and that the entangled PPT states form a three-dimensional manifold which is isomorphic to a prism. In the present work we investigate the important special case of $`3N_2`$ systems with arbitrary $`N_2`$.
The method developed in NtensorN enables the treatment of the case of equal spins $`j_1=j_2`$. In this paper we extend this method to arbitrary spins $`j_1`$ and $`j_2`$. For the analysis of entanglement under $`\mathrm{SO}(3)`$-symmetry it is advantageous to replace the transposition used in the PPT criterion by another unitarily equivalent operation, namely by the antiunitary transformation of the time reversal. The reason for this fact is that the operation of the time reversal of states commutes with the representations of the rotation group.
There are two natural representations of rotationally invariant states. The first one uses the fact that any invariant state can be written as a unique convex linear combination of the projections $`P_J`$ onto the eigenspaces of the total angular momentum $`J`$ of the composite spin system. The advantage of this representation is that it leads to very simple conditions expressing the positivity and the normalization of physical states. However, the set of the PPT states is most easily determined in another representation which employs the irreducible spherical tensor operators of spin-$`j`$ particles. We will construct a complete system of invariant operators $`Q_K`$ which are built out of the spherical tensors of rank $`K`$. Any invariant state of the composite spin system can then be written as a unique linear combination of the $`Q_K`$. The introduction of the invariant operators $`Q_K`$ generalizes the ideas of Schliemann SCHLIEMANN1 ; SCHLIEMANN2 , who has developed a representation of $`\mathrm{SU}(2)`$-invariant states by means of spin-spin correlators and has formulated various separability conditions and sum rules in terms of these correlators.
The paper is organized as follows. The representations of $`\mathrm{SO}(3)`$-invariant states in terms of the invariant operators $`P_J`$ and $`Q_K`$ are constructed in Sec. II. We also derive in this section the linear transformation which connects these representations and show that it is given by an orthogonal matrix whose elements are determined by Wignerโs $`6`$-$`j`$ symbols. The behaviour of states under partial time reversal and the construction of the set of the invariant separable states are discussed in Sec. III.
The general theory is then applied in Sec. IV to the case of $`3N_2`$ systems with arbitrary $`N_2`$. We prove that the PPT criterion represents a necessary and sufficient separability condition for $`3N_2`$ systems if and only if $`N_2`$ is odd. Thus, for integer spins $`j_2`$ all PPT states are separable, while for half-integer spins $`j_2`$ there always exist entangled PPT states. This fact has already been conjectured by Hendriks HENDRIKS on the basis of a detailed numerical investigation. We also show that for half-integer $`j_2`$ the boundary of the separability region is curved. Finally, Sec. V contains a discussion of the results and some conclusions. In particular, we construct an optimal entanglement witness for the case of half-integer spins and exploit this witness to design a protocol which allows the detection of entangled PPT states through measurements of the total angular momentum.
## II Representations of $`\mathrm{SO}(3)`$-invariant states
We consider two particles with spins $`j_1`$ and $`j_2`$ and corresponding angular momentum operators $`\widehat{๐}^{(1)}`$ and $`\widehat{๐}^{(2)}`$. The Hilbert space $`^{N_1}`$ of the first particle is spanned by the common eigenstates $`|j_1,m_1`$ of the square of $`\widehat{๐}^{(1)}`$ and of $`\widehat{๐}_z^{(1)}`$, where $`N_12j_1+1`$ and $`m_1=j_1,\mathrm{},+j_1`$. Correspondingly, the Hilbert space $`^{N_2}`$ of the second particle is spanned by the eigenstates $`|j_2,m_2`$, where $`N_22j_2+1`$ and $`m_2=j_2,\mathrm{},+j_2`$.
The Hilbert space of the total system composed of both particles is given by the tensor product $`^{N_1}^{N_2}`$. The angular momentum operator of the composite system is defined by:
$$\widehat{๐ฑ}=\widehat{๐}^{(1)}I+I\widehat{๐}^{(2)},$$
(1)
where $`I`$ denotes the unit matrix. A state of the composite system is described by a density matrix on the product space, i. e., by a positive operator $`\rho `$ on $`^{N_1}^{N_2}`$ with unit trace: $`\rho 0`$, $`\mathrm{tr}\rho =1`$.
The irreducible unitary representation of the group $`\mathrm{SO}(3)`$ of proper rotations $`R`$ on the state space of a particle with spin $`j`$ will be denoted by $`D^{(j)}(R)`$. The transformation of the states of the composite spin system is then given by the product representation $`D^{(j_1)}(R)D^{(j_2)}(R)`$. A state $`\rho `$ of the combined system is said to be rotationally invariant or $`\mathrm{SO}(3)`$-invariant if the relation
$$\left[D^{(j_1)}(R)D^{(j_2)}(R)\right]\rho \left[D^{(j_1)}(R)D^{(j_2)}(R)\right]^{}=\rho $$
holds true for all proper rotations $`RSO(3)`$.
We shall use two different representations of rotationally invariant states. The first one employs the projection operators
$$P_J=\underset{M=J}{\overset{+J}{}}|JMJM|,$$
(2)
where $`|JM`$ denotes the common eigenstate of the square of the total angular momentum $`\widehat{๐ฑ}`$ and of its $`z`$-component $`\widehat{J}_z`$, i. e., we have $`\widehat{๐ฑ}^2|JM=J(J+1)|JM`$ and $`\widehat{J}_z|JM=M|JM`$. The operator $`P_J`$ projects onto the manifold which is spanned by the eigenstates belonging to a fixed value $`J`$ of the total angular momentum. According to the triangular inequality $`J`$ takes on $`N_1`$ different values which may be integer or half-integer valued:
$$J=j_2j_1,j_2j_1+1,\mathrm{},j_2+j_1.$$
(3)
It follows from Schurโs lemma that any invariant state $`\rho `$ can be written as a linear combination of the $`P_J`$:
$$\rho =\frac{1}{\sqrt{N_1N_2}}\underset{J}{}\frac{\alpha _J}{\sqrt{2J+1}}P_J.$$
(4)
Here, the $`\alpha _J`$ are real parameters and we have introduced convenient normalization factors of $`\sqrt{N_1N_2}`$ and $`\sqrt{2J+1}`$. In order for Eq. (4) to represent a true density matrix the $`\alpha _J`$ must of course be positive and normalized appropriately:
$`\alpha _J`$ $``$ $`0,`$ (5)
$`\mathrm{tr}\rho `$ $`=`$ $`{\displaystyle \underset{J}{}}\sqrt{{\displaystyle \frac{2J+1}{N_1N_2}}}\alpha _J=1.`$ (6)
Any invariant state $`\rho `$ is thus uniquely characterized by a real vector $`๐ถ`$ in an $`N_1`$-dimensional parameter space $`^{N_1}`$ which will be referred to as $`\alpha `$-space. The conditions of the positivity and of the normalization of $`\rho `$ are expressed by the relations (5) and (6). We denote the set of all vectors $`๐ถ`$ whose components $`\alpha _J`$ satisfy these relations by $`S^\alpha `$. Being isomorphic to the set of invariant states, $`S^\alpha `$ is of course a convex set. We infer from Eqs. (5) and (6) that $`S^\alpha `$ represents an $`(N_11)`$-dimensional simplex.
A useful alternative representation of the invariant states is obtained by use of a complete system of irreducible spherical tensor operators (see, e. g. EDMONDS ; SCHWINGER ). The tensor operators which act on the state space of the particle with spin $`j_i`$ are written as $`T_{K_iq_i}^{(i)}`$, where $`i=1,2`$. The index $`K_i=0,1,\mathrm{},2j_i`$ denotes the rank of the tensor operator. For a given rank $`K_i`$ the index $`q_i`$ takes on the values $`q_i=K_i,K_i+1,\mathrm{},+K_i`$. We thus have $`(2K_i+1)`$ tensor operators $`T_{K_iq_i}^{(i)}`$ of rank $`K_i`$ which transform under rotations according to an irreducible representation of the rotation group. The explicit definitions of the tensors and a brief summary of their properties are given in Appendix A.
Using the tensor operators one defines Hermitian operators $`Q_K`$ acting on the state space of the composite spin system:
$$Q_K=\underset{q=K}{\overset{+K}{}}T_{Kq}^{(1)}T_{Kq}^{(2)},$$
(7)
where the index $`K`$ takes on $`N_1`$ different integer values:
$$K=0,1,\mathrm{},2j_1.$$
(8)
It follows from the transformation properties of the tensor operators that all $`Q_K`$ are invariant under rotations. For instance, the operator $`Q_0`$ is proportional to the identity, $`Q_0=\frac{1}{\sqrt{N_1N_2}}II`$, while $`Q_1`$ is proportional to the invariant scalar product $`\widehat{๐}^{(1)}\widehat{๐}^{(2)}`$ of the spin vectors.
The $`Q_K`$ defined by Eq. (7) form a complete system of operators. This means that any rotationally invariant Hermitian operator can be represented as a unique linear combination of the $`Q_K`$ in a way analogous to Eq. (4):
$$\rho =\frac{1}{\sqrt{N_1N_2}}\underset{K}{}\frac{\beta _K}{\sqrt{2K+1}}Q_K.$$
(9)
Here, we have again introduced appropriate normalization factors and real parameters $`\beta _K`$ which form a vector $`๐ท`$ in an $`N_1`$-dimensional parameter space $`^{N_1}`$ referred to as $`\beta `$-space. The operators $`Q_K`$ satisfy $`\mathrm{tr}\{Q_KQ_K^{}\}=(2K+1)\delta _{KK^{}}`$. This fact follows directly from the orthogonality relation (76) for the spherical tensors. The $`Q_K`$ for $`K0`$ are therefore traceless which leads to the normalization condition
$$\mathrm{tr}\rho =\beta _0=1.$$
(10)
The sets $`\{P_J\}`$ and $`\{Q_K\}`$ represent complete systems of invariant operators. The corresponding parameter vectors $`๐ถ`$ and $`๐ท`$ must therefore be related by a linear transformation $`^{N_1}^{N_1}`$. We write
$$๐ท=L๐ถ,$$
(11)
where $`L`$ is an $`(N_1\times N_1)`$ matrix. To find the elements of this matrix we use Eqs. (4) and (9) to get
$$\underset{J}{}\frac{\alpha _J}{\sqrt{2J+1}}P_J=\underset{K}{}\frac{\beta _K}{\sqrt{2K+1}}Q_K.$$
(12)
Multiplying this equation by $`Q_K^{}`$ and taking the trace we find that the elements of $`L`$ are given by
$$L_{KJ}=[(2K+1)(2J+1)]^{1/2}\mathrm{tr}\{Q_KP_J\}.$$
(13)
This can be expressed as
$$L_{KJ}=\sqrt{(2K+1)(2J+1)}(1)^{j_1+j_2+J}\left\{\begin{array}{ccc}j_1& j_2& J\\ j_2& j_1& K\end{array}\right\}.$$
(14)
The curly brackets denote a $`6`$-$`j`$ symbol introduced by Wigner WIGNER into the quantum theory of angular momentum. A proof of the relation (14) is given in Appendix B. The $`6`$-$`j`$ symbols are scalar quantities which are defined through invariant sums over products of Clebsch-Gordan coefficients. They describe the transformation between different coupling schemes for the addition of three angular momenta EDMONDS . Their properties have been studied in great detail and many closed formulae, recursion relations and sum rules are known. In particular, it follows from the sum rules that $`L`$ represents an orthogonal $`(N_1\times N_1)`$ matrix.
The above results lead to the conclusion that the set of $`\mathrm{SO}(3)`$-invariant states is represented in $`\beta `$-space by the set
$$S^\beta =LS^\alpha .$$
(15)
The set $`S^\beta `$ is again an $`(N_11)`$-dimensional simplex which may be constructed by determining the images of the extreme points of $`S^\alpha `$ under the orthogonal transformation $`L`$.
The introduction of two parameter spaces is motivated by the following observations. On the one hand, the set of states is most easily constructed as a subset in $`\alpha `$-space. This is due to the fact that the representation of Eq. (4) corresponds to the spectral decomposition of $`\rho `$ and, therefore, the requirement of the positivity of $`\rho `$ immediately leads to the simple condition (5). On the other hand, the representation (9) of states in $`\beta `$-space is much more suitable for the construction of the set of separable states, which is due to the fact that the operation of the partial time reversal is diagonal in the $`Q_K`$-representation.
## III Invariant separable states
A state $`\rho `$ of the composite spin system is said to be separable if it is possible to write this state as a convex linear combination of product states:
$$\rho =\underset{i}{}\lambda _i\rho _i^{(1)}\rho _i^{(2)},\lambda _i0,\underset{i}{}\lambda _i=1,$$
(16)
where the $`\rho _i^{(1)}`$ and $`\rho _i^{(2)}`$ are normalized states of the first and of the second spin, respectively WERNER . It is clear that the set in $`\beta `$-space which represents the invariant and separable states is a convex subset of $`S^\beta `$. This subset will be denoted by $`S_{\mathrm{sep}}^\beta `$.
Following the work of Vollbrecht and Werner VOLLBRECHT we define a projection super-operator ($`\mathrm{SO}(3)`$ twirling) by means of
$$\mathrm{\Pi }\rho =๐RU(R)\rho U(R)^{},$$
(17)
where $`U(R)D^{(j_1)}(R)D^{(j_2)}(R)`$ and the integration is extended over all group elements $`RSO(3)`$. The twirl operation maps any state $`\rho `$ of the composite spin system to an $`\mathrm{SO}(3)`$-invariant state $`\mathrm{\Pi }\rho `$. Moreover, if $`\rho `$ is separable then also $`\mathrm{\Pi }\rho `$ is separable. In terms of the invariant operators $`P_J`$ or $`Q_K`$ the action of the twirl operation may be expressed by
$$\mathrm{\Pi }\rho =\underset{J}{}\frac{\mathrm{tr}\{P_J\rho \}}{2J+1}P_J=\underset{K}{}\frac{\mathrm{tr}\{Q_K\rho \}}{2K+1}Q_K.$$
(18)
It is known that any invariant separable state is a convex linear combination of $`\mathrm{\Pi }`$-projections of pure product states. Given a pure product state
$$\rho =|\phi ^{(1)}\phi ^{(2)}\phi ^{(1)}\phi ^{(2)}|,$$
(19)
Eq. (18) shows that the corresponding parameters $`\alpha _J`$ and $`\beta _K`$ of its projection $`\mathrm{\Pi }\rho `$ are given by
$`\alpha _J`$ $`=`$ $`\sqrt{{\displaystyle \frac{N_1N_2}{2J+1}}}\phi ^{(1)}\phi ^{(2)}|P_J|\phi ^{(1)}\phi ^{(2)},`$ (20)
$`\beta _K`$ $`=`$ $`\sqrt{{\displaystyle \frac{N_1N_2}{2K+1}}}\phi ^{(1)}\phi ^{(2)}|Q_K|\phi ^{(1)}\phi ^{(2)}.`$ (21)
We introduce into Eq. (21) the definition (7) of the $`Q_K`$ and define the functions
$`\stackrel{~}{\beta }_K[\phi ^{(1)},\phi ^{(2)}]`$
$`=\sqrt{{\displaystyle \frac{N_1N_2}{2K+1}}}{\displaystyle \underset{q=K}{\overset{+K}{}}}\phi ^{(1)}|T_{Kq}^{(1)}|\phi ^{(1)}\phi ^{(2)}|T_{Kq}^{(2)}|\phi ^{(2)}.`$
Let us further define $`W^\beta `$ as the range of the parameter vector $`๐ท`$ whose components are given by these functions, where $`|\phi ^{(1)}^{N_1}`$ and $`|\phi ^{(2)}^{N_2}`$ run independently over all normalized states of the first and of the second spin, respectively:
$$W^\beta =\left\{๐ท\right|\beta _K=\stackrel{~}{\beta }_K[\phi ^{(1)},\phi ^{(2)}],\phi ^{(1,2)}=1\}.$$
(23)
The set of separable states is then equal to the convex hull of $`W^\beta `$:
$$S_{\mathrm{sep}}^\beta =\text{hull}\left(W^\beta \right).$$
(24)
This means that $`S_{\mathrm{sep}}^\beta `$ is equal to the smallest convex set which contains $`W^\beta `$.
Within this formulation the problem of constructing $`S_{\mathrm{sep}}^\beta `$ reduces to the determination of the convex hull of the range of the functions $`\stackrel{~}{\beta }_K`$. Even for the present case of a highly symmetric state space this is, in general, an extremely difficult task. A strong necessary condition for separability is the Peres-Horodecki criterion PERES ; HORODECKI96a . According to this criterion a necessary condition for a given state $`\rho `$ to be separable is that its partial transposition is a positive operator: $`T_2\rho (IT)\rho 0`$. Here, $`TB=B^T`$ denotes the transposition of the operator $`B`$ on $`^{N_2}`$ which is defined in terms of the basis states of the second spin by means of $`j_2,m_2|B^T|j_2,m_2^{}=j_2,m_2^{}|B|j_2,m_2`$. The partial transposition $`T_2`$ is then defined by $`T_2(AB)=AB^T`$.
The operation of taking the partial transposition destroys the rotational invariance of states, i. e., if $`\rho `$ is invariant under rotations the partially transposed state $`T_2\rho `$ is generally not $`\mathrm{SO}(3)`$-invariant. However, there exist another operation which is unitarily equivalent to $`T_2`$ and which does map rotationally invariant operators to rotationally invariant operators. This operation will be denoted by $`\vartheta _2=I\vartheta `$. It involves the antiunitary time reversal transformation $`\vartheta `$ of the second spin and will therefore be referred to as partial time reversal.
According to Wignerโs representation theorem WIGNER the action of the time reversal transformation $`\vartheta `$ on an operator $`B`$ can be expressed as:
$$\vartheta B=VB^TV^{}=\tau B^{}\tau ^1.$$
(25)
In the first expression $`T`$ denotes again the transposition and $`V`$ is a unitary matrix which represents a rotation of the coordinate system about the $`y`$-axis by the angle $`\pi `$. In the second expression of Eq. (25) $`\tau `$ denotes the operator $`\tau =V\tau _0`$ which is composed of the $`\pi `$-rotation $`V`$ and of the operator $`\tau _0`$ of the complex conjugation. The operator $`\tau `$ is antiunitary and satisfies
$$\tau ^2=(1)^{2j_2}.$$
(26)
$`\vartheta `$ is a positive but not completely positive map. It is unitarily equivalent to the transposition $`T`$ and, hence, the Peres-Horodecki criterion can be expressed by
$$\vartheta _2\rho (I\vartheta )\rho 0.$$
(27)
A great advantage of the representation of states in $`\beta `$-space is that the operators $`Q_K`$ have a very simple behaviour under the map $`\vartheta _2`$. Namely, as is shown in Appendix A they are eigenoperators of the partial time reversal: $`\vartheta _2Q_K=(1)^KQ_K`$. In $`\beta `$-space the map $`\vartheta _2`$ therefore induces a reflection of the coordinate axes corresponding to the odd values of $`K`$:
$$\vartheta _2:\beta _K(1)^K\beta _K.$$
(28)
We thus get the image $`\vartheta _2S^\beta `$ of $`S^\beta `$ by reversing the signs of the odd coordinates.
We define $`S_{\mathrm{ppt}}^\beta `$ as the set of states which are positive under $`\vartheta _2`$ or, equivalently, under $`T_2`$ (PPT states). This set is equal to the intersection of $`S^\beta `$ with its image $`\vartheta _2S^\beta `$. According to the Peres-Horodecki criterion the set of separable states is a subset of the set of PPT states. Hence, we have
$$S_{\mathrm{sep}}^\beta S_{\mathrm{ppt}}^\beta =S^\beta \vartheta _2S^\beta .$$
(29)
We note three properties which turn out to be useful in the construction of the set of separable states.
(1) The functions defined by Eq. (III) are invariant under simultaneous rotations of the input arguments:
$$\stackrel{~}{\beta }_K[D^{(j_1)}(R)\phi ^{(1)},D^{(j_2)}(R)\phi ^{(2)}]=\stackrel{~}{\beta }_K[\phi ^{(1)},\phi ^{(2)}].$$
(30)
This property is an immediate consequence of the rotational invariance of the operators $`Q_K`$.
(2) The range $`W^\beta `$ defined in Eq. (23) is obviously invariant under the partial time reversal $`\vartheta _2`$. This means that $`๐ทW^\beta `$ implies $`\vartheta _2๐ทW^\beta `$.
(3) There exist two distinguished separable states. These are the state given by the parameter vector $`๐ถ`$ with components
$$\alpha _J=\sqrt{\frac{N_1N_2}{2J_{\mathrm{max}}+1}}\delta _{J,J_{\mathrm{max}}},J_{\mathrm{max}}j_1+j_2,$$
(31)
and the partially time reversed state given by $`๐ถ^{}=\vartheta _2๐ถ`$. To proof this statement we consider a pure product state $`\rho `$ of the form of Eq. (19) with $`|\phi ^{(1)}=|j_1,+j_1`$ and $`|\phi ^{(2)}=|j_2,+j_2`$. We then have the obvious relation $`|J=J_{\mathrm{max}},M=+J_{\mathrm{max}}=|\phi ^{(1)}\phi ^{(2)}`$ and, hence,
$$\phi ^{(1)}\phi ^{(2)}|P_J|\phi ^{(1)}\phi ^{(2)}=\delta _{J,J_{\mathrm{max}}}.$$
(32)
Equation (20) then immediately leads to Eq. (31). This means that the pure product state $`\rho `$ is mapped under the twirl operation to the separable state $`\mathrm{\Pi }\rho =\frac{1}{2J_{\mathrm{max}}+1}P_{J_{\mathrm{max}}}`$ corresponding to the maximal value of the total angular momentum $`J_{\mathrm{max}}`$. It follows from point (2) that also the partially time reversed state is separable.
The point $`๐ถ`$ given by Eq. (31) is an extreme point of the simplex $`S^\alpha `$ and its image $`๐ถ^{}`$ is an extreme point of $`\vartheta _2S^\alpha `$. Thus, $`๐ถ`$ and $`๐ถ^{}`$ are extreme points of $`S_{\mathrm{ppt}}^\alpha `$. It follows that the corresponding points $`๐ท=L๐ถ`$ and $`๐ท^{}=L๐ถ^{}`$ in $`\beta `$-space belong to $`W^\beta `$ and represent extreme points of $`S_{\mathrm{ppt}}^\beta `$.
As an illustration of the above analysis consider a $`2N_2`$ system for which $`j_1=\frac{1}{2}`$ and $`j_2`$ is arbitrary. As has been demonstrated by Schliemann SCHLIEMANN1 the PPT criterion is a necessary and sufficient separability condition in this case. Within the present formulation this statement can be proven as follows. We first note that the index $`K`$ takes on the two values $`K=0,1`$ such that $`๐ท`$ is a two-dimensional vector. Because of the normalization condition (10) we only need a single parameter $`\beta _1`$ to characterize uniquely an invariant state of a $`2N_2`$ system. It follows that $`S^\beta `$ can be represented by an interval of the $`\beta _1`$-axis, and $`S_{\mathrm{ppt}}^\beta `$ by a sub-interval of this interval. Since an interval has exactly two extreme points (its endpoints) we conclude with the help of point (3) above that the extreme points of $`S_{\mathrm{ppt}}^\beta `$ belong to $`W^\beta `$. By the relation (24) the sets $`S_{\mathrm{ppt}}^\beta `$ and $`S_{\mathrm{sep}}^\beta `$ therefore coincide. This shows that the PPT criterion is indeed necessary and sufficient for separability.
## IV $`3N`$ systems
Let us now consider the case $`j_1=1`$ ($`N_1=3`$) and $`j_2`$ arbitrary, i. e. the case of $`3N_2`$ systems. For convenience we write $`NN_2=2j_2+1`$. Since $`J`$ takes on the values $`J=j_21`$, $`j_2`$ and $`j_2+1`$, $`๐ถ`$ is a three-vector
$$๐ถ=\left(\begin{array}{c}\alpha _{j_21}\\ \alpha _{j_2}\\ \alpha _{j_2+1}\end{array}\right).$$
(33)
The set $`S^\alpha `$ of invariant states is given by the relations:
$$\alpha _{j_21},\alpha _{j_2},\alpha _{j_2+1}0$$
(34)
and
$$\sqrt{\frac{N2}{3N}}\alpha _{j_21}+\sqrt{\frac{1}{3}}\alpha _{j_2}+\sqrt{\frac{N+2}{3N}}\alpha _{j_2+1}=1.$$
(35)
We observe that $`S^\alpha `$ is a 2-simplex, i. e. a triangle whose vertices are given by the following parameter vectors $`๐ถ`$:
$$\left(\begin{array}{c}0\\ 0\\ \sqrt{\frac{3N}{N+2}}\end{array}\right),\left(\begin{array}{c}\sqrt{\frac{3N}{N2}}\\ 0\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \sqrt{3}\\ 0\end{array}\right).$$
(36)
In order to transform to $`\beta `$-space we first determine the matrix $`L`$ by means of the formulae (102)-(104):
$`L=`$
$`\left[\begin{array}{ccc}\sqrt{\frac{N2}{3N}}& \sqrt{\frac{1}{3}}& \sqrt{\frac{N+2}{3N}}\\ \sqrt{\frac{(N2)(N+1)}{2N(N1)}}& \sqrt{\frac{2}{(N1)(N+1)}}& \sqrt{\frac{(N1)(N+2)}{2N(N+1)}}\\ \sqrt{\frac{(N+1)(N+2)}{6N(N1)}}& \sqrt{\frac{2(N2)(N+2)}{3(N1)(N+1)}}& \sqrt{\frac{(N1)(N2)}{6N(N+1)}}\end{array}\right].`$
The extreme points of $`S^\beta `$ are found by applying this matrix to the vectors given in Eq. (36). Since $`\beta _0`$ is identically equal to $`1`$ by the normalization condition (10) we can represent points in $`\beta `$-space by two coordinates $`(\beta _1,\beta _2)`$. One finds that $`S^\beta `$ is a triangle in the $`(\beta _1,\beta _2)`$-plane with the vertices:
$$A=(\sqrt{\frac{3(N1)}{2(N+1)}},\sqrt{\frac{(N1)(N2)}{2(N+1)(N+2)}}),$$
(38)
$$B=(\sqrt{\frac{3(N+1)}{2(N1)}},\sqrt{\frac{(N+1)(N+2)}{2(N1)(N2)}}),$$
(39)
$$C=(\sqrt{\frac{6}{(N1)(N+1)}},\sqrt{\frac{2(N2)(N+2)}{(N1)(N+1)}}).$$
(40)
The image $`\vartheta _2S^\beta `$ of $`S^\beta `$ under the partial time reversal is obtained by reversing the sign of the coordinate $`\beta _1`$. Consequently, $`S_{\mathrm{ppt}}^\beta `$ is a polygon with the four vertices $`A`$, $`A^{}`$, $`D`$ and $`E`$, where $`A`$ is given by Eq. (38) and:
$`A^{}`$ $`=`$ $`(\sqrt{{\displaystyle \frac{3(N1)}{2(N+1)}}},\sqrt{{\displaystyle \frac{(N1)(N2)}{2(N+1)(N+2)}}}),`$ (41)
$`D`$ $`=`$ $`(0,\sqrt{{\displaystyle \frac{2(N1)(N2)}{(N+1)(N+2)}}}),`$ (42)
$`E`$ $`=`$ $`(0,\sqrt{{\displaystyle \frac{(N+1)(N1)}{2(N+2)(N2)}}}).`$ (43)
Here, $`A^{}=\vartheta _2A`$ is the image of $`A`$ under $`\vartheta _2`$, while $`D`$ and $`E`$ are the intersections of the lines $`AC`$ and $`AB`$ with the $`\beta _2`$-axis, respectively. The case $`N=4`$ is illustrated in Fig. 1. Similar pictures are obtained for other values of $`N`$. Examples are shown in Fig. 2. Note that the origin of the $`(\beta _1,\beta _2)`$-plane describes the state $`\rho =\frac{1}{3N}II`$ of maximal entropy.
To construct the set $`S_{\mathrm{sep}}^\beta `$ of separable states we have to investigate the functions:
$`\stackrel{~}{\beta }_1[\phi ^{(1)},\phi ^{(2)}]`$
$`=\sqrt{N}{\displaystyle \underset{q=1}{\overset{+1}{}}}\phi ^{(1)}|T_{1q}^{(1)}|\phi ^{(1)}\phi ^{(2)}|T_{1q}^{(2)}|\phi ^{(2)}`$
and
$`\stackrel{~}{\beta }_2[\phi ^{(1)},\phi ^{(2)}]`$
$`=\sqrt{{\displaystyle \frac{3N}{5}}}{\displaystyle \underset{q=2}{\overset{+2}{}}}\phi ^{(1)}|T_{2q}^{(1)}|\phi ^{(1)}\phi ^{(2)}|T_{2q}^{(2)}|\phi ^{(2)}.`$
We distinguish two cases, namely the cases of odd and of even $`N`$.
###### Theorem 1
For integer spins $`j_2=1,2,3,\mathrm{}`$ one has $`S_{\mathrm{ppt}}^\beta =S_{\mathrm{sep}}^\beta `$. Hence, for all $`3N`$ systems with odd $`N`$ the PPT criterion represents a necessary and sufficient condition for the separability of rotationally invariant states.
To proof this theorem we show that the vertices $`A`$, $`A^{}`$, $`D`$ and $`E`$ of the polygon $`S_{\mathrm{ppt}}^\beta `$ belong to $`W^\beta `$. The statement $`S_{\mathrm{ppt}}^\beta =S_{\mathrm{sep}}^\beta `$ then follows immediately from Eq. (24).
The point $`A`$ corresponds to the parameter vector $`๐ถ`$ given by Eq. (31). It follows that this point as well as the point $`A^{}=\vartheta _2A`$ belong to $`W^\beta `$. Hence, it suffices to verify that $`D,EW^\beta `$.
To show that $`EW^\beta `$ we choose the states
$$|\phi ^{(1)}=|1,m_1=0,|\phi ^{(2)}=|j_2,m_2=0.$$
(46)
According to the selection rules for the matrix elements of the tensor operators (78) and to Eq. (85) we have that $`\phi ^{(1)}|T_{1q}^{(1)}|\phi ^{(1)}=0`$ for $`q=0,\pm 1`$ and, therefore,
$$\stackrel{~}{\beta }_1=0.$$
(47)
On the other hand, the non-vanishing matrix elements of the second-rank tensors are given by \[see Eq. (87)\]:
$$\phi ^{(1)}|T_{20}^{(1)}|\phi ^{(1)}=\frac{2}{\sqrt{6}},$$
(48)
and
$$\phi ^{(2)}|T_{20}^{(2)}|\phi ^{(2)}=\frac{2\sqrt{5}j_2(j_2+1)}{\sqrt{(N+2)(N+1)N(N1)(N2)}},$$
(49)
which yields:
$`\stackrel{~}{\beta }_2`$ $`=`$ $`\sqrt{{\displaystyle \frac{3N}{5}}}\phi ^{(1)}|T_{20}^{(1)}|\phi ^{(1)}\phi ^{(2)}|T_{20}^{(2)}|\phi ^{(2)}`$ (50)
$`=`$ $`\sqrt{{\displaystyle \frac{(N+1)(N1)}{2(N+2)(N2)}}}.`$
We see from Eqs. (47), (50) and (43) that $`(\stackrel{~}{\beta }_1,\stackrel{~}{\beta }_2)=E`$ and, hence, that the point $`E`$ belongs to $`W^\beta `$.
To show that also $`D`$ belongs to $`W^\beta `$ we take the states
$$|\phi ^{(1)}=|1,0,|\phi ^{(2)}=|j_2,+j_2.$$
(51)
Since the state $`|\phi ^{(1)}`$ is the same as before, Eqs. (47) and (48) hold true. Instead of Eq. (49), however, we get
$$\phi ^{(2)}|T_{20}^{(2)}|\phi ^{(2)}=\frac{2\sqrt{5}[3j_2^2j_2(j_2+1)]}{\sqrt{(N+2)(N+1)N(N1)(N2)}}.$$
(52)
This gives
$`\stackrel{~}{\beta }_2`$ $`=`$ $`\sqrt{{\displaystyle \frac{3N}{5}}}\phi ^{(1)}|T_{20}^{(1)}|\phi ^{(1)}\phi ^{(2)}|T_{20}^{(2)}|\phi ^{(2)}`$ (53)
$`=`$ $`\sqrt{{\displaystyle \frac{2(N1)(N2)}{(N+1)(N+2)}}}.`$
A comparison with Eq. (42) shows that $`(\stackrel{~}{\beta }_1,\stackrel{~}{\beta }_2)=DW^\beta `$. This concludes the proof of the theorem.
Let us now turn to the case of half-integer spins $`j_2`$, i. e., we assume that $`N`$ is even. Of course, we again have that $`A`$ and $`A^{}`$ belong to $`W^\beta `$. But also $`DW^\beta `$ because the state $`|j_2,+j_2`$ exists for integer as well as for half-integer spins $`j_2`$. The argument following Eq. (51) can thus also be applied in the present case. It follows that $`S_{\mathrm{sep}}^\beta `$ contains at least the triangle $`AA^{}D`$ (see Fig. 3).
On the other hand, the state $`|j_2,m_2=0`$ exists, of course, only for integer spins $`j_2`$. Instead of (46) we consider the states
$$|\phi ^{(1)}=|1,0,|\phi ^{(2)}=|j_2,+1/2,$$
(54)
which lead to
$$\stackrel{~}{\beta }_1=0,\stackrel{~}{\beta }_2=\sqrt{\frac{(N+2)(N2)}{2(N+1)(N1)}}.$$
(55)
This shows that the point
$$F=(0,\sqrt{\frac{(N+2)(N2)}{2(N+1)(N1)}})$$
(56)
belongs to $`W^\beta `$. Hence, $`S_{\mathrm{sep}}^\beta `$ contains at least the polygon with the vertices $`A`$, $`A^{}`$, $`D`$ and $`F`$.
We introduce the straight line $`h`$ which intersects the point $`F`$ and which is parallel to the $`\beta _1`$-axis (see Fig. 3). We are going to demonstrate that $`S_{\mathrm{sep}}^\beta `$ lies entirely below this line. The line $`h`$ is thus tangential to $`S_{\mathrm{sep}}^\beta `$ and corresponds to an optimal entanglement witness (see Sec. V). To show this we employ the rotational invariance of the functions $`\stackrel{~}{\beta }_K`$ \[see Eq. (30)\] to obtain a suitable parametrization of the states of the first spin $`j_1=1`$. Namely, by an appropriate rotation $`R`$ any state of this spin can be brought into the following form:
$$|\phi ^{(1)}=\sqrt{r}|1,+1+\sqrt{1r}|1,1,$$
(57)
where we omit an irrelevant overall phase factor and $`r`$ is a real parameter taken from the interval $`[0,1]`$. Invoking the rotational invariance we may assume without restriction that $`|\phi ^{(1)}`$ is of this form. The state space of the first spin $`j_1`$ has thus only a single relevant parameter $`r[0,1]`$.
By use of the representation (57) the quantities $`\stackrel{~}{\beta }_1`$ and $`\stackrel{~}{\beta }_2`$ become functions of the parameter $`r`$ and of the state vector $`|\phi ^{(2)}`$ of the second spin. Inserting Eq. (57) into Eq. (IV) and using Eqs. (85) and (86) of Appendix A we get
$$\stackrel{~}{\beta }_1[r,\phi ^{(2)}]=\sqrt{\frac{N}{2}}(2r1)\phi ^{(2)}|T_{10}^{(2)}|\phi ^{(2)}.$$
(58)
The function $`\stackrel{~}{\beta }_2`$ is found by substituting the expression (57) into Eq. (IV) and by using Eqs. (87)-(A). One finds that $`\stackrel{~}{\beta }_2`$ can be written as the expectation value
$$\stackrel{~}{\beta }_2[r,\phi ^{(2)}]=\phi ^{(2)}|H(\lambda )|\phi ^{(2)}$$
(59)
of the Hermitian $`(N\times N)`$ matrix
$$H(\lambda )H_0+\lambda H_1.$$
(60)
Here, we have defined
$$H_0=\sqrt{\frac{N}{10}}T_{20}^{(2)},H_1=\frac{1}{2}\sqrt{\frac{3N}{5}}\left(T_{22}^{(2)}+T_{22}^{(2)}\right),$$
and introduced the parameter
$$\lambda =2\sqrt{r(1r)},0\lambda 1.$$
(61)
For a given value of $`\lambda `$ the function $`\stackrel{~}{\beta }_2`$ defined by Eq. (59) is certainly smaller than or equal to the largest eigenvalue of $`H(\lambda )`$ which we denote by $`\epsilon _0(\lambda )`$. We are going to demonstrate below that $`\epsilon _0(\lambda )`$ is a monotonically increasing function of $`\lambda `$ and attains its maximum at $`\lambda =1`$:
$$\epsilon _0(1)=\sqrt{\frac{(N+2)(N2)}{2(N+1)(N1)}}.$$
(62)
Hence, we have
$$\stackrel{~}{\beta }_2[r,\phi ^{(2)}]\epsilon _0(1)$$
(63)
for all $`r`$ and $`|\phi ^{(2)}`$. Note that $`\epsilon _0(1)`$ is equal to the $`\beta _2`$-coordinate of the point $`F`$ \[see Eq. (56)\]. This shows that, as claimed, all points of $`W^\beta `$ and, hence, all points of $`S_{\mathrm{sep}}^\beta `$ lie below the line $`h`$.
To prove that $`\epsilon _0(\lambda )`$ is a monotonically increasing function of $`\lambda `$ we denote the eigenvalues of $`H(\lambda )`$ by $`\epsilon _n(\lambda )`$, where $`n=0,1,2\mathrm{}`$, and $`n=0`$ labels the largest eigenvalue. With the help of Eq. (83) one verifies that $`H(\lambda )`$ is invariant under time reversal. It follows that if $`|\phi `$ is an eigenstate of $`H(\lambda )`$ then also the time reversed state $`\tau |\phi `$ is an eigenstate with the same eigenvalue. Since $`j_2`$ is half-integer valued the states $`|\phi `$ and $`\tau |\phi `$ are orthogonal. In fact, using the antiunitarity of $`\tau `$ and Eq. (26) we get
$$\tau \phi |\phi =\tau ^2\phi |\tau \phi ^{}=(1)^{2j_2}\tau \phi |\phi =\tau \phi |\phi ,$$
which shows that $`\tau \phi |\phi =0`$.
All eigenvalues $`\epsilon _n(\lambda )`$ are thus two-fold degenerate and we write the corresponding eigenstates as $`|\phi _{n,k}(\lambda )`$, where the index $`k=1,2`$ labels the eigenstates corresponding to the same eigenvalue: $`|\phi _{n,2}(\lambda )=\tau |\phi _{n,1}(\lambda )`$. We remark that the two-fold degeneracy is analogous to the Kramers degeneracy according to which the energy levels of an invariant system of an odd number of spin-$`\frac{1}{2}`$ particles are at least two-fold degenerate (see, e. g. SAKURAI ).
The Hellman-Feynman theorem now yields
$$\frac{d\epsilon _0}{d\lambda }=\phi _{0,1}(\lambda )|H_1|\phi _{0,1}(\lambda ).$$
(64)
In particular, we have
$$\frac{d\epsilon _0}{d\lambda }|_{\lambda =0}=0.$$
(65)
On differentiating Eq. (64) once again we find:
$$\frac{d^2\epsilon _0}{d\lambda ^2}=2\underset{n0,k}{}\frac{|\phi _{n,k}(\lambda )|H_1|\phi _{0,1}(\lambda )|^2}{\epsilon _0(\lambda )\epsilon _n(\lambda )}0.$$
(66)
This shows that $`\epsilon _0(\lambda )`$ is a convex function of $`\lambda `$ with zero derivative at $`\lambda =0`$. It follows that $`\epsilon _0(\lambda )`$ increases monotonically. Some examples of the behaviour of this functions are shown in Fig. 4.
It remains to verify Eq. (62). We first note that $`H(1)`$ can be written with the help of Eqs. (87) and (A) in terms of the spin operator $`\widehat{๐}^{(2)}`$ as:
$`H(1)`$ $`=`$ $`2\sqrt{{\displaystyle \frac{2}{(N+2)(N+1)(N1)(N2)}}}`$ (67)
$`\times \left(\left[\widehat{๐}^{(2)}\right]^23\left[\widehat{j}_y^{(2)}\right]^2\right).`$
The largest eigenvalue of this matrix is given by
$`\epsilon _0(1)`$ $`=`$ $`2\sqrt{{\displaystyle \frac{2}{(N+2)(N+1)(N1)(N2)}}}`$ (68)
$`\times \left(j_2(j_2+1){\displaystyle \frac{3}{4}}\right).`$
Using $`N=2j_2+1`$ one shows that this equation coincides with Eq. (62).
We finally demonstrate that the boundary of $`S_{\mathrm{sep}}^\beta `$ is differentiable at the point $`F`$ \[see Eq. (56)\]. To this end, we construct a smooth curve which belongs to $`W^\beta `$ and passes the point $`F`$. Consider the following fixed state of the second spin:
$$|\phi ^{(2)}=\frac{1}{\sqrt{2}}|\widehat{j}_y^{(2)}=+1/2+\frac{i}{\sqrt{2}}|\widehat{j}_y^{(2)}=1/2.$$
(69)
This is an eigenstate of the matrix $`H(1)`$ \[Eq. (67)\] corresponding to the largest eigenvalue $`\epsilon _0(1)`$. Since $`|\phi ^{(2)}`$ is fixed the functions $`\stackrel{~}{\beta }_1`$ and $`\stackrel{~}{\beta }_2`$ depend only on the parameter $`r`$ and describe a curve in the $`(\beta _1,\beta _2)`$-plane. Writing $`r(1+\mu )/2`$ and determining the matrix elements one finds:
$`\stackrel{~}{\beta }_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{3N^2}{8(N+1)(N1)}}}\mu ,`$ (70)
$`\stackrel{~}{\beta }_2`$ $`=`$ $`{\displaystyle \frac{\epsilon _0(1)}{4}}\left(1+3\sqrt{1\mu ^2}\right),`$ (71)
where $`1\mu +1`$. The curve described by these equations represents the upper half of an ellipse in the $`(\beta _1,\beta _2)`$-plane (see Fig. 3). It intersects the point $`F`$ and lies entirely in $`W^\beta `$. Since $`F`$ is the only point of $`h`$ belonging to $`W^\beta `$, it follows that the boundary of the separability region must be curved and that it is differentiable at the extreme point $`F`$, the line $`h`$ being the tangent. Summarizing, we have shown:
###### Theorem 2
For half-integer spins $`j_2=\frac{3}{2},\frac{5}{2},\frac{7}{2},\mathrm{}`$ the set $`S_{\mathrm{sep}}^\beta `$ of separable states is a true subset of the set of PPT states. Hence, for all $`3N`$ systems with even $`N`$ the PPT criterion is only necessary and there always exist entangled PPT states. The line $`h`$ represents the tangent to $`S_{\mathrm{sep}}^\beta `$ at the extreme point $`F`$. The set $`S_{\mathrm{sep}}^\beta `$ is bounded by the straight lines $`AD`$ and $`A^{}D`$ and by a concave curve which passes the points $`A`$, $`A^{}`$ and $`F`$.
## V Discussion and conclusions
The state space structure of rotationally invariant spin systems has been analyzed in this paper. The set of invariant states has been represented by means of two systems of invariant operators, namely by the projections $`P_J`$ onto the total angular momentum manifolds and by the invariant operators $`Q_K`$ composed of the spherical tensors. The transformation between both representations was found to be given by a matrix $`L`$ which is determined by certain $`6`$-$`j`$ symbols of Wigner. The $`Q_K`$-representation is particularly useful in applying the PPT criterion for separability because the $`Q_K`$ are eigenoperators of the partial time reversal. The method has been demonstrated to lead to a complete classification of separability of $`3N`$ systems. We have shown that the PPT criterion is necessary and sufficient for all system with odd $`N`$, while entangled PPT states exist for systems with even $`N`$.
Some remarks on the structure of the state space in the limit $`N\mathrm{}`$ might be of interest. In this limit the value of the second spin $`j_2`$ becomes arbitrary large. We infer from Eqs. (39)-(42) that the point $`B`$ then converges to the point $`A^{}`$, and $`C`$ to $`D`$. At the same time $`F`$ converges to $`E`$ \[see Eqs. (43) and (56)\]. Hence, as $`N`$ increases the set $`S_{\mathrm{ppt}}^\beta `$ approaches the set $`S^\beta `$ and $`S_{\mathrm{sep}}^\beta `$ approaches $`S_{\mathrm{ppt}}^\beta `$. This behaviour is also indicated in Fig. 2. The limit $`N\mathrm{}`$ thus corresponds to a kind of classical limit in which all invariant states have a positive partial transpose and are separable.
The line $`h`$ constructed in Sec. IV leads to an entanglement witness which we denote by $`๐ฒ`$. An entanglement witness is a Hermitian operator which satisfies $`\mathrm{tr}\{๐ฒ\sigma \}0`$ for any separable state $`\sigma `$, and $`\mathrm{tr}\{๐ฒ\rho \}<0`$ for at least one non-separable state $`\rho `$ HORODECKI96a ; TERHAL . The hyperplane $`h`$ corresponding to an entanglement witness $`๐ฒ`$ is defined by $`\mathrm{tr}\{๐ฒ\rho \}=0`$. In the case of $`3N`$ systems $`h`$ is a one-dimensional line and the witness is, in fact, optimal LEWENSTEIN00 because $`h`$ is tangential to the region of separable states. We have formulated the witness in $`\beta `$-space. Transforming back to $`\alpha `$-space one easily shows that the entanglement witness corresponding to $`h`$ may be written in terms of the projections $`P_J`$ as:
$$๐ฒ=\frac{1}{N2}P_{j_21}+P_{j_2}+\frac{1}{N+2}P_{j_2+1}.$$
(72)
This expression leads to the following physical interpretation of $`๐ฒ`$. Suppose one carries out a measurement of the total angular momentum $`J`$ on some invariant state $`\rho `$. If $`\rho `$ is separable the inequality
$$\frac{p_{j_21}}{N2}+p_{j_2}+\frac{p_{j_2+1}}{N+2}0$$
(73)
must be satisfied, where $`p_J=\mathrm{tr}\{P_J\rho \}`$ denotes the probability of finding the value $`J`$. In other words, if the inequality (73) is violated the state $`\rho `$ must necessarily be entangled.
We exploit the witness (72) to design a prescription for the detection of entangled PPT states in $`3N`$ systems with even $`N`$ (bound entanglement HORODECKI98 ). A given state $`\rho `$ is positive under partial transposition if and only if the corresponding point $`(\beta _1,\beta _2)`$ lies below the line through $`A^{}`$ and $`E`$, and above the line through $`A^{}`$ and $`D`$ (see Fig. 1). If we transform to $`\alpha `$-space this yields the conditions
$$\frac{2p_{j_21}}{N1}+\frac{(N^25)p_{j_2}}{(N+1)(N1)}+\frac{2p_{j_2+1}}{N+1}0$$
(74)
and
$$\frac{2p_{j_21}}{(N1)(N2)}\frac{2p_{j_2}}{N1}+p_{j_2+1}0.$$
(75)
These inequalities are equivalent to the PPT condition (27). Hence, entangled PPT states can be detected in the following way: Suppose again that a total angular momentum measurement is performed on some state $`\rho `$. If one finds that the measurement outcomes, i. e. the probabilities $`p_J`$, satisfy the inequalities (74) and (75) and violate the inequality (73) then the state $`\rho `$ must be an entangled PPT state.
The witness $`๐ฒ`$ defined in Eq. (72) does not detect all entangled PPT states. As has been shown in Sec. IV a part of the boundary of the region of the separable states is curved and, therefore, one needs an infinite number of linear entanglement witnesses. The upper boundary of $`S_{\mathrm{sep}}^\beta `$ can of course be described by means of a suitable nonlinear equation. A possible way to derive the latter is to construct the envelope of appropriate families of curves of the type given by Eqs. (70) and (71).
The considerations of Sec. IV reveal that for $`3N_2`$ systems half-integer spins are crucial for the emergence of entangled PPT states. The entanglement structure of systems involving half-integer spins is thus quite different from those with integer spins. It seems that this is closely connected to the fact that pure states which are invariant under time reversal only exist for integer spins, while for half-integer spins a given pure state is always orthogonal to its time reversed state. A clear physical interpretation of this result and its generalization to arbitrary $`N_1N_2`$ systems is of great interest. The next step to further investigate this point could be to study $`4N_2`$ systems, which is possible by the method developed in this paper.
###### Acknowledgements.
The author would like to thank J. Schliemann and F. Petruccione for helpful comments and stimulating discussions.
## Appendix A Spherical tensor operators
We define here the irreducible spherical tensor operators $`T_{Kq}`$ acting on the state space $`^N`$ of a particle with spin $`j`$, where $`N=2j+1`$, $`K=0,1,\mathrm{},2j`$, and $`q=K,\mathrm{},+K`$. The tensor operators $`T_{K_iq_i}^{(i)}`$ for $`i=1,2`$ used in the main text are obtained by setting $`j=j_1`$ or $`j=j_2`$.
The spherical tensor operators $`T_{Kq}`$ represent a complete system of operators on $`^N`$. This means that any operator on the state space of the spin-$`j`$ particle may be written as a unique linear combination of the $`T_{Kq}`$. Moreover, the tensors are orthonormal with respect to the Hilbert-Schmidt inner product:
$$\mathrm{tr}\left\{T_{K^{}q^{}}^{}T_{Kq}\right\}=\delta _{KK^{}}\delta _{qq^{}}.$$
(76)
For a fixed $`K`$ the $`(2K+1)`$ operators $`T_{Kq}`$ represent the spherical components of a tensor of rank $`K`$. They transform according to an irreducible representation of $`\mathrm{SO}(3)`$ which corresponds to the angular momentum $`K`$:
$$D^{(j)}(R)T_{Kq}D^{(j)}(R)^{}=\underset{q^{}=K}{\overset{+K}{}}D_{q^{}q}^{(K)}(R)T_{Kq^{}}.$$
(77)
For instance, the $`T_{1q}`$ behave as components of a vector, and the $`T_{2q}`$ as components of a second-rank tensor.
The matrix elements of the tensors may be defined in term of Wignerโs $`3`$-$`j`$ symbols as WIGNER ; EDMONDS
$$j,m|T_{Kq}|j,m^{}=\sqrt{2K+1}(1)^{jm}\left(\begin{array}{ccc}j& j& K\\ m& m^{}& q\end{array}\right).$$
(78)
The $`3`$-$`j`$ symbols are closely related to the Clebsch-Gordan coefficients:
$`j_1,m_1;j_2,m_2|JM=`$ (81)
$`\sqrt{2J+1}(1)^{j_1j_2+M}\left(\begin{array}{ccc}j_1& j_2& J\\ m_1& m_2& M\end{array}\right).`$
According to the selection rules for the $`3`$-$`j`$ symbols the matrix element (78) is equal to zero for $`mm^{}q0`$. In particular, we have $`T_{00}=\frac{1}{\sqrt{N}}I`$.
The matrix elements (78) of the tensor operators are real and one has $`T_{Kq}^{}=T_{Kq}^T=(1)^qT_{K,q}`$. It follows that the $`T_{Kq}`$ are eigenoperators of the time reversal transformation $`\vartheta `$ which was defined in Eq. (25). In fact, using the transformation behaviour (77) of the tensors and the fact that a rotation by $`\pi `$ about the $`y`$-axis is represented by the unitary matrix
$$D_{q^{}q}^{(K)}(\pi )=(1)^{Kq^{}}\delta _{q^{},q},$$
(82)
one finds
$$\vartheta T_{Kq}=VT_{Kq}^TV^{}=(1)^KT_{Kq}.$$
(83)
As a consequence the operators $`Q_K`$ which have been introduced in Eq. (7) are eigenoperators of the partial time reversal $`\vartheta _2=I\vartheta `$:
$$\vartheta _2Q_K=(1)^KQ_K.$$
(84)
We finally list the non-vanishing matrix elements of the tensor operators needed in Sec. IV:
$$j,m|T_{10}|j,m=2m\sqrt{\frac{3}{N(N1)(N+1)}},$$
(85)
$$j,m|T_{11}^{}|j,m+1=\sqrt{\frac{6(jm)(j+m+1)}{N(N1)(N+1)}},$$
(86)
$$j,m|T_{20}|j,m=\frac{2\sqrt{5}[3m^2j(j+1)]}{\sqrt{(N+2)(N+1)N(N1)(N2)}},$$
(87)
$`j,m|T_{21}^{}|j,m+1=`$
$`\sqrt{5}(1+2m)\sqrt{{\displaystyle \frac{6(jm)(j+m+1)}{(N+2)(N+1)N(N1)(N2)}}},`$
$`j,m|T_{22}^{}|j,m+2=`$
$`\sqrt{5}\sqrt{{\displaystyle \frac{6(jm1)(jm)(j+m+1)(j+m+2)}{(N+2)(N+1)N(N1)(N2)}}}.`$
## Appendix B Proof of relation (14)
The starting point is given by Eq. (13). We insert into this equation the definitions (2) and (7) for the invariant operators $`P_J`$ and $`Q_K`$, and introduce complete sets of product basis states $`|j_1,m_1;j_2,m_2`$. This yields a multiple sum over products of two Clebsch-Gordan coefficients and two matrix elements of the tensor operators. By use of Eqs. (78) and (81) the Clebsch-Gordan coefficients as well as the matrix elements of the spherical tensors can be written in terms of the $`3`$-$`j`$ symbols. We also use the selection rules for the $`3`$-$`j`$ symbols and their symmetry properties. This procedure leads to the following sum over $`4`$-fold products of $`3`$-$`j`$ symbols:
$`L_{KJ}`$ $`=`$ $`\sqrt{(2K+1)(2J+1)}(1)^{j_1+j_2+J}\times `$
$`{\displaystyle \underset{\{m_i\}}{}}\chi (\{m_i\})\times `$
$`\left(\begin{array}{ccc}j_1& j_2& J\\ m_1& m_2& m_3\end{array}\right)\left(\begin{array}{ccc}j_1& j_1& K\\ m_1& m_5& m_6\end{array}\right)\times `$ (95)
$`\left(\begin{array}{ccc}j_2& j_2& K\\ m_4& m_2& m_6\end{array}\right)\left(\begin{array}{ccc}j_2& j_1& J\\ m_4& m_5& m_3\end{array}\right),`$ (100)
where $`\chi (\{m_i\})`$ is a phase factor:
$`\chi (\{m_i\})`$ $`=`$ $`(1)^{j_1+m_1}(1)^{j_2+m_2}(1)^{J+m_3}\times `$
$`(1)^{j_2+m_4}(1)^{j_1+m_5}(1)^{K+m_6}.`$
The sum over the quantum numbers $`m_1,\mathrm{},m_6`$ in Eq. (B) exactly corresponds to a certain $`6`$-$`j`$ symbol of Wigner WIGNER . A general $`6`$-$`j`$ symbol involves six angular momenta and is written as
$$\left\{\begin{array}{ccc}j_1& j_2& j_3\\ j_4& j_5& j_6\end{array}\right\}.$$
(101)
The sum of Eq. (B) is equal to the $`6`$-$`j`$ symbol (101) with $`j_3=J`$, $`j_4=j_2`$, $`j_5=j_1`$ and $`j_6=K`$. Hence, we see that Eq. (B) reduces to Eq. (14). We remark that a similar technique has been used in Ref. NtensorN in order to derive an expression for the matrix which represents the partial time reversal $`\vartheta _2`$ in the $`P_J`$-representation.
By use of the formulae for the $`6`$-$`j`$ symbols EDMONDS we find that the first three rows of $`L`$ are given by
$$L_{0J}=\sqrt{\frac{2J+1}{N_1N_2}},$$
(102)
and
$$L_{1J}=2\sqrt{3(2J+1)}\frac{j_1(j_1+1)+j_2(j_2+1)J(J+1)}{\sqrt{(N_11)N_1(N_1+1)(N_21)N_2(N_2+1)}},$$
(103)
$$L_{2J}=2\sqrt{5(2J+1)}\frac{3X(X1)4j_1(j_1+1)j_2(j_2+1)}{\sqrt{(N_12)(N_11)N_1(N_1+1)(N_1+2)(N_22)(N_21)N_2(N_2+1)(N_2+2)}},$$
(104)
where $`Xj_1(j_1+1)+j_2(j_2+1)J(J+1)`$. |
warning/0506/physics0506152.html | ar5iv | text | # Determination of the chemical potential using energy-biased sampling
## I Introduction
The chemical potential is a central quantity underpinning many physical and chemical processes, such as phase equilibria, osmosis, thermodynamic stability, binging affinity and so on Lu et al. (2003). However, its evaluation by computer simulation is more complicated and time-consuming than for other intensive thermodynamic quantities, such as the pressure $`P`$ or temperature $`T`$. While $`P`$ and $`T`$ can be evaluated from averages over mechanical properties of molecules (forces, velocities and positions), the chemical potential is a thermal average and therefore it requires sampling the phase space of the system. Indeed, computing the chemical potential is a special case of the more general problem of computing a free-energy difference $`A_1A_0`$ between two states (labelled as 0 and 1), a problem for which the inherent difficulty is well understood Allen and Tildesley (1987); Frenkel and Smith (2002); Kollman (1993); Lu et al. (2003). Free energy perturbation (FEP) is an important category of methods for free energy calculation; we refer to the recent works by Lu et al. Lu et al. (2003) and by Shirts and Pande Shirts and Pande (2005) for review and comparisons. As explained by Lu et al. Lu et al. (2003), the general working equation for FEP methods can be cast as
$$\mathrm{exp}[\beta (A_1A_0)]=\frac{w(u)\mathrm{exp}[\beta u/2]_0}{w(u)\mathrm{exp}[\beta u/2]_1},$$
(1)
with $`\beta =1/k_BT`$ and $`uU_1U_0`$ the energy difference between both systems; $`K_B`$ is the Boltzmann constant. The angular brackets denote ensemble averages performed on the system labelled by the subscript โ0โ or โ1โ. The weighting function $`w(u)`$ is arbitrary and differs for each method introduced in the literature.
The chemical potential is the free energy difference between two thermodynamic states differing by the presence of a single molecule. In other words, the chemical potential is $`A_1A_0`$ where $`A_1=A(N+1,V,T)`$ and $`A_0=A(N,V,T)`$. Here $`A(N,V,T)`$ is the Helmholtz free energy of the system which depends on the number of molecules N, the volume $`V`$ and temperature of the system. In order to express the averages of Eq. (1) in terms of one-dimensional integrals of the energy difference $`u`$ one can then introduce the following distribution functions Deitrick et al. (1989)
$`f(u)`$ $`=`$ $`{\displaystyle \delta \left(uU_1+U_0\right)_0V^1๐๐ซ},`$ (2)
$`g(u)`$ $`=`$ $`\delta \left(uU_1+U_0\right)_1,`$ (3)
where $`\delta (.)`$ is the Dirac delta function. In Eq. (2), $`U_1=U_1(๐^N,๐ซ)`$, where $`๐^N`$ is the configuration of the first N molecules and $`๐ซ`$ denotes the configuration of the N+1 molecule. Note that in Eq. (2) the N+1 molecule acts as a โtest-moleculeโ which probes the system โ0โ (i.e. the system with N molecules), but does no interact with it. Therefore $`f(u)`$ is the probability density of the N molecule ensemble increasing in potential energy by an amount $`u`$ if this test-molecule were randomly inserted into the ensemble. Conversely, $`g(u)`$ is the probability density of the (N+1)-molecule ensemble decreasing in potential energy by an amount $`u`$ if a randomly selected real molecule were removed from the ensemble.
ยฟFrom Eq. (1)-(3) an expression for the excess chemical potential $`\mu =A_1A_0\mu _{id}`$ (where $`\mu _{id}`$ is the ideal gas chemical potential Frenkel and Smith (2002)) can be derived in terms of the $`f`$ and $`g`$ distributions Shing and Gubbins (1982); Deitrick et al. (1989); Lu et al. (2003)
$$\mathrm{exp}(\beta \mu )=\frac{w(u)g(u)๐u}{w(u)f(u)\mathrm{exp}(\beta u)๐u}.$$
(4)
A good choice of the weighting function $`w(u)`$ is key for the efficiency of the method. For instance, the Widom method Frenkel and Smith (2002); Allen and Tildesley (1987) ($`w(u)=1`$) is known to provide very poor convergence at large densities. The Widom method is a single stage FEP, meaning that sampling is only performed in the reference system โ0โ (i.e., in the $`f`$ distribution, see Eq. (4)). As discussed by Lu et al. Lu et al. (2003), multiple staging provides much better efficiency. The efficiency is generally defined as the reciprocal of the product of the variance of the estimator multiplied by its cost $`n_{cost}`$ (that is, the total number of energy evaluations performed by the algorithm)
$$\epsilon =(n_{cost}\mathrm{๐
๐๐}[\beta \mu ])^1.$$
(5)
Bennett Bennett (1976) showed that the variance of Eq. (4) is minimised if the weighting function is $`w(u)=[\beta (uc)]`$, where $`(x)=1/(1+\mathrm{exp}(x))`$ is the Fermi function and $`c`$ is an arbitrary constant. The Bennett estimator is then
$$\beta \mu =\mathrm{ln}\left(\frac{[\beta (uc)]_g}{[\beta (uc)]_f}\right)+\beta c,$$
(6)
where the subscripts $`g`$ and $`f`$ indicate (simple) averages over the distributions $`g(u)`$ and $`f(u)`$. The value of $`c`$ providing the minimum variance and maximum overlap is $`c=\mu `$ and to evaluate $`\mu `$ using the optimum $`c(=\mu )`$ one requires to use a self-consistent procedure, iterating the value of $`c`$ in Eq. (6) and resetting $`c=\mu `$ until $`[\beta (uc)]_g=[\beta (uc)]_f`$. In practise, this step only requires a small number of iterations. Recent publications Lu et al. (2003); Shirts and Pande (2005) demonstrate that the Bennett method remains the best general method to compute the chemical potential for many applications.
Note that the Bennett method is a two-stage FEP and therefore it also requires sampling of the system โ1โ. In the case of the determination of the chemical potential this system has N+1 molecules and $`g(u)`$ is obtained from its single-molecule energy distribution. However this extra requirement is not really a drawback. Lu et al. Lu et al. (2003) showed that, provided $`N>O(100)`$, the $`g`$average can be evaluated in the same simulation as is used to sample the $`f`$ distribution (system โ0โ) without any noticeable loss in accuracy. The $`g`$ distribution (constructed from the energy of the real particles) is thus a byproduct of the simulation so the average $`_g`$ does not demand any extra computational cost.
Another group of methods for determination of the chemical potential are based on biased instead of uniform sampling. In particular, cavity-biased methods first select spherical cavities of minimum radius $`R_c`$ (a free parameter) in which to insert the test-molecule. This accelerates the evaluation of the ensemble average in dense phases because the low-energy configurations of the test-molecule (with large Boltzmann factors) are usually located in larger cavities with less steric hindrance. Variations of this method have been proposed by several authors; these include the Cavity Insertion Widom method (CIW) due to Mezei and coworkers Jedlovszky and Mezei (2000), the Excluded Volume Map Sampling by Deitrick et al. Deitrick et al. (1989) and the method proposed by Pohorille and Wilson Pohorille and Wilson (1996). The cavities are located by a grid search over the whole simulation cell. A cavity centre is assigned at each grid point whose distance to the closest particle is greater than $`R_c`$. In order to correct the bias introduced in sampling only inside the cavities one also has to calculate the probability of finding a cavity, which is obtained in the same grid-search step. A drawback of the cavity-biased method is that it is only indirectly related to the test-particle energy via the excluded volume. This fact introduces a certain inaccuracy in the estimation of the chemical potential, as it can depend on the value of the cavity radius $`R_c`$ selected. For instance, the CIW has recently been used to calculate the chemical potential of several species across a lipid bilayer Jedlovszky and Mezei (2000). As a test calculation the authors estimated the chemical potential of water in water and reported variations of about 1 Kcal/mol as $`R_c`$ was varied from $`2.6\AA `$ to $`2.8\AA `$. Also, using $`R_c[2.6,2.9]\AA `$ resulted in uncertainties of about 2 Kcal/mol in estimates of the excess chemical potential of some species across the lipid layer. Note that the important region of the cavity-biased method is constructed over the translational degrees of freedom of a โcoarse-grainedโ spherical molecule with an effective radius. This means that it can only be applied to small solutes with spherical or roughly spherical shapes Deitrick et al. (1989).
In this work we present an energy-biased method for the estimation of the chemical potential and reconstruction of the energy distribution $`f(u)`$ in dense phases. The idea is to restrict the sample to an important region defined by the set of bounded domains in the configurational space of the test-molecule where the energy $`u`$ is smaller than a given free parameter $`u_w`$. We denote as an energy-well each compact subdomain within the test-molecule energy-landscape for which $`u<u_w`$. Note that the present approach retains the main benefit of the cavity-biased method, but provides an exact evaluation of the energy distribution $`f(u)`$ and the chemical potential, because the energy-wells are defined directly in terms of the energy landscape. Moreover our energy-biased method does not assume any particular molecular shape and therefore it may be used for non-spherical molecules and can coherently sample over rotational degrees of freedom as well.
We also note that the number of stages are not limited to two. When systems 0 and 1 are very different it may be impossible within the simulation time to sample the importance region of the two systems. In this case it is more efficient to compute the total free energy difference by using a set of intermediate states. The energy bias method can be applied on each of these intermediate state transitions at the cost of performing independent simulations for each state. Other approaches include, for instance, slow and fast growth methods where the system is changed from one state to another within a certain simulation time $`\tau `$ (large for slow growth). The fast growth method consists of sampling rapid transformation from many simulations which are then combined by using Jarzynski nonequilibrium work relation Jarynski (1997) to obtain the total free energy difference.
The rest of the paper proceeds as follows. The energy-biased method is explained in Sec. II, while in Sec. III we derive an analytical expression for the efficiency of the method and estimate the optimal parameter $`u_w`$ by maximising the efficiency. In Sec IV the method is tested in liquid argon at high density (modelled as Lennard-Jones atoms) where it is used to reconstruct the test-particle energy distribution $`f(u)`$ and the chemical potential. We also demonstrate the gain in efficiency obtained with energy-biased sampling with respect to uniform sampling. We conclude with a summary of our findings in Sec. V. Finally in Appendix A we briefly explain the Hit&Run algorithm which efficiently samples bounded regions of arbitrary shape immersed in an arbitrary number of dimensions.
## II Overview of the method
As stated in the introduction, energy-biased sampling consists of uniform sampling of the importance region defined by the set of subdomains in the test-molecule configurational space where its potential energy is less than $`u_w`$. The probability density is therefore given by
$$h(u)=\{\begin{array}{cc}f(u)/F_w& uu_w\\ 0& u>u_w,\end{array}$$
(7)
where the normalisation factor $`F_w_{\mathrm{}}^{u_w}f(u)๐u`$ is the cumulative probability of the unbiased distribution $`f(u)`$ and $`u_w`$ is an arbitrary energy (free parameter).
Note that the energy-biased distribution of Eq. (7) can be straightforwardly combined with any of the popular methods to calculate the chemical potential from Eq. (4). We shall use the Bennett method due to its excellent performance. Introducing the weighting function $`w(u)=[\beta (cu)]`$ in Eq. (6) and using Eq. (7), one obtains the energy-biased Bennett estimator for $`\beta \mu `$,
$$\beta \mu =\mathrm{ln}\left(\frac{_c_g}{F_w_c_h}\right)+\beta c,$$
(8)
where we have introduced the notation $`_c[\beta (uc)]`$ to indicate that after the ensemble average we still have a function of $`c`$. As before, the subscript $`h`$ indicates the average over the biased distribution of Eq. (7).
Sampling from the energy probability distribution $`h(u)`$ requires a more careful consideration of the energy landscape of the system. We indicate by $`๐ซ`$ a configuration of the (N+1)th molecule and by $`๐`$ the configuration of the remaining N molecules. For a simple argon fluid $`๐ซD`$ where $`DR^3`$, while for a 3 sites flexible water model like TIP3P $`DR^9`$, which includes the three Euler angles determining the molecule orientation, the H-O-H angle and the two H-O distances.
As shown in Fig. (1), the region
$$A_{u_w}=\{๐ซD:u(๐ซ,๐)<u_w\}$$
(9)
is composed of many disconnected bounded regions of different sizes such that $`A_{u_w}=_\alpha A_{u_w}^\alpha `$, where each $`A_{u_w}^\alpha `$ is now a connected region. Of course, for $`u_w\mathrm{}`$ we have that all the regions $`A_{u_w}^\alpha `$ connect and $`A_{\mathrm{}}^\alpha =D`$, the entire domain. The sampling algorithm must reproduce a uniform probability distribution
$$p_{u_w}(๐ซ)=\frac{1}{\mathrm{\Omega }(A_{u_w})},$$
(10)
where $`\mathrm{\Omega }(A_{u_w})`$ is the volume of the region.
For a given energy bias $`u_w`$, the algorithm for selecting configurations $`๐ซ`$ according to Eq. (10) can be described in terms of two main steps which are applied iteratively:
1. Locate a compact energy-well $`A_{u_w}^\alpha `$ in the configurational space D, where $`u<u_w`$.
2. Sample the energy-well $`A_{u_w}^\alpha `$ with a uniform probability density.
The simplest procedure for locating energy wells in step (1) is to perform a random search over the whole configurational space until a fixed number of cavities is found. This procedure, however, does not avoid the probability of exploring the same well more than once, and we observed that it can easily lead to highly correlated data. Instead we perform step (1) by choosing points on a grid within the whole configurational space of the test-molecule. In the case of the Lennard-Jones fluid, the three-dimensional configurational space is probed at the nodes of a Cartesian grid of size $`n_x\times n_y\times n_z`$, where $`n_\alpha `$ is the number of nodes along the coordinate $`\alpha `$. We observed that the minimum distance between nodes that guarantees statistically independent samples is around $`0.5\sigma `$.
An energy well is found at each node where the energy of the test-molecule is $`u<u_w`$. Then, the locations of each of these nodes are used as starting configurations for independent well samplings. In this way we ensure that we are sampling different cavities for each explored configuration (snapshot) of the system. Note that using grid-sampling the number of cavities found per snapshot is a fluctuating quantity.
The search requires an average of $`n_0=1/F_w`$ energy evaluations to locate one well (i.e. one configuration with energy $`u<u_w`$.) During this same step (1) one can calculate the cumulative probability $`F_w`$ from the estimator $`m/n_0`$, with $`n_0`$ being the total number of samples (Bernoulli trials) and $`m`$ the number of successful trials with $`u<u_w`$, i.e., the total number of energy-wells found. This number $`m/n_0`$ converges to $`F_w`$ as $`n_0\mathrm{}`$ and, for a finite number of statistically independent trials $`n_0`$, its variance is $`(1F_w)F_w/n_0`$. In practise, the estimation of $`F_w`$ requires the number of unbiased samples to be $`n_0>>1/F_w`$; this condition also ensures that a significant number of energy-wells ($`m>0`$) are to be found.
Step (2) of the loop mentioned above requires a procedure to sample in an unbiased way the interior of each energy well. This is a delicate step because any bias incurred in sampling the importance region will be transfered to the estimator for $`\beta \mu `$, resulting in inaccuracy of the method. To tackle this problem we use the so-called Hit&Run algorithm Smith (1984), which is explained in Appendix A.
## III Efficiency and optimal parameters of the method
We now calculate the efficiency of the method and provide a way of choosing the optimal value of the parameter $`u_w`$ by maximising the efficiency. We also compare the efficiency of the estimator in Eq. (8) based on energy-biased sampling with that of the standard Bennett algorithm of Eq. (6).
### III.1 Energy-biased Bennett method
The variance of the Bennett method can be cast in terms of the probability densities $`f(u)`$ and $`g(u)`$. Starting from Eq. (6), after some algebra the variance of the Bennett method assumes the form
$$\mathrm{๐
๐๐}_B[\beta \mu ]=\frac{1}{n_0[\beta (uc)]_f},$$
(11)
where $`n_0`$ is the number of insertions used to sample the complete configurational space of the test-particle. Note that the computational cost of the standard Bennett method is $`n_0`$, so according to (5) and Eq. (11) its maximum efficiency is given by
$$\epsilon _B=_c_f.$$
(12)
Let us now consider the variance of the estimator in Eq. (8), which is the sum of the variance of the estimator for $`F_w`$ and the estimator for the ensemble average
$$\mathrm{๐
๐๐}_{EB}[\beta \mu ]=\mathrm{๐
๐๐}[\mathrm{ln}F_w]+\frac{1}{n_w_c_h}\frac{1}{n_0F_w}+\frac{1}{n_w_c_h},$$
(13)
where we have used the relation $`\mathrm{๐
๐๐}[\mathrm{ln}(F_w)]\mathrm{๐
๐๐}[F_w]/F_w^2=(1F_w)/(n_0F_w)1/(n_0F_w)`$, for $`F_w<<1`$. Here $`n_0`$ is the number of random insertions in the entire configurational space and $`n_w`$ is the number of independent samples within the importance region $`u<u_w`$.
The probability of finding an energy-well with $`u<u_w`$ using uniform sampling over the whole configurational space is $`F_w`$, so the number of cavities found after $`n_0`$ trials is $`m=F_wn_0`$. If the number of statistically independent samples per well is $`s`$, the total number of independent samples within the restricted configurational space $`u<u_w`$ is
$$n_w=n_0sF_w.$$
(14)
We note that the number of independent samples per well $`s`$ depends on the fluid considered and, of course, on the biasing energy $`u_w`$. In Appendix B we provide a way of estimating $`s`$ from the outcome of the data obtained from Hit&Run sampling. Inserting Eq. (14) into Eq. (13) one obtains for the energy-biased algorithm
$$\mathrm{๐
๐๐}_{EB}[\beta \mu ]=\frac{1}{n_0}\left(\frac{1}{F_w}+\frac{1}{s_c_f}\right).$$
(15)
In deriving Eq.(15) we used that $`_c_f=F_w_c_h`$ up to a negligible amount. This can be seen by noticing that the function $`[\beta (uc)]`$ in the integrand of $`_c_f=_{\mathrm{}}^{\mathrm{}}f(u)[\beta (uc)]๐u`$ decays exponentially for $`u>c`$. Hence, in any practical case ($`u_w>c`$) most of the integral weight comes from $`u<u_w`$, for which the energy-biased reconstruction of the energy profile $`f(u)`$ is exact (see Fig. 3).
We now evaluate the cost, which is given by the total number of energy evaluations of the test molecule needed to obtain $`n_w`$ samples:
$$n_{cost}=n_0+n_w/๐,$$
(16)
where $`๐<1`$ is the acceptance ratio of the Hit&Run sampling algorithm, defined in Appendix A. Introducing Eq.(14) into Eq. (16) we obtain
$$n_{cost}=n_0\left(1+\frac{sF_w}{๐}\right).$$
(17)
For the energy-biased algorithm the efficiency is $`\epsilon =(n_{cost}\mathrm{๐
๐๐}_{EB}[\beta \mu ])^1`$. Using Eq.(15) and Eq.(17) one obtains
$$\epsilon _{EB}^1=\frac{1}{F_w}+\frac{1}{s_c_f}+\frac{s}{๐}+\frac{F_w}{๐_c_f}.$$
(18)
By maximising the efficiency $`\epsilon =\epsilon (F_w)`$ in Eq. (18) with respect to $`F_w`$, one obtains the optimal value $`F_w^{opt}`$ and the maximum efficiency $`\epsilon _{EB_{\mathrm{max}}}=\epsilon _{EB}(F_w^{opt})`$:
$`F_w^{opt}`$ $`=`$ $`\sqrt{๐_c_f}`$ (19)
$`\epsilon _{EB_{\mathrm{max}}}^1`$ $`=`$ $`2{\displaystyle \frac{1}{\sqrt{๐_c_f}}}+{\displaystyle \frac{s}{a}}+{\displaystyle \frac{1}{s_c_f}}.`$ (20)
Finally, we compare the efficiency of the energy-biased algorithm with that provided by the Bennett algorithm, given by $`\epsilon _B=_c_f`$. According to Eq. (20) the ratio of efficiencies is given by
$$\frac{\epsilon _B}{\epsilon _{EB_{\mathrm{max}}}}=2\sqrt{\frac{_c_f}{๐}}+\frac{s_c_f}{๐}+\frac{1}{s}.$$
(21)
Equation (21) yields the range of values of $`_c_f`$ for which the energy-biased Bennett estimator for $`\beta \mu `$ method is more efficient than the standard (unbiased) Bennett algorithm. Note that for $`s=\sqrt{๐/_c_f}`$ the efficiency ratio given by Eq. (21) reaches its minimum value, $`\epsilon _B/\epsilon _{EB_{\mathrm{max}}}=4\sqrt{_c_f/๐}`$, and therefore $`\epsilon _B<\epsilon _{EB}`$ if $`_c_f>๐/16`$. Hence the energy-biased method is suited for fluids at high densities or low temperatures or for molecular fluids with low insertion probability. In this regime $`_c_f<<๐/16`$ and the dominant term in Eq. (21) is $`1/s`$, hence $`\epsilon _{EB_{\mathrm{max}}}s\epsilon _B`$. In other words, the maximal efficiency of the present energy-biased method is limited by the average number $`s`$ of independent samples that can be obtained within one energy-well. As shown in Appendix B, for the Lennard-Jones fluid we have observed that in the most unfavourable case (high density and low temperature) $`s[510]`$.
### III.2 Reconstruction of the energy distribution
We now show that the reconstruction of $`f(u)`$ using the energy-biased procedure (EB) is faster and more efficient than that obtained using any unbiased sampler which uniformly explores the whole configurational space. To that end we consider the evaluation of the cumulative probability $`F(u)=_{\mathrm{}}^uf(u^{})๐u`$ for $`u<u_w`$ (i.e. for $`F(u)<F_w`$). We shall compare the variance of two estimators for $`F`$: one based on uniform insertion over the whole domain and the other based on the energy-biased procedure. The variance of the unbiased estimator is simply $`\mathrm{๐
๐๐}(F)=F(1F)/n_0`$ and for low energies ($`F<<1`$) its efficiency is $`1/F`$. The expected value of the energy-biased estimator is $`HF_w`$, where $`H(u)=_{\mathrm{}}^uh(u^{})๐u^{}`$ is the cumulative probability of the biased distribution in Eq. (7). This estimator is constructed as a product of two statistically independent fluctuating variables and its variance is Goodman (1960)
$`\mathrm{๐
๐๐}_{EB}(F)`$ $`=`$ $`\mathrm{๐
๐๐}(HF_w)=F_w^2\mathrm{๐
๐๐}(H)`$
$`+`$ $`H^2\mathrm{๐
๐๐}(F_w)+\mathrm{๐
๐๐}(F_w)\mathrm{๐
๐๐}(H).`$
Using $`\mathrm{๐
๐๐}(H)=H(1H)/n_w`$ and $`\mathrm{๐
๐๐}(F_w)=F_w(F_w1)/n_0`$ one obtains
$$\mathrm{๐
๐๐}_{EB}=\frac{F_w(F_w1)H^2}{n_0}+\frac{H(1H)F_w^2}{n_w}+\frac{F_wH(1H)}{n_0n_w}.$$
(23)
Note that, as expected, for $`H1`$ one recovers the variance of the unbiased insertion method. The interesting part of the energy distribution is the importance region, located in the low energy range, where $`H<<1`$. In this regime one can make the approximation $`1H1`$. Using $`F=HF_w`$ and $`n_w=n_0sF_w`$, one gets
$$\mathrm{๐
๐๐}_{EB}=\frac{F}{n_0}\left(\frac{F}{F_w}+\frac{1}{s}+\frac{1}{n_0sF_w}\right).$$
(24)
Note that the term in brackets is the reduction in variance with respect to uniform unbiased sampling. Because $`F_w`$ is evaluated from $`n_0`$ probes, this means that necessarily $`n_0>>1/F_w`$ so the third term inside the brackets is much smaller than unity. On the other hand, for the low energy range considered $`F<<F_w`$ and one finally concludes that $`\mathrm{๐
๐๐}_{EB}\mathrm{๐
๐๐}(F)/s`$, where $`\mathrm{๐
๐๐}(F)F/n_0`$ is the variance obtained in the unbiased uniform sampling of the whole domain.
The cost associated with the energy-biased procedure is $`n_{cost}=n_0(1+sF_w/๐)`$. In the case of a Lennard-Jones liquid we have found that $`๐0.17`$ and $`sO(10)`$, while the optimal cumulative probability is $`F_w10^3`$. This means that, in practical situations, $`sF_w/๐1`$ and $`n_{cost}n_0`$. Thus, according to Eq. (24) the energy-biased sampling procedure is around $`s`$ times faster than a uniform unbiased (grid or random) sampler in reconstructing the low energy range of $`f(u)`$. As before, $`s`$ is the average number of independent samples taken per well.
## IV Results
In order to confirm the foregoing theoretical relations about efficiency and variance reduction, we performed molecular dynamics simulations of a Lennard-Jones liquid at high density and low temperature ($`\rho =0.0236\AA ^3`$ and $`T=84`$K). These simulations were performed in a cubic periodic box of side $`L=10\sigma `$. We used the standard Verlet method Allen and Tildesley (1987) to integrate Newtonโs equations of motion, incorporating a Langevin thermostat Kremer and Grest (1990) to keep the system in the NVT ensemble.
During the simulation, the iterative loop (1)+(2) explained in Sec. II was performed $`m`$ times per time interval $`\delta t_{samp}=0.5\tau `$, which corresponds to about three times the collision time. The search for wells performed in step (1) was done by probing at the nodes of a Cartesian grid comprising $`15^3`$ nodes. This ensured that the explored cavities are independent. All the cavities found in step (1) were sampled using the Hit&Run algorithm (see Appendix A).
### IV.1 Estimation of the chemical potential
One way to measure the efficiency of the method is to evaluate the convergence of the estimated value of the chemical potential for an increasing number of test-particle probes $`n_{cost}`$. Convergence can be calculated from the difference between successive values of $`\mu _n`$, where $`n(=n_{cost})`$ indicates the total number of evaluations of the test-particle energy. Figure 2 shows how this difference decreases in calculations based on both the energy-biased and the unbiased samples. These calculations correspond to liquid argon with number density $`\rho =0.0236\AA ^3`$ and temperature $`T=84`$K (these values correspond to $`\rho =0.92\sigma ^3`$ and $`T=0.7`$ in Lennard-Jones units), for which the average of the Fermi function is $`_c_f=8.9\times 10^6`$. According to Eq. (19) the optimum value of $`F_w`$ is $`0.0012`$, which corresponds to $`u_w14.19`$ Kcal/mol. We selected the predicted optimum parameter ($`u_w=14.19`$ Kcal/mol) and performed $`d=15`$ samples per well. As can be seen in Fig. 2, for equal numbers of energy probes ($`n=n_{cost}`$), the average difference between successive estimates of the chemical potential via the energy-biased method is about five times smaller than that obtained with the unbiased sampler. As predicted by Eq. (21), such a gain in efficiency is consistent with the average number $`s`$ of independent samples per well (see table 2), which for this simulation was $`s5`$.
Evaluations of the chemical potential for Lennard-Jones (LJ) fluids are shown in Table 1 together with the estimated efficiency of each calculation. For a LJ fluid with $`\rho =0.02360\AA ^3`$ and $`T=84`$K the numerically obtained net gain is around 7, which coincides with the prediction in Eq. (21) using $`s=7`$. For illustrative purposes we also analysed a case for which the efficiency of our implementation of the energy-biased sampling is similar to the uniform-unbiased Bennett method. For instance, $`_c_f=0.0102`$ for $`\rho =0.01755\AA ^3`$ and $`T=178.5`$K. Using $`๐=0.165`$ and the (optimum) number of samples $`s=\sqrt{๐/_c_f}4`$ in Eq. (21) one obtains $`\epsilon _B/\epsilon _{EB_{\mathrm{max}}}1`$; our numerical calculations, with $`u_w=7.33`$ and $`d=8`$, confirmed this conclusion. We note that for any value of $`u_w`$ considered the energy-biased estimation of the chemical potential $`\mu `$ agrees within about $`0.01`$ Kcal/mol with the unbiased Bennett result. This is illustrated in Table 2 where we show the estimated $`\mu `$ for the higher density liquid, using several values of $`u_w`$.
### IV.2 Reconstruction of the energy distribution $`f(u)`$
In Fig. 3 we compare the reconstructed energy distribution $`f(u)`$ at energies $`u<u_w`$ with that computed from an unbiased method, which consists of a large number of random insertions within the entire configurational space. Figure 3 clearly illustrates that the energy-biased method exactly reproduces the unbiased distribution $`f(u)`$ for energies smaller that $`u_w`$. This attractive feature is a consequence of the fact that it is easy to exactly correct for the bias in terms of the cavity energies. This is not true for the accessible volume of the molecule, as in cavity-biased procedures Jedlovszky and Mezei (2000); Deitrick et al. (1989).
In order to illustrate the above conclusion we show in Fig. 4 the estimation of the cumulative probability $`F(u)`$ versus the total number of test-particle energy probes used for the evaluation. The particular case shown corresponds to $`u=5`$ Kcal/mol, for a LJ liquid at $`\rho =0.0236\AA ^3`$ and $`T=84`$K. The energy-biased sampling was done using $`u_w=14.19`$ Kcal/mol and $`d=15`$ samples per well, and for this calculation we obtained $`s5`$ (see Appendix and Table 2). Compared with the unbiased procedure, the reduction of variance provided by the energy-biased sampler is immediately apparent on inspection of Fig. 4. A numerical evaluation of the variance of each data set in Fig. 4 provides: $`\mathrm{๐
๐๐}_{EB}=4.14\times 10^5/n_0`$, while the (best) result for the algorithm based on uniform unbiased sampling is $`\mathrm{๐
๐๐}(F)=F/n_0=1.9\times 10^4/n_0`$. Hence the net gain in efficiency is about 4.6, in agreement with the value of $`s=5`$ obtained from the independent correlation analysis explained in Appendix B. As shown in Table 1, the estimated net gain in the evaluation of the chemical potential compared with the unbiased Bennett method is $`7\pm 1`$, which is close to the estimate $`s5`$ obtained from the analysis of the cumulative probability.
## V Conclusion
We have presented a new method for sampling the energy of a test-molecule in order to calculate single-particle ensemble averages and, in particular, the chemical potential. The method, called energy-biased sampling, restricts the important region to the bounded domains in the test-molecule energy-landscape where the test-molecule energy $`u`$ is smaller than a given free parameter $`u_w`$. This energy-biased sampling retains the principal benefit of cavity-biased methods Jedlovszky and Mezei (2000); Deitrick et al. (1989) in the sense that, by sampling only within regions with a significant Boltzmann factor, convergence is greatly accelerated with respect to uniform sampling. Furthermore, because the energy-biased sampling is accurately defined in terms of the test-particle energy it has some important benefits: first, it allows accurate reproduction of the test-particle energy distribution $`f(u)`$ and the chemical potential; second, it is possible to sample cavities of arbitrary shape (not only spherical ones) and to generalise the cavity dimensionality to include the rotational degrees of freedom in the energy-well reconstruction; finally, and rather importantly, it enables one to combine the sampling results with standard free energy perturbation (FEP) formulae. In particular, we combined it with the Bennett method Bennett (1976) which minimises the variance of the estimator and has proved to be the best method in the literatureLu et al. (2003); Shirts and Pande (2005). Energy-biased sampling is a general protocol to bias the sampling and consists of two sequential steps: (1) searching and (2) sampling the interior of energy-wells. In this work we have implemented these two steps using relatively simple algorithms: uniform unbiased search and Hit&Run sampling. However we note that other solutions are also possible. For instance, non-uniform sampling of the importance region may surely increase the efficiency of the present method. In dense systems, the searching step becomes the most difficult one and a more effective extension of this method could be to perform a biased search (using, for instance, some variation of the usher algorithm Delgado-Buscalioni and Coveney (2003); De Fabritiis et al. (2004)) so as to significantly increase the probability of finding favourable cavities for insertion of the test particle. These extensions are left for future studies.
###### Acknowledgements.
This research was supported by the EPSRC Integrative Biology project GR/S72023 and by the EPSRC RealityGrid project GR/67699. R.D-B acknowledges support from the European Commission via the MERG-CT-2004-006316 grant and from the Spanish research grants FIS2004-01934 and CTQ2004-05706/BQU.
## Appendix A Sampling bounded regions with the Hit&Run algorithm
There exists a relatively large literature on sampling a bounded connected region (see for instance Ref. Liu (2001) and references therein). In this work we have used the so-called Hit&Run algorithm for its simplicity and good performance Liu (2001). The Hit&Run sampler is a special Monte Carlo Markov chain which draws numbers from an assigned distribution Smith (1984); Liu (2001) $`p(๐ซ)`$, where $`๐ซA`$ lies within a bounded connected region of an n-dimensional space $`AR^n`$. In our case, $`p(๐ซ)`$ is a uniform probability density over the region $`A_{u_w}^\alpha `$ such that
$$p(๐ซ)=\frac{1}{\mathrm{\Omega }(A_{u_w}^\alpha )}.$$
(25)
The Hit&Run algorithm starts from a point $`๐ซ_\mathrm{๐}`$ within the bounded region $`A`$ and performs the following steps:
1. Choose a random direction $`๐`$ and find the intersections of the cavity border with the line $`๐ซ(\lambda )=๐ซ_0+\lambda ๐`$, where $`\lambda `$ is a real number. As the cavity $`A`$ is bounded the intersection is composed by two points $`๐ซ(\lambda ^+)`$ and $`๐ซ(\lambda ^{})`$ (here $`\lambda ^+>0`$ and $`\lambda ^{}<0`$).
2. Select a point $`๐ซ_1`$ within the segment ($`๐ซ(\lambda ^+)`$, $`๐ซ(\lambda ^{})`$), i.e.,
$$๐ซ_1=๐ซ(\lambda ^{})+\xi (๐ซ(\lambda ^+)๐ซ(\lambda ^{}))$$
(26)
where $`\xi (0,1)`$ is a uniformly distributed random number.
3. Sample at $`๐ซ_1`$, set $`๐ซ_1๐ซ_0`$ as the new starting point and go to (i).
The above procedure is repeated to obtain the desired number of samples $`d`$. In our case the starting point for the sample chain, $`๐ซ_0`$, is the test-particle configuration returned by the algorithm for energy-well searching ($`U(๐ซ_0,๐)<u_w`$). In order to locate the borders of the energy well $`๐ซ(\lambda ^+)`$ and $`๐ซ(\lambda ^{})`$ we use the following procedure. Starting from $`๐ซ_0`$ we cross the well along the line defined by the random unit vector $`๐`$ moving in steps of size $`\delta s`$, i.e., according to
$$๐ซ(k)=๐ซ_0+k\delta s๐,$$
(27)
with $`k`$ being an integer starting from $`k=\pm 1`$. The energy is computed at each point $`๐ซ(k)`$ until one crosses the edges of the well at $`k=k^+`$ and $`k=k^{}`$ (for which $`u(๐ซ(๐ค^\pm ),๐)>u_w`$). An approximate location of the cavity borders is provided by setting $`\lambda ^\pm =k^\pm `$. We used typically $`\delta s0.3\AA `$ and required, on average, about five iterations to cross the well in one random direction (this value depends on the density and $`u_w`$). Note that the acceptance ratio is $`๐=k^+k^{}^1`$ and for the high density cases considered here $`๐0.17`$.
## Appendix B Optimal number of sampling directions
It is possible to reduce the cost without increasing the variance by setting the number of samples per cavity $`d`$ equal to or somewhat larger than $`s`$, the average number of independent samples per cavity. Note that the number of statistically independent samples within one cavity is $`s=d/\tau _c`$, where $`\tau _c`$ is an empirically estimated autocorrelation length of the whole chain of data. This number $`\tau _c`$ can be estimated from the large $`m`$ limit of the quantity $`m\mathrm{๐
๐๐}[^{(m)}]/\mathrm{๐
๐๐}[]`$, where $`_c=[\beta (uc)]`$ is the Fermi function evaluated at a single energy $`u`$ and $`^{(m)}`$ denotes the mean of $`m`$ consecutive $``$ values.
The value of $`s`$ can be estimated by performing several Hit&Run samplings with an increasing number of directions per cavity $`d>s`$, then computing $`\tau _c`$ for the chain of samples and evaluating $`d/\tau _c`$, which should be nearly independent of $`d`$. We carried out this evaluation of $`s`$ for varying values of $`u_w`$ within the same system and for fixed $`u_w`$ and varying density. The results of this study, reported in Table 2, clearly indicate that $`s`$ does not greatly vary for a broad range of values of the cavity-border energy $`u_w`$. In fact, at low and moderate values of $`u_w`$ the energy-cavities are isolated and their average size (in $`\AA `$) grows quite slowly with $`u_w`$. This is due to the steepness of the hard-core part of the Lennard-Jones potential. Above a certain energy $`u_w`$ the cavities become connected and a steep rise in the average size of the energy-cavities is observed. This is reflected in the value of $`s`$. As shown in Table 2 for $`u_w=14.19`$ Kcal/mol we obtained $`s4.5`$ and $`s11`$ for two calculations using $`d=15`$ and $`d=100`$ respectively. We obtained a relatively close value $`s7`$ for twice as large an energy limit $`u_w=28.38`$ Kcal/mol. However using $`u_w=165.53`$ Kcal/mol the average number of independent samples increased up to $`25`$, reflecting the more complex shape and larger volume of these energy cavities. In summary, for the optimum range of values of $`u_w[1030]`$ Kcal/mol we find $`s[510]`$ in the case of the Lennard-Jones liquid. |
warning/0506/astro-ph0506737.html | ar5iv | text | # The NGC 5846 Group: Dynamics and the Luminosity Function to ๐_๐
=-12
## 1. Introduction
The mass function of dark matter halos (i.e., their abundance as a function of mass) is an important ingredient in constraining cosmological parameters using galaxies and clusters of galaxies (Frenk et al. 1990; Kauffmann & White 1993; Haiman et al. 2001; Reiprich & Bรถhringer 2002; Hoekstra et al. 2002). Very often, an assumed theoretical formula for the mass function is used to fit observations so that cosmological quantities such as the mean matter density of the universe $`\mathrm{\Omega }_m`$ may be constrained (Kochanek et al. 2001; Haiman et al. 2001; Henry 2004). However, almost all these methods neglect the fact that the abundance of small, $`10^910^{10}M_{}`$ halos is poorly understood. In fact, the very N-body simulations used to derive forms for this fitting function, e.g. Sheth & Tormen (1999), often produce an abundance of such halos (sometimes referred to โsatellite galaxiesโ) that is far in excess of the observations (Moore et al. 1999; Kazantzidis et al. 2004; van den Bosch et al. 2005). One possibility is that the relationship between the light and mass distributions is poorly understood at small masses (Gao et al. 2004)โat the scale of dwarf galaxies. It is therefore critically important to constrain the *luminosity* function of these galaxies (Trentham & Tully 2002; Tully et al. 2002; Trentham et al. 2005), thus providing a crucial check on the simulations that attempt to model galaxy formation.
It is becoming clear that the relationship between the amount of starlight (or gas) that we see and the mass of dwarf halos is not simple. There is evidence that the amount of light associated with halo mass varies strongly with environment (Tully 2005). It may be that there are dark matter halos that have retained baryons in only undetectably small amounts. Various astrophysical processes can lead to a separation of dark matter and baryons, e.g., the ram pressure stripping processes such those observed in the famous โbullet clusterโ 1E 0657-56 (Markevitch et al. 2002). The relationship between the luminosity function and the mass function might be complex, but the luminosity function may nevertheless retain signatures of specific astrophysical processes or of the underlying dark matter spectrum itself.
The present article is a contribution within a long term program to provide a better definition of the faint end of the luminosity function of galaxies. The general properties of the program are as follows. First, the observations should reach very faint absolute magnitudes. This goal can only be achieved if the targets are nearby. Second, good statistics are needed in conditions that provide control of volume completion. The issue of volume completion is addressed by obtaining complete samples to an apparent magnitude limit in groups selected to have minimal contamination problems. A single group may or may not provide adequate statistics by itself. Observations of numerous groups may be required. Third, a wide variety of environments should be sampled in order to constrain the possibility of environmental dependencies. This requirement imposes a need for many nights of observations.
The program depends critically on the recent availability of panoramic digital cameras. The imaging material discussed in this paper was acquired with the Canada-France-Hawaii Telescope (CFHT) 12K detector (the predecessor of the current MegaCam). The program also depends on some manner of confirmation that galaxy candidates are group members, hence relevant to the construction of the luminosity function. Confirmation is most reliably provided by redshifts, and this paper includes the results of spectroscopic observations of relatively faint, low surface brightness dwarfs with the LRIS instrument on the Keck I Telescope..
Initial attention was given to well populated clusters in both the high density (Trentham 1998b, 1998a) and low density (Trentham et al. 2001, hereafter TTV01) regimes. These observations suggested that there are significant variations with environment (Tully et al. 2002)โthat the luminosity function of dwarf galaxies rises steeply in dense environments, but remains flat in the field. These results prompted an exploratory sampling of a wide variety of locations within the Local Supercluster with the Subaru Telescope SuprimeCam imager (Trentham & Tully 2002, hereafter TT02). It became clear that much more sky needed to be observed to build up meaningful statistics. This paper presents results from the first group in the program to receive full coverage: a tight, well defined knot of early type galaxies surrounding the elliptical galaxies NGC 5846 and NGC 5813.
## 2. Observations
### 2.1. NGC 5846: A Well-Isolated Group of Galaxies
This group is readily apparent because of its high density contrast and relatively isolated location. It has been studied as a group of galaxies in both the X-ray and the optical (Tully 1987; Haynes & Giovanelli 1991; Nolthenius 1993; Giuricin et al. 2000; Trinchieri & Goudfrooij 2002; Mulchaey et al. 2003). The distance to the system is taken to be 26.1 Mpc from an average over several sources, but most heavily reflective of the Surface Brightness Fluctuation measurement of Tonry et al. (2001). These authors give a distance to NGC 5846 itself of $`25\pm 4`$ Mpc and a distance to NGC 5813 of $`32\pm 3`$ Mpc, estimates we consider compatible with the two being at a common distance. The group members are overwhelmingly of early type, dominated by the ellipticals NGC 5846 ($`M_R=22.5`$) and NGC 5813 ($`M_R=22.2`$). Within this distance, the group is the third most massive knot of early type galaxies (after the Virgo and Fornax clusters). The virial mass based on velocity information to be discussed in a later section, is $`8\times 10^{13}M_{}`$, and the ratio of mass to light at R band is $`320M_{}/L_{}`$.
The isolation of the NGC 5846 Group is extremely favorable. It lies well off the main plane of the Local Supercluster, with no known structure to the foreground and very little in the background until the Hercules Supercluster at 10,000 km s<sup>-1</sup>. Figure 1a is a histogram of all Sloan Digital Sky Survey (Abazajian et al. 2004, hereafter SDSS) and NASA Extragalactic Database (NED hereafter) velocities of galaxies within $`3^{}`$ (1.4 Mpc) of the center of NGC 5846 (not including our survey below). Remarkably, there are no galaxies at all within $`3000<cz<6,000`$ km s<sup>-1</sup>.
Figure 2 shows the distribution of all but two of the galaxies that are established to be within 3ยฐ of the group on the basis of a measured velocity; the two remaining objects are described only as H I detections in the literature, without optical counterparts. The shown galaxies have velocities in the range $`900<cz<2700`$ km s<sup>-1</sup> .
Over the range of environments that will be explored in this program, the NGC 5846 Group lies in the regime of high density and intermediate mass. It possesses 4 galaxies brighter than $`L^{}`$ (3 E/S0 and an Sb). In this paper, 324 galaxies are identified as probable or possible members, extending in faintness down to $`M_R10`$.
### 2.2. Wide Field Imaging
Observations of the NGC 5846 Group were made with the CFH12K CCD camera in queue mode during 11 nights between 16 March 2002 and 10 June 2002. An overall rectangular area of $`220^{}\times 180^{}`$ was surveyed with a small hole to avoid the glare of a 4th magnitude star. Two control fields were observed $`9^{}`$ north, in a region with no known objects foreground of $`10,000`$ km s<sup>-1</sup>. The CFHT12K detector is a mosaic of 12 CCD detectors providing a field of $`42^{}\times 28^{}`$, oriented in this experiment with the long axis E-W. The observations were tiled with half-field overlaps and dithers so that gaps between CCD chips were almost entirely covered in subsequent exposures and most of the area was observed twice. In total, 67 x 9 minute exposures were taken, all in the Cousins R band, covering 10.05 square degrees. Two of the 11 nights were non-photometric but photometry could be propagated across the entire survey region through the half-field overlaps. Seeing was 0.7โ1.0 arcsec as mandated by the queue request. Images of the individual dwarfs will be available online via the CFHT image cutout service (planned for the future).
Members of the NGC 5846 Group range from spectroscopically accessible high surface brightness objects ($`\mu _R<20`$ mag arcsec<sup>-2</sup>) to spectroscopically challenging very low surface brightness objects ($`\mu _R>20`$ mag arcsec<sup>-2</sup>). Experience has shown that most low luminosity galaxies are low surface brightness, although there are exceptions (Drinkwater et al. 2003). Candidates can be isolated on the basis of this property. Morphological criteria can then be applied to further evaluate the probability of group membership. The details of the procedures to chose candidates have been described in TTV01 and TT02. The essence of the process is a culling of the very large background population with a concentration index threshold. Known dwarfs lie in a distinctly lower concentration index regime that giant galaxies. Nonetheless, substantial numbers of background contaminants manage to pass the concentration index screen and must be culled based on morphological criteria. This latter step requires inspection of images with the following considerations in mind. If there is evidence for a bulge or major bar, or spiral structure, or tidal disruption then such objects that pass the concentration filter are probably background. If instead an object is diffuse except possibly for a semi-stellar nucleation or patchy structure then it is probably nearby, hence a group member.
In TTV01 and TT02 there is a description of a rating scheme developed to characterize the probability of group membership for galaxies which have passed the surface brightness criterion (1: probable member; 2: possible member; 3: conceivable member; 4: almost certainly background). Additionally, a rating 0 is given to galaxies that were identified as group members on the basis of redshifts before the survey began, whatever their surface brightness, and a rating 5 is given to galaxies initially rejected as candidates because they lie above the concentration index cut but subsequently associated with the group by a redshift. The availability of velocity information for a significant fraction of the sample now permits a reevaluation of the validity of the morphology-based rating scheme. In ยง2.4 below we show that galaxies rated 1 and 2 are essentially always found to be group members, of order half the galaxies rated 3 are found to be group members, but almost none of the galaxies rated 4 are group members. Five high surface brightness galaxies (identified with a rating 5) are found to be members.
Before considering the SDSS and out Keck data, 318 galaxies are identified in the survey region with ratings 0-3. After including the SDSS, the number increases to 324, because 6 priority 4-5 galaxies not originally included in our sample are confirmed by the SDSS as members. These objects are identified in Table 1. The magnitudes presented in this table are isophotal $`R`$-band magnitudes extracted to an isophote of 25.2 mag arcsec<sup>-1</sup>. Their distribution on the sky is shown in Figure 2.1. There is a strong enhancement in the surface number density of both confirmed members and candidates surrounding the elliptical NGC 5846 and a secondary enhancement surrounding the elliptical NGC 5813.
The two control fields $`9^{}`$ N were chosen to lie off the filament containing the NGC 5846 Group and in the direction of the Local Void (Tully & Fisher 1987). No galaxies were found in these fields that could be rated 1โ3. On the basis of the detection rate in the NGC 5846 area, one would anticipate 9 such candidates in these fields.
### 2.3. Spectrocopic Observations
The NGC 5846 Group lies within the area with published spectroscopic information from Data Release 3 of the SDSS. The NASA/IPAC Extragalactic Database (NED) and the SDSS provide velocities for 64 of our targets, all of them brighter than $`R17`$ or $`M_R15`$. There is essentially no prior literature information for candidates at fainter magnitudes or for those with very low surface brightnesses. In order to probe these regimes, we undertook observations using the blue side of the Low Resolution Imaging Spectrograph (Oke et al. 1995, LRIS hereafter), on Keck I Telescope. The spectrograph, equipped with a 1โณ slit and a 600 lines in<sup>-1</sup> grating, has high blue quantum efficiency, with an overall system throughput of 56% at 5000 ร
. Data were acquired on May 4, 2003 and June 13, 2004. Both nights were hampered by cirrus with extinction of 0.5-2 mag; nevertheless we obtained suitable spectra for 17 and 13 galaxies on the two nights, respectively, of which a total of 19 are newly reported members. The poor conditions limited attempts to observe the faintest dwarfs; however, good spectra were acquired for objects as faint as $`R=18.8`$ or $`M_R=13.4`$. An LRIS spectrum and CFH 12K image are shown in Figure 3 of the faintest galaxy with a redshift.
### 2.4. Spectroscopic Membership Confirmation
The availability of a large number of new velocities provides a way of evaluating the membership rating scheme based on a concentration parameter and a qualitative judgment based on morphology. The NGC 5846 Group provides an environment that is particularly well-suited to this evaluation, because there is negligible confusion from the foreground or near background. Table 2 provides a summary of how things have turned out. We found all 26 surveyed galaxies with new velocities rated 1 (probable) and 2 (possible) to be group members. Of the galaxies rated 4 (likely background) only one relatively large galaxy has been revealed by spectroscopy to be a member. Among the thousands of galaxies in the survey region excluded by the concentration criteria, 304 have measured redshifts; of these, only 5 galaxies have been demonstrated to be group members.
Thus candidates rated 1 and 2 ought to be group members and, by contrast, very few group members emerge among galaxies rated 4 or those excluded because of high concentration. That leaves the galaxies rated 3 (conceivable member) to be considered. We find that velocity measurements confirmed 16 galaxies with rating 3 to be members and 7 galaxies to be background.
In summary, the combination of the quantitative concentration parameter and the qualitative morphological evaluation leads to good membership discrimination. There remains a grey area with the candidates rated 3. High surface brightness objects elude discovery in the imaging survey but can be found with a spectroscopic survey. Having said all this, the spectroscopic confirmation is complete only to $`M_R15`$ and sampled only to $`M_R13.3`$. There cannot be complete confidence that the rating scheme works among the fainter galaxies that extend down to $`M_R10.5`$. It would be surprising, though, if there is a population of high surface brightness galaxies in the group at these faint magnitudes that makes a significant numeric contribution to the luminosity function.
Henceforth we refer to all member galaxies with redshifts as โspectroscopically confirmed membersโ; when we refer to priority 0-3 galaxies, we mean only those without a redshift.
### 2.5. Indicative Membership: Spatial Correlation
Concentrations of galaxies toward NGC 5846 and NGC 5813 are seen in the galaxy projections of Fig. 2 and in Fig. 2.1 introduced in the next section. A comparison of angular 2-point correlations among various sub-samples gives hints of different degrees of clustering. Sub-samples with relatively weak correlations are probably contaminated by non-members. The correlation that exists among all 324 galaxies rated 0-3 is shown in the upper left panel of Figure 4. The normalization in the correlation is achieved by comparison with 1000 Monte Carlo random populations of the area of the photometric survey. The entire area is overdense with respect to a fair sample of the universe so the amplitude of the normalization is given no meaning.
Next consider the correlation shown in the same panel for 199 galaxies either with membership confirmed by redshift or rated 1โ2. The spectroscopic evidence suggests that most of the sample receiving these rating are members. Hence the correlation function should be fairly representative of the true global function for the group. We refer to this function as the โreference correlationโ hereafter. The reference correlation is more peaked than the function describing the ensemble of the candidates. The ensemble sample must be contaminated by objects drawn from the background.
Figure 4b shows the correlation for the 83 galaxies established to be group members on the basis of velocities. This correlation function is noisier because of smaller numbers but is comparable to the function for the 199 member priority 0-2 sample. The reference correlation is actually slightly steeper which tends to confirm that the galaxies rated 1 and 2 are overwhelmingly members. A different situation is found in the correlation shown for the 125 galaxies rated 3. The much flatter distribution is evidence that quite a few galaxies rated 3 are non-members. The spectroscopic information already discussed revealed that only 16 of 23 candidates rated 3 (70%) are members.
In the Figure 4c, the correlation function is shown for the 61 galaxies with the morphological designation dE,N, nucleated dwarf ellipticals. Only 12 of these have membership confirmed by velocities. Most lie at faint magnitudes. The correlation analysis suggests that, overwhelmingly, galaxies of this morphology are members.
The spectroscopic information extends to galaxies as faint as $`M_R=13.4`$. In total there are 38 priority 3 candidates brighter than this limit. The correlation analysis for these galaxies appears in Figure 4d. The correlation is intermediate between the reference function and that shown by the entire rating 3 sample in panel (b). We create a simple model in which 30% of the galaxies in the reference correlation are replaced by randomly distributed objects. This concoction of 70% correlated and 30% uncorrelated components provides a good description of the distribution seen in the bright priority 3 galaxies. This result agrees with the spectroscopic information available for half the objects in question.
There is no velocity information concerning the fainter candidates, and from the poorer correlation seen in Figures 4e-4f it can be inferred that many of the fainter rating 3 targets are drawn from the background. In fact, the early and late morphological types rated 3 have significantly different correlation characteristics. The distinct distributions are revealed in Figure 4e (early types) and 4f (late types). The 77 galaxies typed dE, dE,N, and dE/I show the correlation seen in panel (e) that is well described by a mix of objects that are 50% correlated and 50% uncorrelated (uncertainty $`10\%`$). By contrast, the late types (mostly dI, a few VLSB) show no correlation to the group. To summarize, the correlation distribution of the rating 3 candidates seen in the filled symbols in Fig. 4b can be decomposed into 70% group members among those brighter than $`M_R=13.4`$, 50% group members among fainter early types, and essentially no group members among fainter late types.
One final consideration is morphological segregation of the galaxies. It is well known that in groups and clusters, early-type galaxies are more densely clustered than late-type galaxies (Dressler 1980; Postman & Geller 1984; Helsdon & Ponman 2003). If the fraction of early-type galaxies in the three priority groups differs greatly, the morphology-density effect would significantly influence the correlation functions. This would lead us to misinterpret differences among the correlation functions as differences in the membership probabilities of the galaxies. Fortunately, Table 2 shows that the morphological content of the three priority classes is quite similar. The weighted average early type fraction in the priority 0-2 galaxies is 0.82, and for the priority 3 galaxies it is 0.79. The difference in the correlation functions is too great to be explained just by this small difference in the early-type fraction.
Based on this simple model, we can calculate an estimate of the total group membership within our survey limits: 83 spectroscopically confirmed members; 32 priority 1 members; 84 priority 2 members; 60-80% of 14 bright priority 3 members; and 40-60% of the 84 faint early-type priority 3 members. The total estimate is $`251\pm 10`$ group members.
## 3. Structure of the NGC 5846 System
### 3.1. Broad X-ray and Optical Picture
Combining the SDSS data with our observations, we are able to provide a broad picture of the environment of the NGC 5846 system. As evident in Figures 1 and 2.1, the system possesses complex substructure, but remains remarkably well isolated from other groups or clusters.
Another clue to the physical state of the system is the distribution of X-ray emission across the group. In elliptical-dominated groups, a hot, $`1`$ keV optically thin plasma often constitutes the largest baryonic component of the groupโs mass. Archival ROSAT, ASCA, Chandra, and XMM-Newton observations of the NGC 5846 region exist. The ROSAT All-Sky Survey (RASS) covers the entire 10 sq. deg. field of our observations. However, with an average exposure time of only 355s, only the X-ray emission closest to the two brightest galaxies, NGC 5846 and NGC 5813, is detected in the survey. We show the 0.5-2.0 keV RASS emission in Figure 2.1. A weak X-ray source $`40\mathrm{}`$ to the north of NGC 5846 is visible, but it is not associated with any group galaxies, and no optical sources in either the POSS or the SDSS appear to coincide with it. Also shown is the adaptively smoothed light distribution of the region. Each member galaxy is represented by a Gaussian on the sky with standard deviation equal to the distance to the third nearest member galaxy. The Gaussians are then added together. Figure 2.1 makes it clear that the two highest surface brightness features in the group, in both the X-ray and the optical, are the regions within $`300`$ kpc of the two brightest galaxies.
The broad picture revealed in both the optical and X-ray surface brightness maps is unusual. The NGC 5846 and 5813 galaxies are two roughly equal peaks of emission at both wavelengths, so the system as a whole appears binary. However, there are also other ellipticals nearly as bright, such as NGC 5838 and NGC 5831, without significant X-ray emission but with significant concentrations of smaller galaxies surrounding them. We will explore the dynamics of the system in greater detail in ยง3.3.
### 3.2. Two Subgroups: NGC 5846 and NGC 5813
A central question in understanding the NGC 5846 system regards the existence of โgroup-scaleโ X-ray emission. Both NGC 5846 and NGC 5813 possess diffuse X-ray emission from a hot, optically thin plasma, suggesting that they mark the center of the two largest dark matter concentrations in the system. But is the X-ray emitting gas confined to the galaxies only, emitting X-rays properly only as the hot interstellar medium of the elliptical galaxies (Eskridge et al. 1995), or is there evidence of a true intracluster medium (ICM) in which the galaxies are embedded? An extended ICM would make it clear that the NGC 5846 and 5813 galaxies are not only bright, X-ray emitting ellipticals; they are also the centers of much more massive dark matter halos. Previous analyses have considered *ROSAT* data (Bรถhringer et al. 2000; Ikebe et al. 2002; Osmond & Ponman 2004), but no comparisons with high-quality surface brightness profiles (such as from the SDSS) or discussions of the *XMM-Newton* observations of NGC 5846 exist.
To evaluate the extent of the X-ray emission, we examine *XMM-Newton* data for NGC 5846 (superior in sensitivity and field-of-view to the *Chandra* data) and the *ASCA* data for NGC 5813 (where neither *XMM-Newton* nor *Chandra* data is available). We extract X-ray surface brightness profiles from the available data. In the case of *ROSAT* and *ASCA*, we extract pre-calibrated photon images and exposure maps from the *ROSAT* All-Sky Survey Data Browser<sup>1</sup><sup>1</sup>1http://wave.xray.mpe.mpg.de/rosat/data-browser and the NASA HEASARC data archive<sup>2</sup><sup>2</sup>2http://heasarc.gsfc.nasa.gov/W3Browse/, respectively. For *XMM-Newton*, only images from the highest-sensitivity camera on board the telescope, the EPIC pn, were available. We use the *XMM-Newton* Software Analysis System (SAS)<sup>3</sup><sup>3</sup>3http://xmm.vilspa.esa.es/sas/ to reduce these data.
In the case of NGC 5846, the RASS contributes a useful signal outside the $`20\mathrm{}`$ diameter field-of-view of *XMM-Newton*. In the case of NGC 5813, *ASCA* provides superior data to the RASS. For the RASS, we use the 0.5-2.0 keV data to minimize the noise from the unrelated X-ray background. For the *XMM-Newton* EPIC pn we also use photons with 0.5-2.0 keV energies. For the NGC 5813 observations by the *ASCA*, we use the full-band images covering photon energies 0.7-10 keV. We also extract optical surface brightness profiles from calibrated $`g`$-band imaging data publicly available from the SDSS.
The profiles appear in Figure 5. In order that our analysis is not affected by the point spread function (PSF) of each instrument, we ignore data within a radius 4 times the full width of the PSF at half-maximum; the PSF deformation of the profiles should be negligible beyond these radii (Mohr et al. 1999). We fit the profiles with simple azimuthally symmetric models to evaluate their extent. While the models may not be correct in detail, they are useful tools in characterizing the distribution of light in the group. To fit the SDSS $`g`$-band optical light distribution, we use a modified Sรฉrsic (Sรฉrsic 1968; Trujillo et al. 2004) profile:
$$\mathrm{\Sigma }_g(r)=\mathrm{\Sigma }_0(r/r_e)^\gamma \mathrm{exp}\left[b\left(r/r_e\right)^\alpha \right],$$
(1)
where $`r`$ is the projected distance to the galaxy center, $`r_e`$ is the half-light radius, b is a constant, and $`\alpha `$ and $`\gamma `$ are the characteristic slopes. The case $`\gamma =0,\alpha =1/4`$ corresponds to the widely known de Vaucouleurs et al. (1976) profile. The Sรฉrsic profile ($`\gamma =0`$, $`\alpha `$ free) is known to fit a subset of the bright elliptical galaxies observed with the Hubble Space Telescope (HST). Others, however, require a nonzero inner slope $`\gamma `$ for an acceptable fit (Trujillo et al. 2004) <sup>4</sup><sup>4</sup>4In Trujillo et al. (2004), these so-called โcore-Sรฉrsicโ galaxies have inner profiles of the form $`[1+(r/r_b)^\delta ]^{\gamma /\delta }`$. However, typical values of $`r_b`$ are $`50100`$ pc, requiring much finer resolution to resolve than the data we discuss here possess. Therefore we approximate the inner region of the Trujillo et al. (2004) profile as a simple power law with slope $`\gamma `$..
The X-ray light distribution is characterized by a double $`\beta `$-model (Mohr et al. 1999):
$`\mathrm{\Sigma }_x(r)`$ $`=`$ $`\mathrm{\Sigma }_1(1+r^2/r_1^2)^{3\beta _1+1/2}`$ (2)
$`+`$ $`\mathrm{\Sigma }_2(1+r^2/r_2^2)^{3\beta _2+1/2}.`$
The single $`\beta `$-model was originally used to describe data observed by the *Einstein* observatory (Jones & Forman 1984). Higher resolution *Chandra* observations have since revealed that emission from relaxed clusters is often more complex and requires either a broken power law or the two-component $`\beta `$-model shown above (Buote & Lewis 2004).
The fit results in Table 3 suggest substantial differences between the two galaxies. The $`\gamma =0`$ optical fit for NGC 5846 is of good quality. The brightest group member is consistent with being a de Vaucouleurs et al. (1976) model galaxy ($`\alpha =1/4`$). For NGC 5813 the pure $`\gamma =0`$ Sรฉrsic profile is not a good fit. With $`\chi ^2/\nu =364/291`$, the fit is rejected at better than 99.7% confidence. However, with $`\gamma `$ free, we obtain a good fit. Thus, the dominant galaxies in the group, though possessing similar Hubble types (E0 and E1 for NGC 5846 and 5813, respectively), exhibit different light profiles in detail. NGC 5813โs inner stellar surface brightness profile is consistent with a power law of index $`1.2`$, but that of NGC 5846 flattens near the center. Such variations in profile shapes are also present in field ellipticals observed by HST (Trujillo et al. 2004), though the reasons for the existence of two populations are not yet clear.
One clue to understanding the differences in the light profiles of the two galaxies may lie in their merging histories. Detailed HST observations of the central few arcsec in NGC 5813 reveal a dusty circumnuclear disk (Tran et al. 2001). The inner region of NGC 5846 contains no disk, but X-ray and optical filamentary structures instead (Goudfrooij & Trinchieri 1998; Tran et al. 2001; Trinchieri & Goudfrooij 2002). One interpretation is that the filaments in NGC 5846 are relics of interaction with smaller galaxies (Goudfrooij & Trinchieri 1998). The fact that NGC 5813 contains an undisturbed circumnuclear disk suggests that it may have had fewer recent mergers than NGC 5846. For this reason NGC 5813 retains a steeper, more undisturbed stellar profile than does its sister galaxy.
The X-ray data reveal that the intracluster medium is more extended than the optical light distribution. We obtain acceptable fits for a double $`\beta `$-model in NGC 5846, while for NGC 5813 a single $`\beta `$-model suffices. The ratio of the X-ray to the optical half-light radii for NGC 5846 and NGC 5813 are $`1.3\pm 0.3`$ and $`3.1\pm 0.8`$, respectively; the inferred 90% light radii have ratios of $`>2.2`$ and $`>8.2`$, respectively. Interestingly, though NGC 5813 is optically less bright NGC 5846, it contains a more luminous and extended X-ray halo.
### 3.3. Dynamics
On dynamical grounds, the NGC 5846 Group can be expected to consist of an evolved core surrounded by an infall region. Given the mass of the group, the infall region is expected to extend to $`10^{}`$ radius, considerably beyond the boundaries of the survey region illustrated in Fig. 2. The interest of the current study is with the dense, dynamically evolved core. The dimensions of that core can be inferred from the radial density and velocity distribution of galaxies in the vicinity of the group which are shown in Figure 6. It is seen in the top panel of that plot that the surface number density of galaxies declines smoothly with radius out to $`1.8^{}`$ then abruptly drops a factor $`4`$ to a roughly constant plateau. In the lower panel it is seen that the velocity dispersion within radial shells drops by roughly a factor 2 inside $`1.8^{}`$. Declining velocity dispersion profiles are typical of dynamically evolved clusters (Mahdavi et al. 1999; Biviano & Katgert 2004; Sand et al. 2004; Mahdavi & Geller 2004). Beyond 1.8ยฐ, locally averaged velocity dispersions drop to the low levels see in the Local Volume (Karachentsev et al. 2003).
The properties demonstrated in Fig. 6 are as expected from a collapsed group. The $`1.8^{}`$ dimension can be inferred to mark the caustic of second turn-around (Bertschinger 1985). Objects bound to the group decouple from the cosmic expansion at a radius of first turn-around (the zero-velocity surface around the group) (Sandage 1986), collapse then reexpand to a second turnaround, then continue to oscillate while exchanging orbital energy with other group components. An observable cusp might be anticipated from the recent arrivals that have only had time to pass once through the group and just reach second turnaround. All other galaxies that have collapsed should lie interior (with rare slingshot exceptions).
If one gives consideration to other dynamically evolved groups in the Local Supercluster in a similar manner, one can infer a plausible second turnaround caustic by looking for an outer boundary to the distribution of early type systems. One of the authors (RBT) has undertaken this task, and has determined the following unpublished relation between the group velocity dispersion, $`\sigma _V`$, and the radius of the apparent caustic of second turnaround, $`r_{2t}`$, from observations of 7 dense groups and clusters in the Local Supercluster:
$$\sigma _V/r_{2t}=390\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1.$$
(3)
The conditions in the NGC 5846 Group are consistent with this relationship. Because each value of $`r_{2t}`$ can be used to define a value of $`\sigma _V(r_{2t})`$ using the velocity dispersion of the galaxies within $`r_{2t}`$, the above expression is a nonlinear equation in one variable. Solving the equation numerically, we find that 87 galaxies with velocities giving a dispersion of $`\sigma _V=320\mathrm{km}\mathrm{s}^1`$, defining a circle of projected radius $`r_{2t}=0.84`$ Mpc $`=1.84^{}`$, are included within the projected circle. This value of $`r_{2t}`$ is consistent with the radius at which the velocity dispersion profile levels off (Figure 6).
The $`r_{2t}`$ circle is superimposed on Fig. 2, with the center (at 226.40,+1.79) chosen to optimize the $`r_{2t}`$ enclosure. We conduct a virial analysis based on these 87 galaxies, with no luminosity weighting. We apply the median virial mass estimator (Heisler et al. 1985), which involves no luminosity weighing and does not require the determination of a velocity or spatial center. 10000 Monte Carlo simulations of the group are used to derive the errors on the estimated mass. We find the virial mass to be $`8.3\pm 0.3\times 10^{13}M_{}`$.
To further explore the group dynamics, we consider all 100 galaxies with known redshifts within 3ยฐ (1.8 Mpc) of NGC 5846, the galaxies plotted in Fig. 2. We construct a map of the mean velocity and velocity dispersion across the group. We define two functions on the sky: $`v_n(\alpha _{2000},\delta _{2000})`$ is the mean velocity of the $`n`$ closest group members to the position $`(\alpha _{2000},\delta _{2000})`$; $`\sigma _n(\alpha ,\delta )`$ is the velocity dispersion of the $`n`$ closest members to that position. We show $`v_{10}(\alpha _{2000},\delta _{2000})`$ and $`\sigma _{10}(\alpha _{2000},\delta _{2000})`$ along with the RASS X-ray contours in Figure 7.
The maps of $`v_{10}`$ and $`\sigma _{10}`$ are complex in character. There is considerable structure in the velocity dispersion across the face of the group. The regions associated with the bright galaxies NGC 5846, NGC 5813, and NGC 5831 show large velocity dispersions, $`500`$ km s<sup>-1</sup>, while the north-south region at $`\alpha _{2000}=`$ 15:03:30 (exactly between the NGC 5846 and NGC 5813 subgroups) shows a low velocity dispersion, $`300`$ km s<sup>-1</sup>. The spatial clumping and the velocity dispersion signatures suggest that NGC 5846 and NGC 5813 are the centers of distinct sub-structures. The persistence of this substructure is interesting because the characteristic crossing-time between these two centers (separation divided by mean group dispersion) is only 1.5 Gyr, $`10\%`$ the age of the universe.
The velocity map $`v_{10}`$ shows a general north-south trend in the mean local velocity of the galaxies: the galaxies to the north of the NGC 5846โNGC 5813 axis have lower velocities than those to the south, perhaps indicative of rotation. The complexity of the dynamical maps could not have been predicted from the featureless velocity histogram (Figure 1b). A single-peaked, regular velocity distribution is not necessarily indicative of a relaxed system. The system has almost certainly not reached virial equilibrium.
## 4. Properties of the Member Galaxies
### 4.1. Surface Brightness Scaling Relations
The surface brightness properties of the entire sample of confirmed group members and plausible candidates is demonstrated in Figure 8. The top panel shows the dependence with luminosity of the mean surface brightness within the radius containing half the light of a galaxy, the effective radius, while the bottom panel shows the dependence of the central magnitude through a 300 kpc radius metric aperture. In both plots there is a remarkably clean separation between high and low surface brightness systems. A modest separation is also found with morphological type in these plots. Types Sa and earlier are indicated by circles and lie slightly above types Sab and later, labeled with triangles.
The two plots carry similar information about an overall decrease in surface brightness in proceeding from giant galaxies to dwarfs. There is a small difference between the plots for systems toward the lower luminosity end of the high surface brightness group. If one splits the high surface brightness sample at $`M_R=20.5`$, it is seen that the fainter portion have higher mean surface brightnesses within an effective radius but marginally lower metric central magnitudes than the brighter portion. The increase in mean surface brightness is consistent with the scaling relation found by Kormendy (1977). Trujillo et al. (2001) point out that the ratio of central to mean surface brightness varies in a well correlated way with the Sรฉrsic (1968) parameterization of the radial distribution of light. The increase in mean surface brightness toward fainter luminosities reflects a trend in Sรฉrsic parameter. As for the apparent decrease in the metric central luminosity toward fainter though high surface brightness galaxies, partly this could be an artifact of resolution. The $`R_{300}`$ parameter, the magnitude within a radius of 300 pc = $`2.4^{\prime \prime }`$ was chosen to represent the central flux without constraining to a radius that would be affected by seeing. However, the galaxy N5846โ205 at $`M_R=17.5`$ (labeled 1 in the figure) has an effective radius of only 215 pc = $`1.6^{\prime \prime }`$ so the $`R_{300}`$ measure under-represents the central flux. With NGC 5846A and NGC 5845 (labeled 2 and 3 respectively), effective radii are $`500`$ pc. These unusual objects are given attention in the following section.
### 4.2. Small High Surface Brightness Objects
The galaxies that are identified by numbers in Fig. 8 are remarkable for their high central densities and small dimensions. They are N5846โ205 (1), NGC 5846A (2), and NGC 5845 (3). They only come to our attention through spectroscopic confirmation of an appropriate redshift. Their high surface brightnesses would exclude them from our sample of group candidates. Hence, our selection criteria demonstrably fails in this occasion occasion. Is this failure common or rare?
One could ask if there is a relationship with the ultracompact dwarfs that have been found in the Fornax Cluster (Phillipps et al. 2001). Those objects are faint ($`M_R12`$) and small ($`R_e20`$ pc), so small that they are indistinguishable from stars in ground-based imaging. They also only came to attention through spectroscopy. In the case of the NGC 5846 Group, SDSS spectroscopy has provided reasonable completion of non-stellar targets in the field that are brighter than $`R=17`$ ($`M_R=15`$). Hence, an ultracompact dwarf population like that found in Fornax would not be accessed with the current observations in the NGC 5846 region, because such dwarfs would be both too faint and too small.
Still, the SDSS spectroscopy does extend to a faintness limit that is interesting ($`M_R15`$) and precludes that there is a numerically important population of high surface brightness objects in the range $`15>M_R>19`$. Only one high surface brightness galaxy, N5846โ205, is found in this range. This system has properties similar to the Local Group elliptical M32, though N5846โ205, with $`M_R=17.5`$, is brighter by a magnitude. The other two high surface brightness objects that draw attention in Fig. 8, NGC 5846A and NGC 5845, are small in size but have luminosities that put them above the dwarf regime.
These three high surface brightness systems are very close to the core of the group! Their locations are seen in Figure 9. These three galaxies are the closest spectroscopically confirmed group members to NGC 5846, lying at 5, 25, and 55 kpc in projection (the other galaxies shown in blowup boxes in Figure 9 are low surface brightness candidate dwarfs without known velocities). This close proximity to the dominant NGC 5846 of the three most extreme objects in the surface brightness plot suggests very strongly that these objects have been tidally stripped. In this respect, these objects may be very large counterparts to the Fornax ultra compact dwarfs, suspected also to be the victims of tidal stripping or โthreshingโ (Bekki et al. 2003).
These objects are interesting in their own right but they are a distraction from the thread of the current investigation. The SDSS velocity information assures us that high surface brightness galaxies constitute only a tiny fraction of the group population in the interval $`15>M_R>19`$. There is no information, but no reason to suspect, that high surface brightness systems make up an important fraction of the population fainter than $`M_R=15`$.
### 4.3. Luminosity Function
The analysis in the previous sections has provided strong constraints on the group membership down to faint levels. The group is sufficiently populated that the domain of the region that has undergone collapse is reasonably defined ($`r_{2t}=0.84`$ Mpc). There is confirmation of membership for essentially all galaxies brighter than $`M_R=15`$ so the luminosity function at the bright end is quite secure (though statistics at high luminosities are limited). At fainter magnitudes the partial velocity information and the spatial correlation information provide good constraints on membership probabilities.
To constrain the luminosity function of the galaxies, we place objects in 0.5 mag bins and calculate the surface number density per unit absolute magnitude. The result appears in Figure 10. If there is only one galaxy in a bin, we combine that bin with an adjacent one. The luminosity function is fit to a faint limit of $`M_R=12`$. We assume that the dropoff faintward of that magnitude is due to incompleteness. The final luminosity function includes galaxies that are likely members according to the correlation analysis (ยง2.5): all galaxies confirmed by velocities as members, plus all galaxies rated 0โ2, plus 70% of galaxies rated 3 brighter than $`M_R=13.4`$, plus 50% of galaxies of early type rated 3 fainter than $`M_R=13.4`$. Because the number of galaxies per bin is small, we use the Gehrels (1986) asymptotic formula for calculating the error in a bin with $`N`$ members:
$$ฯต_N=1+(0.75+N)^{1/2}.$$
(4)
This formula is accurate to 1.5% for all values of $`N`$.
We find a striking feature: there is a plateau, or even dip, in the LF at $`M_R20`$. If real and not due to noise, this feature cannot be described by a simple Schechter (1976) function. A similar feature exists, with greater significance, in the data described by TT02. The Virgo cluster exhibits a quite significant and similar LF โbumpโ (Trentham & Hodgkin 2002).
We attempt to model this luminosity function using progressively more complicated fitting functions. To begin, in Figure 10a, we fit a simple power law in luminosity,
$$N(M)=N_{19}10^{0.4(\alpha +1)(M+19)},$$
(5)
Where $`N_{19}`$ is the number density of galaxies with $`M_R=19`$, and $`\alpha `$ is the slope defined so that it matches the faint-end slope of a Schechter (1976) function. We find this function to be an acceptable fit to the data (see Table 4). Given, however, that a single power law mass or luminosity distribution is not anticipated by either theoretical models of structure formation or observations, we also fit a Schechter (1976) function:
$`N(M)=`$ $`S(M)=`$ $`N_{}\mathrm{exp}[10^{0.4(MM_{})}]\times `$ (6)
$`10^{0.4(\alpha +1)(MM_{})}`$
This function also yields an acceptable fit. However, the characteristic magnitude $`M_{}`$ is constrained to be $`24.0\pm 3.2`$ at 95% confidence. This value is only marginally consistent with other analyses of the local galaxy luminosity function. For example, the overall SDSS luminosity function (Blanton et al. 2001) has $`M_{}=20.8\pm 0.3`$ at the 68% confidence level (the difference between the $`R`$ and the $`r`$ band photometry is negligible given our error bars). The value of $`\alpha =1.34\pm 0.06`$ is consistent with the overall SDSS value. Figure 10a shows both the single power law and the Schechter function fits to our data.
Because the single Schechter function fit yields somewhat too bright an $`M_{}`$, it is useful to evaluate other luminosity function models. The surface brightness-magnitude diagram (Figure 8) suggests a third way to proceed. There appears to be a significant gap in the surface brightness distribution in the group population. An ensemble of high surface brightness galaxies are clustered separately from the rest of the members (chiefly low-surface brightness dwarfs). We ask whether the two populations have separate luminosity distributions, similarly to what is seen in TT02 and Trentham & Hodgkin (2002). To test this idea, we fit these two populations separately by dividing them into groups with $`M_{R,300}<16`$ and $`M_{R,300}>16`$. In the high surface brightness population, one bin between $`M_R=19`$ and $`M_R=17`$ has zero members; the error in this empty bin is still described by the Gehrels (1986) formulation above (equation 4). The high surface brightness population is fit with a Gaussian,
$$N(M)=N_g\mathrm{exp}[\left(\frac{MM_g}{\sigma _g}\right)]^2,$$
(7)
while the low-surface brightness population is fit with a Schechter function. The results appear in Figure 10b and Table 4.
The Gaussian and Schechter functions provide a good fit to the LF of the two respective populations. The faint-end slope $`\alpha =1.23\pm 0.12`$ is consistent with the faint-end slope of the single Schechter function fit; the difference derives from the fact that the high-luminosity galaxies are excluded, leading to a much fainter $`M_{}=18.5\pm 1.3`$. The well-known correlation between $`M_{}`$ and $`\alpha `$ then leads to a smaller best-fit value of $`\alpha `$.
Our method for the inclusion of the priority 3 members may lead to systematic errors in $`\alpha `$. To constrain these uncertainties, we undertake two further fits. First, we measure $`\alpha `$ for only the spectroscopically confirmed plus the priority 1-2 members, obtaining $`\alpha =1.17\pm 0.15`$. Then we measure $`\alpha `$ for all the spectroscopically confirmed members plus the priority 1-3 members, obtaining $`\alpha =1.38\pm 0.10`$. Thus, we expect both the statistical and systematic errors in $`\alpha `$ to equal be $`0.1`$. Our final adopted $`\alpha `$ is an attempt to encapsulate all the complexities of the sample into a single value. We take the mean between the whole-sample and the low surface brightness fits: $`\alpha =1.3\pm 0.1`$ (statistical) $`\pm 0.1`$ (systematic).
Given the large number of free parameters and the small number of data points, the need for a Gaussian high-luminosity component is not statistically significant. However, a formulation more complicated than the Schechter function is required to describe any saddle between the giants and the dwarfs, a saddle that it seen recurrently in different samples, e.g. TT02 and Trentham & Hodgkin (2002). The present formulation may have some physical sense if the separation between high and low surface brightness systems seen in the scaling relations has some meaning. We have discussed surface brightness bimodality and a possible dynamical interpretation in the context of a sample dominated by late-type disk galaxies (Tully & Verheijen 1997). It is not clear if that discussion has relevance to this sample of predominantly early types. We do not argue that the present sample in itself provides justification for adopting a two-component luminosity function. The matter is something to review as material accumulates for more environments.
Finally, we wish to compare the luminosity function for this group with those found in other environments. To make this comparison we need a number density normalization. As a matter of convenience and to establish a convention, we determine a density in luminous galaxies ($`M_R<17`$) at a radius of 200 kpc from the group center. We are looking for variations in the bright/faint distributions which is why the group density definition is restricted to just the bright members. The metric radius of 200 kpc is chosen because it is representative of the core of our smallest groups. Details of the normalization procedure are discussed by TT02. Figure 11 presents the run of the density of galaxies with radius in the NGC 5846 Group in a fashion analogous to what is found in TT02. The fit to the density distribution of the luminous galaxies (open symbols) gives a normalization of 29 galaxies/Mpc<sup>2</sup> at a group radius of 200 kpc. The resulting renormalized LF, and a comparison with other groups, appears in Figure 12.
## 5. Summary
The long-term goal of this project is to explore the nature of the luminosity function of galaxies over a wide range of environments. In previous studies of environments in the Local Supercluster we have not had the areal coverage to encompass substantial fractions of the target groups. This deficiency is corrected in the present study of the NGC 5846 Group. It is seen in Fig. 2 that the area of our CFHT wide field survey covers almost the entire region subsequently considered to lie within the second turnaround cusp, $`r_{2t}`$. Dwarf galaxies as faint as $`M_R10`$ can be identified in the group, with incompletion setting in at $`M_R12`$.
Almost all galaxies associated with the group with $`M_R<15`$ are spectroscopically confirmed. We identify 324 probable or plausible candidates of which 83 have redshifts. Based on redshift sampling and a correlation analysis, we suggest that $`251\pm 10`$ of the 324 candidates are true members. The number of dwarfs in the NGC 5846 Group is very large. If, following TT02, a dwarf-to-giant ratio is defined as No. galaxies with $`11>M_R>17`$ over No. galaxies with $`M_R<17`$ then *dwarfs / giants* $`=7.3\pm 0.7`$. The error includes the uncertainty in the rating 3 memberships and incompletion near the $`M_R=11`$ limit. This ratio of dwarfs to giants is larger than the values seen in any of the TT02 groups. Figure 13 compares the percentage of dwarf elliptical galaxies ($`11>M_R>17`$) in the NGC 5846 Group with the TT02 groups. The NGC 5846 Group is overwhelmingly populated by early type dwarfs, with $`80\%\pm 5\%`$, comparable to the situation in the Virgo Cluster. About $`1/3`$ of the early-type dwarfs are nucleated.
The large dwarf to giant ratio is reflected in the relatively steep luminosity function at faint magnitudes. The simplified single power law slope of $`\alpha _d=1.34\pm 0.08`$ is significantly steeper than the mean slope for 6 groups of $`\alpha _d=1.19\pm 0.06`$ found by TT02 ($`95\%`$ probability). We also show the luminosity function found for the low density, spiral rich Ursa Major Cluster where the faint end slope is $`\alpha _d1.0`$. All of these observed luminosity functions are much shallower than the modified Press-Schechter mass function expected from the $`\mathrm{\Lambda }CDM`$ hierarchical clustering paradigm (Sheth & Tormen 1999).
The present observations taken with the earlier work already strongly support the proposition that the faint end of the luminosity function of galaxies varies with environment. Alternatively expressed, the ratio of dwarf to giant galaxies varies with environment. The dense, dynamically evolved NGC 5846 Group has a high dwarf/giant ratio and a relatively steep faint end luminosity function. Still, even in this environment there is a dearth of dwarfs compared to the expectations of the $`\mathrm{\Lambda }CDM`$ mass spectrum. There is the implication that astrophysical processes have affected the visible manifestations of low mass halos, and in ways that are more effective at suppression of light in lower density environments.
We thank the anonymous referee for insightful comments. This program involves observations with the Canada-France-Hawaii and Keck telescopes. It is supported by NSF award AST-03-07706. |
warning/0506/gr-qc0506037.html | ar5iv | text | # Hyperbolicity of second-order in space systems of evolution equations
## I Introduction
Research in numerical relativity has recently focused on obtaining a well-posed continuum initial-boundary value problem as a starting point for numerical time evolutions of systems such as a black-hole binary. Well-posedness of an initial-boundary value problem implies that an estimate
$$\delta u(,t)F(t)\left(\delta u(,0)+_0^t\delta g(,\tau )๐\tau \right)$$
(1)
exists, where $`u(x,t)`$ is the solution, $`u(x,0)`$ the initial data, $`g(x,t)`$ appropriate free boundary data, $`\delta `$ denotes a linear perturbation and $`||||`$ stands for appropriate norms (which may involve spatial derivatives), and where $`F(t)`$ is independent of the initial and boundary data. This means that the solution depends continuously on the initial and boundary data. Hyperbolicity is a property of the evolution equations that can be used as an algebraic criterion for well-posedness. We briefly review several notions of hyperbolicity.
Consider a system of quasilinear evolution equations that is first order in both space and time, or
$$\dot{u}=P^i(u)u_{,i}+S(u),$$
(2)
where $`u`$ is a vector of variables and $`P^i`$ are square matrices.
Definition 1: The system (2) is called weakly hyperbolic if the matrix $`P^nn_iP^i`$ has real eigenvalues for any unit vector $`n_i`$.
Definition 2: The system (2) is called strongly hyperbolic if $`P^n`$ is diagonalisable with real eigenvalues for any $`n_i`$, and the matrix $`T_n`$ that diagonalises it and its inverse $`T_n^1`$ depend smoothly on $`n_i`$.
Definition 3: The system (2) is called symmetric hyperbolic if there exists a Hermitian, positive definite matrix $`H`$ such that $`HP^n`$ is Hermitian for any direction $`n_i`$ and where $`H`$ does not depend on $`n_i`$.
The following properties of strongly and symmetric hyperbolic systems give a more practical meaning to the definitions, and we shall use them later to define strong and symmetric hyperbolicity for second-order systems. The key concept for strong hyperbolicity is
Definition 4: A characteristic variable with speed $`\lambda `$ in the $`n_i`$ direction is a linear combination $`๐ฎ`$ of the variables $`u`$ that obeys
$$_t๐ฎ=\lambda _n๐ฎ+\mathrm{},$$
(3)
where $`n_i`$ is normalised with respect to some metric, $`_nn^i_i`$, and the dots denote derivatives transverse to $`n_i`$ with respect to the same metric, and lower order terms.
If we write u as $`๐ฎ=\overline{u}^{}u`$, where $`\overline{u}`$ is a constant vector of coefficients, then
$`_t(\overline{u}^{}u)`$ $`=`$ $`\overline{u}^{}P^n_nu+\mathrm{}=\lambda _n(\overline{u}^{}u)+\mathrm{}`$ (4)
if and only if $`\overline{u}^{}`$ is a left eigenvector of $`P^n`$ or equivalently if $`\overline{u}`$ is an eigenvector of $`P^n`$. Characteristic variables $`๐ฎ`$ of the first-order reduction therefore correspond to left eigenvectors $`\overline{u}^{}`$ of $`P^n`$. This gives us
Lemma 1: A first-order system is strongly hyperbolic if and only if it admits a complete set of characteristic variables with real speeds that depend smoothly on $`n_i`$.
The key concept for symmetric hyperbolicity is that of an energy:
Definition 5: An energy $`ฯต`$ is a quadratic form in $`u`$ that is positive definite in the sense that $`ฯต=0`$ if and only if $`u=0`$, and which is conserved in the sense that there exists a flux $`\varphi ^i`$ quadratic in $`u`$ such that
$$\dot{ฯต}=\varphi _{}^{i}{}_{,i}{}^{}.$$
(5)
With
$$ฯตu^{}Hu,\varphi ^iu^{}HP^iu,$$
(6)
we have
Lemma 2: A linear first-order system with constant coefficients is symmetric hyperbolic if and only if it admits an energy.
For quasilinear systems this energy is conserved in the approximation where $`S(u)`$ is neglected and $`P^i(u)`$ is approximated as constant. (Physically, this corresponds to considering small high-frequency perturbations $`\delta u`$.) When boundaries are present, the time derivative of the energy can be estimated in terms of free boundary data.
Strong hyperbolicity of a first-order system is necessary and sufficient for a well-posed Cauchy problem. The Cauchy problem for a merely weakly hyperbolic system is typically ill-posed in the presence of lower-order terms. Symmetric hyperbolicity implies strong hyperbolicity, and is therefore also sufficient for well-posedness of the Cauchy problem. Furthermore, symmetric hyperbolicity can be used to prove well-posedness of the initial-boundary value problem for a certain class of boundary conditions called maximally dissipative GKO .
Hyperbolicity for equations or systems of equations of higher than first order is less well-established. A definition of weak hyperbolicity exists for systems of arbitrary order, but as for first-order systems, it does not guarantee well-posedness Beig . Alternatively, a quasilinear system that is second order in both space and time, or
$$P^{\mu \nu }(u,u)u_{,\mu \nu }+S(u,u)=0,$$
(7)
is called hyperbolic if $`P^{\mu \nu }`$ is a Lorentzian metric, that is if the principal part of the system is that of a wave equation Wald . Christodoulou has recently generalized the idea introducing the concept of regular hyperbolicity, with less strict positivity requirements on the elliptic block of the principal part reghyp . Both can be used as criteria for well-posedness of the Cauchy problem. The Einstein equations are second order, but they fit these definitions of hyperbolicity only when written in harmonic gauge, and so there are no standard definitions of hyperbolicity immediately applicable to forms of the Einstein equations commonly used in numerical relativity.
One possible approach to the well-posedness of a second-order system is to reduce it to first order by introducing auxiliary variables, and to define the second-order system to be strongly hyperbolic or symmetric hyperbolic if the reduction is. An ad-hoc definition along those lines has been used by Sarbach and co-authors SarbachBSSN ; BeyerSarbach to prove well-posedness of the BSSN formulation of the Einstein equations. We formalise this approach in Section III.
Independently, Nagy, Ortiz and Reula NOR , following Kreiss and Ortiz KreissOrtiz have used a pseudo-spectral reduction to define strong hyperbolicity. This method does not appear to generalise to symmetric hyperbolicity, intuitively because Fourier transforms cannot be carried out on a domain with arbitrary boundary. We briefly review this approach in Section V. By casting it in the notation of Section III, show that the two definitions of strong hyperbolicity are equivalent.
As a third alternative, Gundlach and Martรญn-Garcรญa bssn1 ; bssn2 define strong and symmetric hyperbolicity directly from the second-order system, by focusing on the existence of characteristic variables in strong hyperbolicity, and of an energy in symmetric hyperbolicity. We review this approach in Section IV, and show that its definitions of both strong and symmetric hyperbolicity are equivalent to those using a first-order reduction.
Outside the main line of this paper, we analyse in Section VI the well-posedness of the propagation of any constraints that the original second-order system is subject to. In Section VII we apply our results for symmetric hyperbolic systems to mixed symmetric hyperbolic-parabolic systems. Section VIII summarises our results.
## II The system
In this short Section, we establish notation for the class of system that we want to investigate, and clarify the relation between systems that are second order in both time and space, only in space, or only in space and that only in some of the variables.
We begin systems that are first order in time. Formulations of the Einstein equations based on the ADM formulation are naturally first order in time and second order in space. Some cannot even be written in second-order in time form (for example because they have an odd number of variables). With a non-zero shift, the first-order form may also be preferable for numerical simulations shiftedwave . For simplicity, we restrict attention to linear systems with constant coefficients. These can be considered as the linearisation and frozen coefficients approximation of a nonlinear system.
The class of systems that we are interested in are not uniformly second-order: some variables $`u`$ may appear in the evolution equations without second spatial derivatives. Writing $`u=(v,w)`$, where $`w`$ are those variables that appear only with first derivatives, we consider therefore
$`\dot{v}`$ $`=`$ $`A_1^{ij}v_{,ij}+A_1^iv_{,i}+A_1v+A_2^iw_{,i}+A_2w+a,`$ (8)
$`\dot{w}`$ $`=`$ $`B_1^{ij}v_{,ij}+B_1^iv_{,i}+B_1v+B_2^iw_{,i}+B_2w+b.`$ (9)
Here $`v`$ and $`w`$ are column vectors (not necessarily of the same length) of variables, the capital letters represent constant matrices of the appropriate dimension, and $`a`$, $`b`$ are forcing functions. We assume that while at least one second derivative of every variable $`v`$ appears in the equations, the number of variables $`v`$ has been minimised.
To reduce the system (8,9) to first order, we define the auxiliary variables $`d_iv_{,i}`$. By taking a spatial derivative of (8), we find
$$\dot{d}_i=A_1^{jk}v_{,ijk}+A_1^jv_{,ij}+A_1v_{,i}+A_2^jw_{,ij}+A_2w_{,i}+a_{,i}.$$
(10)
This is of a higher order than we started from, unless $`A_1^{ij}`$ and $`A_2^i`$ both vanish. (An example of a second-order evolution equation that cannot be reduced to first order is the heat equation $`\dot{u}=u^{\prime \prime }`$.) We have HinderPhD
Lemma 3: The general second-order in space, first-order in time linear system that can be reduced to first order by the introduction of auxiliary variables is of the form
$`\dot{v}`$ $`=`$ $`\underset{ยฏ}{A_1^iv_{,i}}+A_1v+\underset{ยฏ}{A_2w}+a,`$ (11)
$`\dot{w}`$ $`=`$ $`\underset{ยฏ}{B_1^{ij}v_{,ij}}+B_1^iv_{,i}+B_1v+\underset{ยฏ}{B_2^iw_{,i}}+B_2w+b.`$ (12)
From now on we refer to this as โtheโ second-order system. We have underlined the highest derivatives. Without loss of generality we assume from now on that $`B_1^{ij}`$ is symmetric.
In order to understand how general the system (11-12) is, it is interesting to convert it into second-order in both space and time form. Taking a time derivative of (11) and using (12) to replace $`\dot{w}`$, we obtain
$$\ddot{v}=A_2B_1^{ij}v_{,ij}+A_1^i\dot{v}_{,i}+A_2B_2^iw_{,i}+A_2B_2w+\mathrm{}$$
(13)
where we have written out all second derivatives and all appearances of $`w`$. We can eliminate the remaining appearances of $`w`$ and $`w_{,i}`$ in terms of $`\dot{v}`$ using (11) if and only if the matrices of the system obey
$$\mathrm{rank}(A_2)=\mathrm{rank}\left(\begin{array}{c}A_2\\ A_2B_2^i\end{array}\right)=\mathrm{rank}\left(\begin{array}{c}A_2\\ A_2B_2\end{array}\right).$$
(14)
When $`A_2`$ is invertible, which in particular implies equal numbers of $`v`$ and $`w`$ variables, these conditions are automatically obeyed. On the other hand, any fully second-order system in a set of variables $`v`$ can be reduced to the form (11-12) by introducing $`\dot{v}w`$. Therefore the class of first-order in time, second-order in space systems (11-12) includes the class of fully second-order systems, but is much bigger.
## III First-order reduction method
### III.1 Parameterised reduction
In reducing (11-12) to first order by defining $`d_iv_{,i}`$, we can write each occurrence of $`v_{,i}`$ also as $`d_i`$, or a mixture of the two, and similarly we can write $`v_{,ij}`$ as $`d_{i,j}`$ or $`d_{j,i}`$. To parameterise these ambiguities, we formally add multiples of the auxiliary constraint
$$c_id_iv_{,i}=0$$
(15)
and its antisymmetrised derivative
$$c_{ij}d_{[j,i]}=c_{[j,i]}=0$$
(16)
to all three equations. We could not add $`c_{(i,j)}`$, or any higher derivatives of the auxiliary constraints, without increasing the order of the system. The general reduction to first order is therefore
$`\dot{v}`$ $`=`$ $`A_1^iv_{,i}+A_1v+A_2w+a`$ (17)
$`+A_3^{ij}c_{ij}+A_3^ic_i,`$
$`\dot{w}`$ $`=`$ $`B_1^{ij}d_{i,j}+B_1^iv_{,i}+B_1v+B_2^iw_{,i}+B_2w+b`$ (18)
$`+B_3^{ij}c_{ij}+B_3^ic_i,`$
$`\dot{d}_i`$ $`=`$ $`A_1^jd_{i,j}+A_1v_{,i}+A_2w_{,i}+a_{,i}`$ (19)
$`+D_{i}^{}{}_{}{}^{k}c_k+D_{i}^{}{}_{}{}^{jk}c_{jk}.`$
From now on, we shall refer to this system as โtheโ reduction. We shall refer to the constant matrices $`A_3^i`$, $`A_3^{ij}`$, $`B_3^i`$, $`B_3^{ij}`$, $`D_{i}^{}{}_{}{}^{j}`$ and $`D_{i}^{}{}_{}{}^{jk}`$ as the reduction parameters. Without loss of generality, we assume that $`A_3^{ij}`$, $`B_3^{ij}`$ and $`D_{k}^{}{}_{}{}^{ij}`$ are antisymmetric in $`i`$ and $`j`$. The terms in the second-order system that become principal terms in the reduction are the ones underlined in (11-12). We shall call these the principal part of the second-order system.
### III.2 Auxiliary constraint evolution
The evolution of $`v`$ and $`d_i`$ implies an evolution of the auxiliary constraints $`c_i`$. This auxiliary constraint system can be written in first-order in space and time form by introducing $`c_{ij}c_{[j,i]}`$ (as already defined above) as auxiliary variables. This results in
$`\dot{c}_i`$ $`=`$ $`A_1^jc_{i,j}A_3^jc_{j,i}A_3^{jk}c_{jk,i}`$ (20)
$`+D_{i}^{}{}_{}{}^{j}c_j+D_{i}^{}{}_{}{}^{jk}c_{jk},`$
$`\dot{c}_{ij}`$ $`=`$ $`A_1^kc_{ij,k}D_{[i|}^{}{}_{}{}^{k}c_{k,|j]}D_{[i|}^{}{}_{}{}^{kl}c_{kl,|j]},`$ (21)
As the right-hand side of this system of linear PDEs is homogeneous, $`c(x,0)=0`$ implies $`\dot{c}(x,0)=0`$. Assuming that the coefficients are constant, as we do in this paper, on taking a Fourier transform in $`x^i`$ we obtain a separate ODE for each wavenumber $`\omega ^i`$, and it follows that $`c(x,t)=0`$ is the unique solution with $`c(x,0)=0`$.
We have shown that if the auxiliary constraints are zero initially, they remain zero at all times. This allows us to prove well-posedness for the reduction to first order, and then restrict it to the subset of solutions that obey the auxiliary constraints in order to infer well-posedness of the original second-order system.
Note that in order to make this argument, well-posedness of the auxiliary constraint system is not required, as we require only existence and uniqueness of the zero solution, not estimates of any non-zero solution.
### III.3 Definition of hyperbolicity
We now focus on the principal part of the first-order reduction,
$$_tuP^i_iu,$$
(22)
where here and in the following $``$ denotes equality up to lower-order terms and now $`u`$ stands for $`(v,w,d_i)`$. The degree of hyperbolicity of the first-order system depends on the the reduction parameters. The appropriate definitions of hyperbolicity are therefore the following:
Definition 1a: The second-order system (11-12) is defined to be weakly hyperbolic if and only if it admits at least one reduction to first order (17-19) that is weakly hyperbolic.
Definition 2a: The second-order system is defined to be strongly hyperbolic if and only if it admits at least one reduction to first order that is strongly hyperbolic.
Definition 3a: The second-order system is defined to be symmetric hyperbolic if and only if it admits at least one reduction to first order that is symmetric hyperbolic.
In the remainder of this Section we derive necessary and sufficient conditions for these definitions to hold, formulated directly in terms of the principal part of the second-order system, without reference to a reduction.
### III.4 2+1 split
We introduce a matrix notation in the groups of variables $`v`$, $`w`$ and $`d`$. In this notation, (22) is
$$_t\left(\begin{array}{c}v\\ w\\ d_j\end{array}\right)P^i_i\left(\begin{array}{c}v\\ w\\ d_k\end{array}\right)$$
(23)
where
$$P^i\left(\begin{array}{ccc}A_1^iA_3^i& 0& A_3^{ik}\\ B_1^iB_3^i& B_2^i& B_1^{ik}+B_3^{ik}\\ A_1\delta _{j}^{}{}_{}{}^{i}D_{j}^{}{}_{}{}^{i}& A_2\delta _{j}^{}{}_{}{}^{i}& A_1^i\delta _{j}^{}{}_{}{}^{k}+D_{j}^{}{}_{}{}^{ik}\end{array}\right).$$
(24)
We expand the tensor indices $`i`$, $`j`$ and $`k`$ in (24) into their components in the direction $`n_i`$ and transverse to it. For this purpose we need a positive definite metric tensor $`\gamma _{ij}`$. This can be chosen to be $`\delta _{ij}`$, or the physical 3-metric for example in applications to general relativity. We can then define the normal component of a tensor such as $`d_nn^id_i`$ where $`n^i\gamma ^{ij}n_j`$ and $`n_i`$ is normalised so that $`n_in_j\gamma ^{ij}1`$, and its transverse part such as $`(\delta _{i}^{}{}_{}{}^{j}n_in^j)d_j`$, which we denote by $`d_A`$. For a reason that will become apparent immediately, we also re-order the rows and columns. We call the resulting matrix $`๐ซ`$. It is related to $`P^n`$ by a unitary transformation that depends on $`n_i`$, but it is not the $`n_i`$ component of any vector of matrices. We have
$$_t\left(\begin{array}{c}w\\ d_n\\ v\\ d_A\end{array}\right)๐ซ_n\left(\begin{array}{c}w\\ d_n\\ v\\ d_B\end{array}\right)+\text{transverse derivatives},$$
(25)
where
$$๐ซ\left(\begin{array}{cccc}B_2^n& B_1^{nn}& B_1^nB_3^n& B_1^{nB}+B_3^{nB}\\ A_2& A_1^n& A_1D_{n}^{}{}_{}{}^{n}& D_{n}^{}{}_{}{}^{nB}\\ 0& 0& A_1^nA_3^n& A_3^{nB}\\ 0& 0& D_{A}^{}{}_{}{}^{n}& A_1^n\delta _{A}^{}{}_{}{}^{B}+D_{A}^{}{}_{}{}^{nB}\end{array}\right).$$
(26)
We write this in shorthand form as
$$๐ซ\left(\begin{array}{cc}๐& \\ 0& ๐\end{array}\right),๐\left(\begin{array}{cc}B_2^n& B_1^{nn}\\ A_2& A_1^n\end{array}\right).$$
(27)
The eigenvalues of $`๐ซ`$ are those of $`P^n`$, and one is diagonalisable if and only if the other is. Therefore we can investigate weak and strong hyperbolicity using $`๐ซ`$. The fact that $`๐ซ`$ has a zero block for all choices of the reduction parameters has several interesting consequences for its eigenvalues and eigenvectors, which respectively determine the weak and strong hyperbolicity properties of the first-order system.
### III.5 Eigenvalues of $`๐ซ`$: weak hyperbolicity
The first of these consequences is that the eigenvalues of $`๐ซ`$ (and hence of $`P^n`$) are the union of the eigenvalues of $`๐`$ and $`๐`$, and are independent of $``$. We will show in Section III.6 that we can choose the reduction parameters so that $`๐`$ has real eigenvalues. This gives us
Lemma 4: The second-order system is weakly hyperbolic according to Definition 1a if and only if $`๐`$ has real eigenvalues for all $`n_i`$.
We can further analyse the eigenvalues of $`๐`$. If we replace $`n_i`$ by $`n_i`$, then $`B_2^n`$ and $`A_1^n`$ change signs, while $`B_1^{nn}`$ and $`A_2`$ do not. This means that if $`\lambda `$ is an eigenvalue of $`๐`$ for some $`n_i`$ then $`\lambda `$ is an eigenvalue of $`๐`$ for $`n_i`$. The eigenvalues must therefore either be $`\lambda =O(n_i)`$ where $`O`$ is odd in $`n_i`$, or they are paired as $`\lambda _\pm =O(n_i)\pm E(n_i)`$ where $`O`$ is odd and $`E`$ is even. As a consequence, if the dimension of $`๐`$ is odd, at least one eigenvalue must be of the form $`\lambda =O(n_i)`$, and by continuity it must vanish for some $`n_i`$.
### III.6 Eigenvectors of $`๐ซ`$: strong hyperbolicity
The reduction is strongly hyperbolic if and only if $`P^n`$ (or equivalently $`๐ซ`$) is diagonalisable. In Appendix A we show that a necessary condition for $`๐ซ`$ to be diagonalisable is that both $`๐`$ and $`๐`$ are diagonalisable. If the sets of eigenvalues of $`๐`$ and $`๐`$ are disjoint, this is also sufficient. If they have any eigenvalues in common, additional necessary criteria arise which involve $``$.
To state these conditions in a simple form, diagonalise $`๐`$ and $`๐`$ simultaneously, so that
$$S^1๐ซS=\left(\begin{array}{cc}\mathrm{\Lambda }_๐& \stackrel{~}{}\\ 0& \mathrm{\Lambda }_๐\end{array}\right)$$
(28)
where $`\mathrm{\Lambda }_๐`$ and $`\mathrm{\Lambda }_๐`$ are diagonal matrices with the eigenvalues of $`A`$ and $`๐`$ respectively in the diagonal. Let all repeated eigenvalues be grouped together in these matrices. Then, for each eigenvalue common to $`๐`$ and $`๐`$, the corresponding block of $`\stackrel{~}{}`$ must vanish if $`๐ซ`$ is to be diagonalisable.
The block $`๐`$ does not contain any reduction parameters, and so is determined by the original second-order system. For a given direction $`n_i`$, blocks $``$ and $`๐`$ are determined completely by the choice of reduction parameters, but there are not enough reduction parameters to make this true for all directions $`n_i`$ at once. (For example, as $`B_1^{ij}`$ is symmetric and $`B_3^{ij}`$ is antisymmetric in $`ij`$, $`B_1^{nB}+B_3^{nB}`$ can be made to vanish for any one direction $`n_i`$, but not for all directions.)
We shall consider one choice of reduction parameters in discussing strong hyperbolicity (here) and another one in discussing symmetric hyperbolicity (in the next subsection). Both choices set
$`A_3^{ij}`$ $`=`$ $`0,`$
$`A_3^i`$ $`=`$ $`A_1^i,`$
$`B_3^i`$ $`=`$ $`B_1^i,`$
$`D_{i}^{}{}_{}{}^{j}`$ $`=`$ $`A_1\delta _{i}^{}{}_{}{}^{j},`$ (29)
This partial choice has the effect of decoupling $`v`$ from the $`w`$ and $`d_i`$. To discuss strong hyperbolicity in the case of three space dimensions, we complete the choice of reduction parameters by
$`B_3^{ij}`$ $`=`$ $`0,`$
$`D_{j}^{}{}_{}{}^{ik}`$ $`=`$ $`\delta _j{}_{}{}^{i}A_{1}^{k}\delta _j{}_{}{}^{k}A_{1}^{i}+i\mu ฯต_{j}^{}{}_{}{}^{ik},`$ (30)
where $`\mu `$ is a real constant, $`ฯต_{ijk}`$ is the totally antisymmetric tensor in three dimensions and $`i=\sqrt{1}`$. This gives
$$๐ซ=\left(\begin{array}{cccc}B_2^n& B_1^{nn}& 0& B_1^{nB}\\ A_2& A_1^n& 0& A_1^B\\ 0& 0& 0& 0\\ 0& 0& 0& i\mu ฯต_A^{nB}\end{array}\right).$$
(31)
The diagonal block $`i\mu ฯต_A^{nB}`$ is diagonalisable with real eigenvalues $`\pm \mu `$. The complexification is unusual, but the term multiplied by $`i\mu `$ is proportional to the auxiliary constraint $`c_{ij}`$, and so has no influence on the original second-order system. If we choose $`\mu `$ large enough, the eigenvalues of this block are distinct from those of the complementary diagonal block (containing $`๐`$ and a zero row and column), and we have found a choice of reduction parameters that makes $`๐ซ`$ diagonalisable if $`๐`$ is diagonalisable. The existence of this choice of reduction parameters completes our proof of
Theorem 1: The second-order system is strongly hyperbolic according to Definition 2a if and only if $`๐`$ is diagonalisable for all $`n_i`$ where the diagonalising matrix and its inverse depend smoothly on $`n_i`$.
Our proof of this theorem assumes three space dimensions, but we suspect that the theorem holds in any number of space dimensions.
Note that this criterion is based only on the coefficients on the second-order system, without explicit reference to the reduction. Note also that for the choice of reduction parameters we have used here, the auxiliary constraint system is strongly hyperbolic, see Appendix B.
### III.7 Symmetric hyperbolicity
According to Definition 3, the reduction is symmetric hyperbolic if and only if there is a Hermitian matrix $`H`$ such that
$$(HP^n)^{}=HP^n$$
(32)
for all $`n_i`$, with $`H`$ independent of $`n_i`$ and positive definite. Note that in the definition (32) we cannot replace $`P^n`$ by $`๐ซ`$.
We again make the partial choice (29) of reduction parameters. $`B_3^{ij}`$ and $`D_{i}^{}{}_{}{}^{jk}`$ will be determined in the following in terms of $`H`$. The resulting form of $`P^i`$ is \[see (23) for the definition of $`P^i`$\]
$$P^i=\left(\begin{array}{ccc}0& 0& 0\\ 0& B_2^i& B_1^{ik}+B_3^{ik}\\ 0& A_2\delta _{j}^{}{}_{}{}^{i}& A_1^i\delta _{j}^{}{}_{}{}^{k}+D_{j}^{}{}_{}{}^{ik}\end{array}\right).$$
(33)
Clearly it is sufficient to find a symmetriser for the $`(w,d_i)`$ block. We parameterise $`H`$ as
$$u^{}Hu=(v^{},w^{},d_m^{})\left(\begin{array}{ccc}1& 0& 0\\ 0& K& L^j\\ 0& L^m& M^{mj}\end{array}\right)\left(\begin{array}{c}v\\ w\\ d_j\end{array}\right)$$
(34)
with $`K`$ and $`M`$ Hermitian and positive, $`K^{}=K`$, $`K>0`$, and $`M^{mj}=M^{jm}`$, $`M>0`$.
The nontrivial, $`(w,d_i)`$ block of $`HP^i`$ is
$$\left(\begin{array}{cc}KB_2^i+L^iA_2& K(B_1^{ik}+B_3^{ik})\\ & +L^kA_1^i+L^jD_{j}^{}{}_{}{}^{ik}\\ L^mB_2^i+M^{mi}A_2& L^m(B_1^{ik}+B_3^{ik})\\ & +M^{mk}A_1^i+M^{mj}D_{j}^{}{}_{}{}^{ik}\end{array}\right).$$
(35)
A necessary condition for this to be Hermitian is for the matrix
$$\left(\begin{array}{cc}KB_2^n+L^nA_2& KB_1^{nn}+L^nA_1^n\\ L^nB_2^n+M^{nn}A_2& L^nB_1^{nn}+M^{nn}A_1^n\end{array}\right)$$
(36)
to be Hermitian for all $`n_i`$. This is just the condition that $`๐`$ admits a symmetriser, or
$$๐=(๐)^{},\text{where}\left(\begin{array}{cc}K& L^n\\ L^n& M^{nn}\end{array}\right),$$
(37)
for all $`n_i`$. If the system is strongly hyperbolic, $`๐`$ always admits a symmetriser formed from its eigenvectors. Nevertheless, (37) is a non-trivial condition because $``$ must form a part of $`H`$, so that its blocks $`L^n`$ and $`M^{nn}`$ are given by $`L^in_i`$ and $`M^{ij}n_in_j`$, where $`K`$, $`L^i`$ and $`M^{ij}`$ do not depend on $`n_i`$. We now show that (37) actually implies that all of (35) is Hermitian for a particular choice of reduction parameters $`B_3^{ij}`$ and $`D_{i}^{}{}_{}{}^{jk}`$. We do this for each block of (35) in turn.
From the top left block of (35) we have
$$(KB_2^i+L^iA_2)^{}=KB_2^i+L^iA_2.$$
(38)
This equation does not contain any reduction parameters, and it is clearly equivalent to the condition of $`KB_2^n+L^nA_2`$ being Hermitian for all $`n_i`$, which is contained in the first diagonal block of (37).
The off-diagonal blocks of (35) give
$$(L^kB_2^i+M^{ki}A_2)^{}=K(B_1^{ik}+B_3^{ik})+L^kA_1^i+L^jD_{j}^{}{}_{}{}^{ik}.$$
(39)
If we denote this equation by $`T^{ik}=0`$, then its symmetric part $`T^{(ik)}=0`$, or
$$B_2^{(i}L^{k)}+A_2^{}M^{(ik)}=KB_1^{ik},+L^{(k}A_1^{i)},$$
(40)
again does not contain any reduction parameters. Furthermore $`T^{(ik)}=0`$ if and only if $`T^{ik}n_in_k=0`$ for all $`n_i`$ (this is a property of all totally symmetric tensors), and this is precisely the condition contained in the off-diagonal terms of (37). The antisymmetric part $`T^{[ik]}=0`$, or
$$B_2^{[i}L^{k]}+A_2^{}M^{[ik]}=KB_3^{ik}+L^jD_{j}^{}{}_{}{}^{ik},$$
(41)
can be solved for $`B_3^{ij}`$ because $`K`$ is by assumption invertible.
Finally, the bottom right block of (35) gives
$`L^m(B_1^{ik}+B_3^{ik})+M^{mk}A_1^i+M^{mj}D_{j}^{}{}_{}{}^{ik}`$
$`=\left[L^k(B_1^{im}+B_3^{im})+M^{km}A_1^i+M^{kj}D_{j}^{}{}_{}{}^{im}\right]^{}.`$ (42)
If we write this as $`T^{mik}=0`$, the totally symmetric part $`T^{(mik)}=0`$, or
$`L^{(m}B_1^{ik)}+M^{(mk}A_1^{i)}=\left[L^{(k}B_1^{im)}+M^{(km}A_1^{i)}\right]^{}.`$ (43)
does not contain any reduction parameters. It is equivalent to $`T^{mik}n_mn_in_k=0`$ for all $`n_i`$, and so vanishes because of the last diagonal block of (37). After solving (41) for $`B_{3}^{}{}_{}{}^{ij}`$ we write (III.7) as
$$U^{mik}+\overline{M}^{mj}D_{j}^{}{}_{}{}^{ik}=U^{kim}+\left(\overline{M}^{kj}D_{j}^{}{}_{}{}^{im}\right)^{},$$
(44)
where we have defined
$`U^{mik}`$ $``$ $`L^mB_1^{ik}+M^{mk}A_1^i`$ (45)
$`+L^mK^1\left(B_2^{[i}L^{k]}+A_2^{}M^{[ik]}\right),`$
$`\overline{M}^{ij}`$ $``$ $`M^{ij}L^iK^1L^j.`$ (46)
$`\overline{M}`$ is Hermitian ($`\overline{M}^{ij}=\overline{M}^{ji}`$) and positive definite \[because $`H`$ is positive definite in particular when restricted to vectors $`(v,w,d_i)=(0,K^1L^kd_k,d_i`$)\], and hence invertible. We define
$$X^{mik}U^{kim}U^{mik},D^{mik}\overline{M}^{mj}D_{j}^{}{}_{}{}^{ik},$$
(47)
and so
$$X^{mik}=X^{kim},$$
(48)
while $`T^{(mik)}=0`$ is equivalent to
$$X^{(mik)}=0.$$
(49)
(44) becomes
$$D^{mik}D^{kim}=X^{mik},$$
(50)
which has the general solution
$`6D^{mik}`$ $`=`$ $`X^{kmi}+2X^{kim}+3X^{ikm}`$ (51)
$`+4X^{imk}+5X^{mik}+Y^{mik},`$
using both (48) and (49). The object $`Y^{mik}`$ must obey
$$Y^{mik}=Y^{[mik]},Y^{mik}=Y^{mik},$$
(52)
but is otherwise arbitrary. It parameterises the part of the $`D_{j}^{}{}_{}{}^{ik}`$ that is not determined by $`H`$, and can be set to zero.
We have shown that a necessary and sufficient condition for $`P^i`$ to admit a symmetriser $`H`$ is for $`๐`$ to admit a symmetriser $``$ that depends on $`n_i`$ in the particular way given in (37). Therefore we have
Theorem 2: A necessary and sufficient condition for the second-order system to by symmetric hyperbolic according to Definition 3a is that (37) holds for all $`n_i`$ for some $`H>0`$ parameterised by (34).
Note that positive definiteness of $``$ does not imply that of $`H`$, because positive definiteness of $`M^{nn}`$ for all $`n_i`$ does not imply positive definiteness of $`M`$. The difference can be expressed in standard terminology as follows: a double quadratic form $`M_{}^{ij}{}_{AB}{}^{}`$ is rank-1 positive if and only if $`M_{}^{ij}{}_{AB}{}^{}n_in_jm^Am^B>0`$ for all $`n_i`$ and $`m^A`$; it is rank-2 positive if and only if $`M_{}^{ij}{}_{AB}{}^{}n_i^An_j^B>0`$ for all $`n_i^A`$. Rank-2 positivity implies rank-1 positivity, but they are not equivalent when the indices $`i,j,A,B`$ belong to spaces of dimension 3 or larger, as shown in the example given in Ball . This suggests it may be useful to introduce an intermediate concept of hyperbolicity based on positivity of $``$ for all $`n_i`$, rather than positivity of $`H`$. This already guarantees that $`๐`$ is diagonalisable, so that the system is strongly hyperbolic. The imposition of rank-1 positivity is also the key ingredient of the definition of โregular hyperbolicityโ by Christodoulou reghyp for second order in both space and time systems, which is also known to yield well-posed problems. We believe there is a connection between regular hyperbolicity and our condition $`>0`$, but have not been able to show it.
## IV Direct second-order method
### IV.1 Strong hyperbolicity
Following bssn1 , we now elevate Lemma 1 to a definition of strong hyperbolicity for second-order in space, first-order in time systems, with the only difference that the $`๐ฎ`$ become linear combinations of $`v_{,i}`$ and $`w`$:
Definition 2b: The system (11-12) is called strongly hyperbolic if and only if is there is a set of characteristic variables $`๐ฎ`$ linear in $`w`$ and $`v_{,i}`$ obeying (3) with real speeds $`\lambda `$ and where the matrix relating $`u`$ and $`๐ฎ`$ and its inverse depend smoothly on $`n_i`$.
This definition has two interesting consequences. The first is that $`v_{,A}`$ can always be considered as a zero speed variable in the direction $`n_i`$ because
$$_t(v_{,A})=_A(_tv)$$
(53)
so that the right-hand side can be considered as a sum of transverse derivatives only. (In fact arbitrary linear combinations of the $`v_{,A}`$ can be given arbitrary speeds.)
The second consequence is that nontrivial characteristic variables are unique only up to adding transverse derivatives. If $`๐ฎ`$ is a characteristic variable with speed $`\lambda `$ then, for any vector $`X^A`$ (made from the $`u`$)
$`_t(๐ฎ+_AX^A)`$ (54)
$``$ $`_t๐ฎ+_A(_tX^A)`$
$``$ $`\lambda _n๐ฎ+\text{transv. deriv.}`$
$``$ $`\lambda _n๐ฎ+\lambda _A(_nX^A)+\text{transv. deriv.}`$
$``$ $`\lambda _n(๐ฎ+_AX^A)+\text{transv. deriv.}`$
and so $`๐ฎ+_AX^A`$ is also a characteristic variable with speed $`\lambda `$. Such calculations rely on commuting partial derivatives to interpret $`_A_n\mathrm{}`$ either as a normal or a transverse derivative, depending on the situation.
The second-order system has no reduction parameters. However, for the purpose of comparing the second-order approach with the reduction approach, we can translate the second-order approach into the notation of a first-order reduction. We account for the fact that $`v_{,i}`$ and $`d_i`$, and $`v_{,ij}`$ and $`v_{,ji}`$, are now identical, by allowing the reduction parameters to depend explicitly on $`n_i`$ (so that they are not tensors.) This allows us to set the blocks $``$ and $`๐`$ of $`๐ซ`$ to arbitrary values for every $`n_i`$, and this allows us to make arbitrary linear combinations of $`v_{,A}`$ and $`v`$ characteristic variables with arbitrary speeds, and to add arbitrary combinations of $`v_{,A}`$ and $`v`$ to any characteristic variables made from $`w`$ and $`v_{,n}`$, as we have discussed above. $`๐`$ being diagonalisable is a necessary condition for $`๐ซ`$ to be diagonalisable, and with $`=0`$ and $`๐=0`$ it is also sufficient. Therefore Definition 2b is equivalent to $`๐`$ being diagonalisable and by Theorem 1 it is then equivalent to Definition 2a.
Alternatively, we can write down the principal part of the second-order system in a 2+1 split as follows:
$`\dot{v}`$ $``$ $`0,`$ (55)
$`\dot{v}_{,n}`$ $``$ $`A_1^j(v_{,n})_{,j}+A_2w_{,n},`$ (56)
$`\dot{v}_{,A}`$ $``$ $`A_1^j(v_{,j})_{,A}+A_2w_{,A},`$ (57)
$`\dot{w}`$ $``$ $`B_1^{nn}(v_{,n})_{,n}+2B_1^{nA}(v_{,n})_{,A}`$ (58)
$`+B_1^{AB}(v_{,A})_{,B}+B_2^iw_{,i},`$
where $`v_{,i}`$ on the right-hand side is now considered a lower order term. Note the way the second derivatives have been ordered differently in $`\dot{v}_{,n}`$ and $`\dot{v}_{,A}`$. In the language of reduction this corresponds to the reduction parameters depending explicitly on $`n_i`$. The matrix $`๐ซ`$ becomes
$$_t\left(\begin{array}{c}w\\ v_{,n}\\ v\\ v_{,A}\end{array}\right)๐ซ_n\left(\begin{array}{c}w\\ v_{,n}\\ v\\ v_{,B}\end{array}\right)+\text{transverse derivatives},$$
(59)
where
$$๐ซ\left(\begin{array}{cccc}B_2^n& B_1^{nn}& 0& 0\\ A_2& A_1^n& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right).$$
(60)
Once again $`๐ซ`$ is diagonalisable if and only if $`๐`$ is. Either way, and taking into account the smoothness conditions, we have
Theorem 3: Definition 2b is equivalent to Definition 2a.
### IV.2 Symmetric hyperbolicity
Following bssn1 , we now elevate Lemma 2 to a definition of symmetric hyperbolicity for second-order systems in
Definition 3b: The second-order system is called symmetric hyperbolic if and only if it admits a positive definite energy $`ฯต`$ and a flux $`\varphi ^i`$, both quadratic in $`w`$ and $`v_{,i}`$, that obey
$$\dot{ฯต}\varphi _{}^{i}{}_{,i}{}^{}.$$
(61)
Note that in obtaining (61) one can make arbitrary use of $`v_{,ij}=v_{,ji}`$.
We parameterise the energy $`ฯต`$ by (34), with $`d_i`$ replaced by $`v_{,i}`$, and parameterise the flux $`\varphi ^i`$ as
$$\varphi ^i=(w^{},v_{,m}^{})\left(\begin{array}{cc}F^i& F^{ik}\\ F^{im}& F^{mik}\end{array}\right)\left(\begin{array}{c}w\\ v_{,k}\end{array}\right),$$
(62)
with $`F^i=F^i`$ and $`F^{kim}=F^{mik}`$. In the second-order system, there are no reduction parameters, and so the non-trivial part of $`HP^i`$ is given by
$$\left(\begin{array}{cc}KB_2^i+L^iA_2& KB_1^{ik}+L^kA_1^i\\ L^mB_2^i+M^{mi}A_2& L^mB_1^{ik}+M^{mk}A_1^i\end{array}\right).$$
(63)
In a first-order system, energy conservation is precisely equivalent to $`HP^i`$ being Hermitian. In the second-order system, the relation between $`HP^i`$ and energy conservation is more complicated, because $`v_{,ik}=v_{,ki}`$.
To see this more clearly, we write out
$`\dot{ฯต}`$ $`=`$ $`2[w^{}(KB_2^i+L^iA_2)w_{,i}`$ (64)
$`+w^{}(KB_1^{ik}+L^{(k}A_1^{i)})v_{,ik}`$
$`+v_{,m}^{}(L^mB_2^i+M^{mi}A_2)w_i`$
$`+v_{,m}^{}(L^mB_1^{ik}+M^{m(k}A_1^{i)})v_{,ik}],`$
$`\varphi ^i_{,i}`$ $`=`$ $`2[w^{}F^iw_{,i}+w^{}F^{(ik)}v_{,ik}`$ (65)
$`+v_{,m}^{}F^{im}w_i+v_{,m}^{}F^{m(ik)}v_{,ik}].`$
Comparing the first terms in $`\dot{ฯต}`$ and $`\varphi ^i_{,i}`$ we have
$`F^i`$ $`=`$ $`KB_2^i+L^iA_2,`$ (66)
$`F^i`$ $`=`$ $`F^i,`$ (67)
which admits a solution $`F^i`$ if and only if (38) is obeyed. Comparing the second and third terms we have
$`F^{(ik)}`$ $`=`$ $`KB_1^{ik}+L^{(k}A_1^{i)},`$ (68)
$`F^{im}`$ $`=`$ $`L^mB_2^i+M^{mi}A_2,`$ (69)
which admits a solution $`F^{ik}`$ if and only if (40) is obeyed. Comparing the fourth terms we have
$`F^{m(ik)}`$ $`=`$ $`L^mB_1^{ik}+M^{m(k}A_1^{i)}S^{mik},`$ (70)
$`F^{kim}`$ $`=`$ $`F^{mik}.`$ (71)
These admit a solution $`F^{mik}`$ if and only if (43) is obeyed. It is clear that (43) is necessary. To demonstrate that it is sufficient we solve explicitly for $`F^{mik}`$. The general solution of (70) is
$$F^{mik}=S^{mik}+J^{mik},$$
(72)
where $`J^{mik}=J^{m[ik]}`$. Defining
$$V^{mik}S^{kim}S^{mik},$$
(73)
the general solution of (71) is
$`6J^{mik}`$ $`=`$ $`V^{kmi}+2V^{kim}+3V^{ikm}`$ (74)
$`+4V^{imk}+5V^{mik}+W^{mik},`$
if and only if $`V^{(mik)}=0`$, which is equivalent to (43). The remaing free coefficient $`W^{mik}W^{[mik]}W^{mik}`$ parameterises terms in $`\varphi ^i`$ whose divergence vanishes identically. It can be set equal to zero without loss of generality.
We have shown that the second-order system admits a conserved energy if and only if (38), (40) and (43) hold, which together are equivalent to (37). This is equivalent to the existence of a conserved energy for the first-order reduction. The two energies are in fact the same under the (unambiguous) identification of $`d_i`$ with $`v_{,i}`$. This means that the two definitions of symmetric hyperbolicity are equivalent, as both are equivalent to the matrix $`H`$ defined by (34) being positive definite, and obeying (37) for all $`n_i`$. We have
Theorem 4: Definition 3b is equivalent to Definition 3a.
Given an energy $`H`$ of the second-order system, a first-order reduction that admits the same energy is given by the reduction parameters (29) and $`B_3^{ij}`$ and $`D_{j}^{}{}_{}{}^{ik}`$ determined in Section III.7. Going the other way, an energy $`H`$ for any first-order reduction is clearly also an energy for the original second-order system. If the second-order system admits a multi-parameter family of energies, then some of these parameters define reduction parameters $`B_3^{ij}`$ and $`D_{j}^{}{}_{}{}^{ik}`$ of the first-order system, and the remainder parameterise the energy of that particular reduction. An example of this split is given in Appendix D.
### IV.3 Maximally dissipative boundary conditions
A first-order symmetric hyperbolic system on a domain $`\mathrm{\Omega }`$ admits an energy
$$E=_\mathrm{\Omega }ฯต๐V$$
(75)
whose time derivative is given by the flux through the boundary,
$$\dot{E}_\mathrm{\Omega }\varphi ^n๐S,$$
(76)
where $`n_i`$ is now the outward pointing normal to the boundary $`\mathrm{\Omega }`$. In Appendix C we show that symmetric hyperbolicity implies strong hyperbolicity and that the boundary flux can be written as
$$\varphi ^n=\underset{\alpha }{}\lambda _\alpha ๐ฎ_\alpha ^2.$$
(77)
where the sum is over a suitable basis of characteristic variables. Therefore the growth of the energy can be controlled by controlling all characteristic variables that are ingoing at the boundary ($`\lambda _\alpha >0`$) (โmaximally dissipative boundary conditionsโ), while ingoing or tangential characteristic variables give negative or zero contributions to $`\varphi ^n`$ and hence $`\dot{E}`$.
The same result holds for a second-order system but, as we have seen, the characteristic variables of the second-order system are not unique: the $`v_{,A}`$ can be given arbitrary speeds, and arbitrary combinations of $`v_{,A}`$ can be added to any characteristic variable. In order to impose maximally dissipative boundary conditions with the desired effect of controlling an energy, we need to control all incoming characteristic variables of an actual first-order reduction. We can do this within the second-order system by imposing boundary conditions on characteristic variables $`๐ฎ_\alpha `$ constructed from $`w`$ and $`v_{,i}`$, but we need to have the correct admixtures of $`v_{,A}`$ in these characteristic variables.
We have seen that in order to show strong hyperbolicity for the second-order system one needs to diagonalise only the matrix $`๐`$. Let us call the characteristic variables of the second-order system given by eigenvectors of $`๐^{}`$, and which are therefore constructed only from by $`w`$ and $`v_{,n}`$, the โshortโ characteristic variables, and let us denote them by $`๐ฎ_\alpha ^{}`$. For proving symmetric hyperbolicity one only needs to find a conserved energy $`ฯต(w,v_{,i})`$.
Assuming that we already have the $`๐ฎ_\alpha ^{}`$ on the one hand, and the energy $`ฯต(w,v_{,i})`$ on the other, the simplest way of generating the required โfullโ characteristic variables $`๐ฎ_\alpha `$ is to make the ansatz $`๐ฎ_\alpha =\mathrm{const}(๐ฎ_\alpha ^{}`$ \+ undetermined multiples of the $`d_A)`$, and to determine the coefficients and overall normalisation for each $`๐ฎ_\alpha `$ so that the full characteristic variables obey
$$ฯต=\underset{\alpha }{}๐ฎ_\alpha ^2.$$
(78)
(76) then follows.
## V Pseudo-differential reduction method
For the purpose of comparison, we now describe the pseudo-differential reduction method of KreissOrtiz ; NOR in our notation. We carry out a Fourier transform in space with wave number $`\omega _i`$ of the second-order system (11-12). We denote the Fourier transforms of $`v`$ and $`w`$ by $`\widehat{v}`$ and $`\widehat{w}`$. We choose the direction $`n_i`$ to be that of the wavenumber $`\omega _i`$ and write $`\omega _i|\omega |n_i`$. We then have $`\widehat{d}_n=i|\omega |\widehat{v}`$ and $`\widehat{d}_A=0`$. We can then use $`\widehat{d}_n`$ to represent $`\widehat{v}`$. The principal part of the Fourier-transformed system can be written as
$$_t\left(\begin{array}{c}\widehat{w}\\ \widehat{d}_n\end{array}\right)i|\omega |๐\left(\begin{array}{c}\widehat{w}\\ \widehat{d}_n\end{array}\right)$$
(79)
where $`๐`$ is the matrix defined above, and the non-principal terms not shown here are homogeneous of order zero in $`|\omega |`$. In our notation, the definition of KreissOrtiz is then
Definition 2c: The second-order system is called strongly hyperbolic if and only if there exists a pseudo-differential reduction to first order of the form (79) where $`๐`$ is diagonalisable with real eigenvalues, and the diagonalising matrix and its inverse depend smoothly on $`n_i`$.
We have set up our notation so that from Theorem 1 we immediately have
Lemma 5: Definition 2c is equivalent to Definition 2a and to Definition 2b.
As the pseudo-differential approach relies in an essential way on Fourier transforms, it does not lend itself to defining a locally conserved energy. Therefore there is no definition of symmetric hyperbolicity in this approach.
## VI The evolution of constraints on the original second-order system
In many applications, the original second-order system is subject to differential constraints, which are conserved under evolution. We shall call these the โoriginal constraintsโ to distinguish them from the โauxiliary constraintsโ $`d_iv_{,i}=0`$ that arise additionally in the process of first-order reduction.
Note that until now we have not mentioned or used original constraints. The reason is that in general one wants to prove well-posedness of the second-order system if the original constraints are obeyed or not. This is important for example if one wants to carry out numerical simulations using โfree evolutionโ where the original constraints are imposed only on the initial data. In the continuum the constraints then remain zero, but in numerical free evolution they are violated through finite differencing error. At a later time one is effectively evolving initial data that do not obey the constraints. Therefore the continuum problem must be well-posed for non-vanishing original constraints as a necessary condition for numerical stability. One may of course use the original constraints to modify the original, second-order system, but here we assume that this has already been done, and the second-order system is fixed.
We now prove that the evolution system of original constraints is strongly hyperbolic if it closes and if the second-order main system is strongly hyperbolic. An equivalent result for first-order systems subject to constraints is given in Reulastrong .
We consider a vector $`c`$ of constraints which are quasilinear of the form
$$cC_1^{ij}v_{,ij}+C_2^iw_{,i},$$
(80)
where the matrix $`C^{ij}`$ is symmetric in $`ij`$. It is clear from (11-12) that the evolution of these constraints is first order in space and time. If the constraint system is closed, its principal part must then be of the form
$$\dot{c}G^ic_{,i}$$
(81)
for a vector of square matrices $`G^i`$. Using the second-order evolution equations (11-12) and comparing the leading order terms in $`v`$ and $`w`$ we find
$`C_1^{(ij}A_1^{k)}+C_2^{(i}B_1^{jk)}`$ $`=`$ $`G^{(i}C_1^{jk)},`$ (82)
$`C_1^{(ij)}A_2+C_2^{(i}B_2^{j)}`$ $`=`$ $`G^{(i}C_2^{j)}.`$ (83)
These identities between totally symmetric matrices hold if and only if their contraction with $`n_i`$ on all indices hold for all $`n_i`$. Writing the $`nnn`$ and $`nn`$ components of these equations in matrix form we have
$$(C_2^n,C_1^{nn})\left(\begin{array}{cc}B_2^n& B_1^{nn}\\ A_2& A_1^n\end{array}\right)=G^n(C_2^n,C_1^{nn})$$
(84)
for all $`n_i`$. We write this in compact notation as
$$C๐=GC.$$
(85)
If the second-order main system is strongly hyperbolic, $`๐`$ is diagonalisable with $`๐=T\mathrm{\Lambda }T^1`$. $`G`$ can always be brought into Jordan form as $`G=SJS^1`$. Then
$$\stackrel{~}{C}\mathrm{\Lambda }=J\stackrel{~}{C},\stackrel{~}{C}S^1CT.$$
(86)
We assume that the rows of $`C`$, and therefore the rows of $`\stackrel{~}{C}`$, are linearly independent. This means that there is no redundancy between the differential constraints, and is similar to an assumption in Reulastrong .
Consider now the first Jordan block of $`J`$ with eigenvalue $`\mu _1`$. For simplicity assume it has size 2. Exceptionally writing out the internal matrix indices, we have
$`\stackrel{~}{C}_{1\alpha }\lambda _\alpha `$ $`=`$ $`\mu _1\stackrel{~}{C}_{1\alpha }+\stackrel{~}{C}_{2\alpha },`$
$`\stackrel{~}{C}_{2\alpha }\lambda _\alpha `$ $`=`$ $`\mu _1\stackrel{~}{C}_{2\alpha }`$ (87)
(no sum over the index $`\alpha `$). From the second equation $`\stackrel{~}{C}_{2\alpha }=0`$ for all $`\alpha `$ such that $`\lambda _\alpha \mu _1`$. Using this result and the first equation, $`\stackrel{~}{C}_{2\alpha }=0`$ precisely for those $`\alpha `$ for which $`\lambda _\alpha =\mu _1`$. By assumption, no row of $`\stackrel{~}{C}`$ vanishes, so $`\stackrel{~}{C}_{2\alpha }`$ cannot all vanish. Therefore, there must be at least one $`\alpha `$ such that $`\lambda _\alpha =\mu _1`$, and the first equation must be absent, that is, the Jordan block has only size 1. Repeating this argument for all Jordan blocks of $`J`$ means that each eigenvalue of $`G`$ coincides with one of $`๐`$, and that $`J`$ is diagonal, that is, $`G`$ is diagonalisable. Writing $`G=S\mathrm{\Lambda }^{}S^1`$ where the diagonal matrix $`\mathrm{\Lambda }^{}`$ is a submatrix of $`\mathrm{\Lambda }`$, we have
$$(SC)๐=\mathrm{\Lambda }^{}(SC).$$
(88)
This means that the rows of $`SC`$ are left eigenvectors of $`๐`$, and parameterise to characteristic variables of the second-order system. We have shown
Theorem 5: The evolution of the original constraints is strongly hyperbolic if the second-order main system is, and its characteristic speeds are then a subset of those of the main system. Furthermore, we can find a basis of characteristic variables for the main system and the constraint system such that for each characteristic variable $`๐_\alpha `$ of the constraint system, there is a characteristic variable $`๐ฎ_\alpha `$ of the main system such that
$$๐_\alpha =_n๐ฎ_\alpha +\text{ transverse derivatives}.$$
(89)
Note that there is no such result for symmetric hyperbolicity.
## VII Symmetric hyperbolic-parabolic systems
Theorem 4.6.2 of GKO asserts the following: Assume we have a vector of variables $`u`$ and another vector of variables $`z`$, which obey a linear system of evolution equations of the form
$`_tu`$ $`=`$ $`D_{11}u+D_{12}z,`$ (90)
$`_tz`$ $`=`$ $`D_{21}u+D_{22}z.`$ (91)
Here the $`D`$ are linear spatial derivative operators whose coefficients can depend on $`t`$ and $`x^i`$. $`D_{11}`$ is a first-order derivative operator such that $`_tu=D_{11}u`$ is symmetric hyperbolic. $`D_{22}`$ is a second-order derivative operator such that $`_tz=D_{22}z`$ is parabolic. $`D_{12}`$ and $`D_{21}`$ are arbitrary first-order derivative operators. Then the coupled system is called mixed symmetric hyperbolic/parabolic. Its Cauchy problem with periodic boundaries is well-posed.
The theorem can be applied straightforwardly to second-order systems. We identify the variables $`u`$ of the theorem with the variables $`(v,w,d_i)`$ of the first-order reduction of what is to be the symmetric hyperbolic subsystem, and then go back to the second-order form of this subsystem by replacing $`d_i`$ with $`v_{,i}`$. The result is the system
$`\dot{v}`$ $`=`$ $`A_1^iv_{,i}+A_1v+A_2w+a+\underset{ยฏ}{Cz},`$ (92)
$`\dot{w}`$ $`=`$ $`B_1^{ij}v_{,ij}+B_1^iv_{,i}+B_1v+B_2^iw_{,i}+B_2w+b`$ (93)
$`+\underset{ยฏ}{D^iz_{,i}}+Dz,`$
$`\dot{z}`$ $`=`$ $`D_{22}z`$ (94)
$`+\underset{ยฏ}{E_1^{ij}v_{,ij}}+E_1^iv_{,i}+E_1v+\underset{ยฏ}{E_2^iw_{,i}}+E_2w.`$
The coupling operators $`D_{12}`$ and $`D_{21}`$ are parameterised by the matrices $`C`$ and $`D`$, and $`E`$, respectively. We have underlined their principal parts to show what order of derivative is allowed in the coupling terms.
Definition 6: A second-order system is called mixed symmetric hyperbolic-parabolic if it is of the form (92-94), such that $`D_{22}`$ is parabolic and the system (11-12) with the same coefficients is symmetric hyperbolic (in the sense of Definition 2a or 2b).
Theorem 4.6.2 of GKO then gives us
Lemma 6: The Cauchy problem with periodic boundary conditions for such a system is well-posed.
## VIII Conclusions
We have formalised the definition of strong or symmetric hyperbolicity of a system of evolution equations that are first order in time and second order in space by reducing them to an equivalent first-order system. We have given necessary and sufficient criteria for the existence of a reduction that is strongly hyperbolic or symmetric hyperbolic. These criteria are formulated entirely in terms of the principal part of the second-order system, without an explicit reference to the reduction.
We have proved that the definitions of strong hyperbolicity based on a first-order reduction, a pseudo-differential reduction, and a direct second-order approach are all equivalent. The definitions of symmetric hyperbolicity based on a first-order reduction and on a direct second-order approach are also equivalent.
In order to analyse the well-posedness of a given second-order system in practice, there are three non-trivial calculations to be carried out, independently of the approach in which one has defined hyperbolicity. Suppressing technical details, they are as follows.
#### Strong hyperbolicity
Strong hyperbolicity of the second-order system is equivalent to diagonalisability, with real eigenvalues, of the matrix $`๐`$.
#### Symmetric hyperbolicity
Symmetric hyperbolicity of the second-order system is equivalent to the existence of an energy and flux quadratic in the $`v_{,i}`$ and $`w`$.
#### Maximally dissipative boundary conditions
In order to impose maximally dissipative boundary conditions, one needs the full characteristic variables of a symmetric hyperbolic first-order reduction. This is done most easily starting from the left eigenvectors of $`๐`$ and the energy $`ฯต(w,v_{,i})`$.
We have established criteria for well-posedness of the second-order system regardless of any constraints it is subject to. However, if the second-order system is strongly hyperbolic, and there is a closed constraint system associated with it, then the constraint system is also strongly hyperbolic, and the characteristic variables of the constraint system are related to a subset of the characteristic variables of the main system.
It is known that a first-order symmetric hyperbolic system coupled to a parabolic system through at most first derivatives of all variables has a well-posed Cauchy problem. We have generalised this result to second-order symmetric hyperbolic systems through a reduction to first order.
Appendix D gives an example of a simple second order system discussed in both the first order reduction approach and the direct second-order approach.
###### Acknowledgements.
We would like to thank Fernando Barbero, Robert Beig, Gioel Calabrese, Ian Hinder, Gabriel Nagy and Olivier Sarbach for discussions and comments on the manuscript, and Louisiana State University, Caltech and the Erwin Schrรถdinger Institute for hospitality. JMM was supported by the Spanish MEC under the research projects BFM2002-04031-C02-02 and FIS2004-01912.
## Appendix A Diagonalisability of matrices with a zero block
Consider the matrix
$$M=\left(\begin{array}{cc}A& B\\ 0& C\end{array}\right)$$
(95)
where the square block $`A`$ has size $`n`$ and the square block $`C`$ has size $`m`$. We now prove that if $`M`$ is diagonalisable then $`A`$ and $`C`$ are both diagonalisable.
The eigenvalues of $`M`$ are clearly the union of those of $`A`$ and $`C`$. The eigenvectors of $`M`$ can be constructed from those of $`A`$ and $`C`$ as follows: Suppose we have
$`Av_i`$ $`=`$ $`\lambda _iv_i,i=1,\mathrm{},n,`$ (96)
$`Cw_j`$ $`=`$ $`\mu _jw_j,j=1,\mathrm{},m.`$ (97)
Then we can form the eigenvectors $`x_i=(v_i,0)`$, which span the invariant subspace of $`M`$, and $`y_j=(u_j,w_j)`$ such that
$`Mx_i`$ $`=`$ $`\lambda _ix_i,i=1,\mathrm{},n,`$ (98)
$`My_j`$ $`=`$ $`\mu _jy_j,j=1,\mathrm{},m,`$ (99)
with the condition $`(A\mu _j)u_j=Bw_j`$ for $`j=1,\mathrm{},m`$. If $`\mu _j`$ is not an eigenvalue of $`A`$ then $`(A\mu _j)`$ can be inverted and there is a unique solution for $`u_j`$. Therefore, if the eigenvalues of $`A`$ and $`C`$ are disjoint, the eigenvectors can be completed and the $`x_i`$ are linearly independent from the $`y_j`$ because they correspond to different eigenvalues. In this case $`M`$ is diagonalisable if and only if both $`A`$ and $`C`$ are diagonalisable.
Now suppose that $`A`$ and $`C`$ share an eigenvalue $`\lambda `$. The transformation
$$UMU^1=\left(\begin{array}{cc}U_1AU_1^1& D\\ 0& U_2CU_2^1\end{array}\right)$$
(100)
with
$$U=\left(\begin{array}{cc}U_1& U_1X\\ 0& U_2\end{array}\right)U^1=\left(\begin{array}{cc}U_1^1& XU_2^1\\ 0& U_2^1\end{array}\right)$$
(101)
and $`D=U_1(XCAX+B)U_2^1`$ brings $`A`$ and $`C`$ into Jordan form simultaneously for suitable $`U_1`$ and $`U_2`$. Without loss of generality we can assume that $`A`$ is a single Jordan block of eigenvalue $`\lambda `$ and rank$`(A\lambda I)=r`$, and that $`C`$ is another Jordan block with the same eigenvalue and rank$`(C\lambda I)=s`$. The matrix $`M`$ is diagonalisable if rank$`(M\lambda I)=0`$. We have rank$`(M\lambda I)r+s`$ because the $`r+s`$ columns of $`(M\lambda I)`$ containing a 1 in the second diagonal are linearly independent, while the matrix $`B`$ could provide additional linearly independent vectors. Therefore if $`M`$ is diagonalisable we must have $`r=s=0`$ for each Jordan block of $`A`$ and $`C`$, and so $`A`$ and $`C`$ are diagonalisable.
## Appendix B Hyperbolicity of the auxiliary constraint system
With the choice of reduction parameters that we have used in the proof of Theorem 1, namely (29), (30), and using the following further auxiliary constraints,
$`C_{ij}`$ $`=`$ $`c_{ij}+c_{[i,j]}=0,`$ (102)
$`C_{ijk}`$ $`=`$ $`c_{ij,k}+c_{jk,i}+c_{ki,j}=0,`$ (103)
the auxiliary constraint system can be reduced to the decoupled system
$`\dot{c}_i`$ $`=`$ $`A_1c_ii\mu ฯต_i{}_{}{}^{jk}c_{j,k}^{},`$ (104)
$`\dot{c}_{ij}`$ $`=`$ $`A_1c_{ij}i\mu ฯต_{[i}{}_{}{}^{kl}c_{kl,|j]}^{}.`$ (105)
This is strongly hyperbolic with speeds $`0,\pm \mu `$ for all $`n_i`$. Furthermore
$`\dot{C}_i`$ $`=`$ $`A_1C_ii\mu ฯต_i{}_{}{}^{jk}C_{j,k}^{},`$ (106)
$`\dot{C}_{ijk}`$ $`=`$ $`A_1C_{ijk},`$ (107)
which is also strongly hyperbolic (in fact, symmetric hyperbolic), and so in this case the auxiliary constraint system is well-posed.
## Appendix C Symmetric hyperbolicity and characteristic variables
Assume that $`P^i`$ admits a symmetriser $`H`$. As $`H`$ is Hermitian and positive definite, there is an invertible matrix S such that
$$H=S^{}S.$$
(108)
From this and
$$HP^i=P^iH$$
(109)
it follows that $`SP^nS^1`$ is Hermitian, for any direction $`n_i`$. Therefore it can be diagonalised by an orthogonal matrix $`R`$ (which generally depends on $`n_i`$), or
$$SP^nS^1=R^1\mathrm{\Lambda }R,$$
(110)
where $`\mathrm{\Lambda }`$ is diagonal. Therefore
$$P^n=T^1\mathrm{\Lambda }T,TRS,$$
(111)
and we have proved that symmetric hyperbolicity implies strong hyperbolicity.
Furthermore, as $`R`$ is orthogonal,
$$H=S^{}(R^{}R)S=T^{}T,$$
(112)
and so there is a preferred basis, namely the rows of $`T`$ (which generally depends on $`n_i`$), of left eigenvectors of $`P^n`$ in which $`H`$ is the unit matrix. In terms of the original basis
$$H=T^{}T,HP^n=T^{}\mathrm{\Lambda }T.$$
(113)
In quadratic forms the same fact can be written as
$$ฯต=\underset{\alpha }{}๐ฎ_\alpha ^2,\varphi ^n=\underset{\alpha }{}\lambda _\alpha ๐ฎ_\alpha ^2.$$
(114)
where the sum is over the characteristic variables in the basis encoded in the rows of $`T`$.
## Appendix D The KWB formulation of the Maxwell equations
We use this formulation of the Maxwell equations to illustrate some of the differences between the reduction approach and the second-order approach, namely the existence of โshortโ and โfullโ characteristic variables, the split of the free parameters of $`H`$ into those that are reduction parameters and those that are not, and the relation between the characteristic variables of the main and constraint systems. We work initially in the first-order reduction approach, and then re-interpret the same calculations in the language of the second-order approach afterwards.
The system has been discussed in KWB ; bssn1 . In flat spacetime, in radiation gauge and in the absence of charges (source terms), it is
$`\dot{A}_i`$ $`=`$ $`E_i,`$ (115)
$`\dot{E}_i`$ $`=`$ $`A_{i,j}^{}{}_{}{}^{,j}+(1a)A_{j,i}^{}{}_{}{}^{,j}+a\mathrm{\Gamma }_{,i},`$ (116)
$`\dot{\mathrm{\Gamma }}`$ $`=`$ $`(b1)E_{}^{i}{}_{,i}{}^{}`$ (117)
where $`v=\{A_i\}`$, $`w=\{E_i,\mathrm{\Gamma }\}`$ are the dynamical variables, and $`a`$ and $`b`$ are constants that parameterise addition of the (โoriginalโ) constraints
$`C_\mathrm{\Gamma }`$ $``$ $`\mathrm{\Gamma }A_{i}^{}{}_{}{}^{,i}=0,`$ (118)
$`C_E`$ $``$ $`E_{i}^{}{}_{}{}^{,i}=0`$ (119)
to the evolution equations. Repeated indices have been raised with the metric $`\delta ^{ij}`$ and are summed over. We shall also use $`\delta _{ij}`$ to decompose tensors into their normal and transverse parts. With $`d_{ij}A_{j,i}`$, we consider the first-order reduction.
$`\dot{A}_i`$ $`=`$ $`E_i,`$ (120)
$`\dot{E}_i`$ $`=`$ $`d_{ji}^{}{}_{}{}^{,j}+a\mathrm{\Gamma }_{,i}+(1a)d_{j}^{}{}_{}{}^{j}{}_{,i}{}^{}`$ (121)
$`+c\left(d_{ij}^{}{}_{}{}^{,j}d_{j}^{}{}_{}{}^{j}{}_{,i}{}^{}\right),`$
$`\dot{\mathrm{\Gamma }}`$ $`=`$ $`(b1)E_{}^{i}{}_{,i}{}^{},`$ (122)
$`\dot{d}_{ij}`$ $`=`$ $`E_{j,i}.`$ (123)
The constant $`c`$ parameterises $`B_3^{ij}`$, and $`D_{i}^{}{}_{}{}^{jk}`$ has been set to zero. For the other reduction parameters we have made the standard choice (29), so that the $`A_i`$ decouple from the $`E_i`$ and $`\mathrm{\Gamma }`$.
The matrix $`๐ซ`$, leaving out the zero rows and columns corresponding to $`A_i`$, is block-diagonal with the blocks
$`\left(\begin{array}{cccc}0& a& a& 1ac\\ b1& 0& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right)`$ $`\left(\begin{array}{c}E_n\\ \mathrm{\Gamma }\\ d_{nn}\\ d_{qq}\end{array}\right),`$ (132)
$`\left(\begin{array}{ccc}0& 1& c\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)`$ $`\left(\begin{array}{c}E_A\\ d_{nA}\\ d_{An}\end{array}\right),`$ (139)
$`\left(0\right)`$ $`\left(\widehat{d}_{AB}\right).`$ (140)
Here $`d_{qq}q^{ij}d_{ij}`$, and $`\widehat{d}_{AB}`$ represents the transverse trace-free object $`q_i^kq_j^ld_{kl}(1/2)q_{ij}d_{qq}`$. The characteristic variables are
$`๐ฎ_0`$ $``$ $`\mathrm{\Gamma }+(b1)d_{nn},`$ (141)
$`๐ฎ_\pm `$ $``$ $`a(\mathrm{\Gamma }d_{nn})+(1ac)d_{qq}\pm \sqrt{ab}E_n,`$ (142)
$`๐ฎ_{\pm A}`$ $``$ $`d_{nA}cd_{An}E_A,`$ (143)
with speeds $`(0,\pm \sqrt{ab},\pm 1)`$, and $`d_{qq}`$, $`d_{An}`$ and $`\widehat{d}_{AB}`$ with zero speed.
The first-order reduction admits the conserved energy
$`ฯต`$ $`=`$ $`E_iE^i+d_{ij}d^{ij}2a\mathrm{\Gamma }d_{i}^{}{}_{}{}^{i}+(2a1ab)(d_{i}^{}{}_{}{}^{i})^2`$ (144)
$`+c_1\left[\mathrm{\Gamma }+(b1)d_{i}^{}{}_{}{}^{i}\right]^2+c[(d_{i}^{}{}_{}{}^{i})^2d_{ij}d^{ji}]`$
with the flux
$`\varphi ^i`$ $`=`$ $`2\left[a(\mathrm{\Gamma }d_{j}^{}{}_{}{}^{j})E^i+d_{j}^{}{}_{}{}^{j}E^id^{ij}E_j\right]`$ (145)
$`+2c(d^{ji}E_jd_{j}^{}{}_{}{}^{j}E^i),`$
where $`c_1`$ is a free parameter in the energy, and $`c`$ is the reduction parameter introduced above.
We now review how one would deal with the same system in a direct second-order approach \[pointing out the relation to the first-order approach in square brackets\]. In order to show strong hyperbolicity, one would diagonalise the matrix $`๐`$. It is block-diagonal with the blocks
$$\left(\begin{array}{ccc}0& a& a\\ b1& 0& 0\\ 1& 0& 0\end{array}\right)\left(\begin{array}{c}E_n\\ \mathrm{\Gamma }\\ A_{n,n}\end{array}\right)$$
(146)
and
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{c}E_A\\ A_{A,n}\end{array}\right).$$
(147)
\[$`๐`$ is the sub-matrix of $`๐ซ`$ obtained by suppressing the rows and columns relating to $`d_{Ai}`$. In a different point of view, we could set the rows and columns relating to $`d_{Ai}`$ in $`๐ซ`$ to zero by allowing the reduction parameters to depend explicitly on $`n_i`$. In our example, this corresponds to setting $`c=1a`$ in (132) but $`c=0`$ in (139).\]
The โshortโ characteristic variables of the second-order system are obtained as eigenvectors of $`๐^{}`$. They are
$`๐ฎ_0^{}`$ $``$ $`\mathrm{\Gamma }+(b1)A_{n,n},`$ (148)
$`๐ฎ_\pm ^{}`$ $``$ $`a(\mathrm{\Gamma }A_{n,n})\pm \sqrt{ab}E_n,`$ (149)
$`๐ฎ_{\pm A}^{}`$ $``$ $`A_{A,n}E_A,`$ (150)
with speeds $`(0,\pm \sqrt{ab},\pm 1)`$. \[These are the characteristic variables of the reduction, minus all terms in $`d_{Ai}=A_{i,A}`$.\]
To show symmetric hyperbolicity of the second-order system, we find an an energy. The most general one is
$`ฯต`$ $`=`$ $`E_iE^i+A_{i,j}A^{i,j}2a\mathrm{\Gamma }A_{i}^{}{}_{}{}^{,i}+(2a1ab)(A_{i}^{}{}_{}{}^{,i})^2`$
$`+c_1\left[\mathrm{\Gamma }+(b1)A_{i}^{}{}_{}{}^{,i}\right]^2+c_2[(A_{i}^{}{}_{}{}^{,i})^2A_{i,j}A^{j,i}]`$
with the flux
$`\varphi ^i`$ $`=`$ $`2\left[a(\mathrm{\Gamma }A_{j}^{}{}_{}{}^{,j})E^i+A_{j}^{}{}_{}{}^{,j}E^iA^{j,i}E_j\right]`$ (152)
$`+2c_2(A^{i,j}E_jA_{j}^{}{}_{}{}^{,j}E^i),`$
where $`c_1`$ and $`c_2`$ are free parameters. \[This is identical to the energy of the reduction, except that $`c_2`$ is now not a reduction parameter. Rather, the term it multiplies is independently conserved if we allow commutation of partial derivatives.\]
In order to impose maximally dissipative boundary conditions, one needs the full characteristic variables $`๐ฎ_\alpha `$. We find this by expressing $`ฯต`$ as a quadratic form in $`๐ฎ_\alpha `$. Comparing $`ฯต`$ with the $`๐ฎ_\alpha ^{}`$ suggests that
$`ฯต`$ $`=`$ $`{\displaystyle \frac{1}{2ab}}\left(๐ฎ_+^2+๐ฎ_{}^2\right)+{\displaystyle \frac{1}{2}}\left(๐ฎ_{+A}๐ฎ^{+A}+๐ฎ_A๐ฎ^A\right)`$ (153)
$`+\left(c_1{\displaystyle \frac{a}{b}}\right)๐ฎ_0^2c_2(๐ฎ_{+A}+๐ฎ_A)A^{n,A}`$
$`+\text{ quadratic in }A_{i,A},`$
with $`๐ฎ_\alpha =๐ฎ_\alpha ^{}`$ \+ multiples of the $`A_{i,A}`$. \[The result is equivalent to the $`๐ฎ_\alpha `$ of the reduction, with $`d_{ij}A_{j,i}`$ and $`cc_2`$.\]
Finally, the constraint evolution system is
$`\dot{C}_\mathrm{\Gamma }`$ $`=`$ $`bC_E,`$ (154)
$`\dot{C}_E`$ $`=`$ $`aC_{\mathrm{\Gamma },i}^{}{}_{}{}^{,i}.`$ (155)
Its non-trivial characteristic variables are
$$๐_\pm =a_nC_\mathrm{\Gamma }\pm \sqrt{ab}C_E=_n๐ฎ_\pm +\text{transv. deriv.},$$
(156)
and these expressions hold for both the second-order system and the reduction. |
warning/0506/math0506239.html | ar5iv | text | # Reconstruction and subgaussian operators
## 0 Introduction
The aim of this article is to investigate the linear โapproximate reconstructionโ problem in $`^n`$. In such a problem, one is given a set $`T^n`$ and the goal is to be able to approximate any unknown $`vT`$ using random linear measurements. In other words, one is given the set of values $`(X_i,v)_{i=1}^k`$, where $`X_1,\mathrm{},X_k`$ are independent random vectors in $`^n`$ selected according to some probability measure $`\mu `$. Using this information (and the fact that the unknown vector $`v`$ belongs to $`T`$) one has to produce, with very high probability with respect to $`\mu ^k`$, some $`tT`$, such that the Euclidean norm $`|tv|\epsilon (k)`$ for $`\epsilon (k)`$ as small as possible. Of course, the random sampling method has to be โuniversalโ in some sense and not tailored to a specific set $`T`$; and it is natural to expect that the degree of approximation $`\epsilon (k)`$ depends on some geometric parameter associated with $`T`$.
Questions of a similar flavor have been thoroughly studied in nonparametric statistics and statistical learning theory (see, for example, \[BBL\] and \[M, MT\] and references therein). This particular problem has been addressed by several authors (see \[CT1, CT2, CT3, RV\] for the most recent contributions), in a rather restricted context. First of all, the sets considered were either the unit ball in $`\mathrm{}_1^n`$ or the unit balls in weak $`\mathrm{}_p^n`$ spaces for $`0<p<1`$ \- and the proofs of the approximation estimates depended on the choice of those particular sets. Second, the sampling process was done when $`X_i`$ were distributed according to the Gaussian measure on $`^n`$ or in \[CT1\] for Fourier ensemble.
In contrast, we present a method which is very general. Our results hold for any set $`T^n`$, and the class of measures that could be used is broad; it contains all probability measures on $`^n`$ which are isotropic and subgaussian, that is, satisfy that for every $`y^n`$, $`๐ผ|X,y|^2=|y|^2`$, and the random variables $`X,y`$ are subgaussian with constant $`\alpha |y|`$ for some $`\alpha 1`$. (see Definition 2.2, below). This class of measures contains, among others, the Gaussian measure on $`^n`$, the uniform measure on the vertices of the combinatorial cube and the normalized volume measure on various convex, symmetric bodies (e.g. the unit balls of $`\mathrm{}_p^n`$ for $`2p\mathrm{}`$).
It turns out that the key parameter in the estimate on the degree of approximation $`\epsilon (k)`$ is indeed geometric in nature. Moreover, the analysis of the approximation problem is centered around the way the random operator $`\mathrm{\Gamma }=_{i=1}^kX_i,e_i`$ (where $`(e_i)_{i=1}^k`$ is the standard basis in $`^k`$) acts on subsets of the unit sphere.
Our geometric approach, when applied to the sets $`T`$ considered in \[CT1\], yields the optimal estimates on $`\epsilon (k)`$, and with better probability estimates (of the order of $`1\mathrm{exp}(ck)`$), even when the sampling is done according to an arbitrary isotropic, subgaussian measure. Moreover, our result is more robust that the one from \[CT1\] in the following sense. The reconstruction method suggested in \[CT1\] is to find some $`yT`$ such that for every $`1ik`$, $`X_i,v=X_i,y`$, and then to show that such $`v`$ and $`y`$ must necessarily be close. From our theorem it is clear that one can choose any $`yT`$ for which $`_{i=1}^kX_i,yv^2`$ is relatively small, which is far more stable algorithmically.
For the moment, let us present a simple version of the main result we prove in this direction, and to that end we require the following notation. Let $`(g_i)_{i=1}^n`$ be independent standard Gaussian random variables. Let $`T^n`$ be a star-shaped set (i.e. for every $`tT`$ and $`0\lambda 1`$, $`\lambda tT`$) and consider the following geometric parameter
$$\mathrm{}_{}(T)=๐ผ\underset{tT}{sup}\left|\underset{i=1}{\overset{n}{}}g_it_i\right|$$
which is, up to a factor of the order of $`2\sqrt{n}`$, the mean width of the body $`T`$. We now define a more sensitive parameter, which as we will see in this article, is the right parameter to control the error term $`\epsilon (k)`$ for a star-shaped set $`T`$:
$$r_k^{}(\theta ,T):=inf\{\rho >0:\rho c\alpha ^2\mathrm{}_{}(T\rho S^{n1})/\theta \sqrt{k}\}.$$
We may now state a version of our main result concerning approximate reconstruction.
Theorem A There exist an absolute constant $`c_1`$ for which the following holds. Let $`T`$ be a star-shaped subset of $`^n`$. Let $`\mu `$ be an isotropic, subgaussian measure with constant $`\alpha `$ on $`^n`$ and set $`X_1,\mathrm{},X_k`$ be independent, selected according to $`\mu `$. Then, with probability at least $`1\mathrm{exp}(c_1k/\alpha ^4)`$, every $`y,vT`$ satisfy that
$$|yv|2\left(\frac{1}{k}\underset{i=1}{\overset{k}{}}\left(X_i,vX_i,y\right)^2\right)^{1/2}+r_k^{}(1/2,TT).$$
The parameter $`r_k^{}(\theta ,T)`$ can be estimated for unit ball of classical normed or quasi-normed spaces (see Section 3). In particular, if $`T=B_1^n`$ then with the same hypothesis and probability as above, one has
$$|yv|2\left(\frac{1}{k}\underset{i=1}{\overset{k}{}}\left(X_i,vX_i,y\right)^2\right)^{1/2}+c\alpha ^2\sqrt{\frac{\mathrm{log}(c\alpha ^4n/k)}{k}}$$
where $`c>0`$ is an absolute constant; this leads to the optimal estimate for $`\epsilon (k)`$ for that set.
The main idea in the proof of this theorem is the fact that excluding a set with exponentially small probability, the random operator $`\frac{1}{\sqrt{k}}_{i=1}^kX_i,`$ is a very good isomorphism on elements of $`T`$ whose Euclidean norm is large enough (see Section 2 for more details).
A question of a similar nature we investigate here focuses on โexact reconstructionโ of vectors in $`^n`$ that have a short support. Suppose that $`z^n`$ is in the unit Euclidean ball, and has a relatively short support $`m<<n`$. The aim is to use a random sampling procedure to identify $`z`$ exactly, rather than just to approximate it. The motivation for this problem comes from error correcting codes, in which one has to overcome random noise that corrupts a part of a transmitted signal. The noise is modelled by adding to the transmitted vector the noise vector $`z`$. The assumption that the noise does not change many bits in the original signal implies that $`z`$ has a short support, and thus, in order to correct the code, one has to identify the noise vector $`z`$ exactly. Since error correcting codes are not the main focus of this article, we will not explore this topic further, but rather refer the interested reader to \[MS, CT2, RV\] and references therein for more information.
In the geometric context we are interested in, the problem has been studied in \[CT2, RV\], where it was shown that if $`z`$ has a short support relative to the dimension $`n`$ and the size of the sample $`k`$, and if $`\mathrm{\Gamma }`$ is a $`k\times n`$ matrix whose entries are independent, standard Gaussian variables, then with probability at least $`1\mathrm{exp}(ck)`$, the minimization problem
$$(P)\mathrm{min}v_{\mathrm{}_1^n}\mathrm{for}\mathrm{\Gamma }v=\mathrm{\Gamma }z,$$
has a unique solution, which is $`z`$. Thus, solution to this minimization problem will pin-point the โnoise vectorโ $`z`$. The idea of using this minimization problem was first suggested in \[CDS\].
We extend this result to a general random matrix whose rows are $`X_1,\mathrm{},X_k`$, sampled independently according to an isotropic, subgaussian measure.
Theorem B Let $`\mathrm{\Gamma }`$ be as above. With probability at least $`1\mathrm{exp}(c_1k/\alpha ^4)`$, any vector $`z`$ whose support has cardinality less than $`c_2k/\mathrm{log}(c_3n/k)`$ is the unique minimizer of the problem $`(P)`$, where $`c_1,c_2,c_3`$ are absolute constants.
Interestingly enough, the same analysis yields some information on the geometry of $`\{1,1\}`$ random polytopes. Indeed, consider the $`k\times n`$ matrix $`\mathrm{\Gamma }`$ whose entries are independent symmetric $`\{1,1\}`$ valued random variables. Thus, $`\mathrm{\Gamma }`$ is a random operator selected according to the uniform measure on the combinatorial cube $`\{1,1\}^n`$, which is an isotropic, subgaussian measure with constant $`\alpha =1`$. The columns of $`\mathrm{\Gamma }`$, denoted by $`v_1,\mathrm{},v_n`$ are vectors in $`\{1,1\}^k`$ and let $`K^+=\mathrm{conv}(v_1,\mathrm{},v_n)`$ be the convex polytope generated by $`\mathrm{\Gamma }`$; $`K^+`$ is thus called a random $`\{1,1\}`$-polytope.
A convex polytope is called $`m`$-neighborly if any set of less than $`m`$ of its vertices is the vertex set of a face (see \[Z\]). The following result yields the surprising fact that a random $`\{1,1\}`$-polytope is $`m`$-neighborly for a relatively large $`m`$. In particular, it will have the maximal number of $`r`$-faces for $`rm`$.
Theorem C There exist absolute constants $`c_1,c_2,c_3`$ for which the following holds. For $`1kn`$, with probability larger than $`1\mathrm{exp}(c_1k)`$, a $`k`$-dimensional random $`\{1,1\}`$ convex polytope with $`n`$ vertices is $`m`$-neighborly provided that
$$m\frac{c_2k}{\mathrm{log}(c_3n/k)}$$
The main technical tool we use throughout this article, which is of independent interest, is an estimate of the behavior of the supremum of the empirical process $`fZ_f=k^1_{i=1}^kf^2(X_i)1`$ indexed by a subset of the $`L_2`$ sphere. That is,
$$\underset{fF}{sup}\left|\frac{1}{k}\underset{i=1}{\overset{k}{}}f^2(X_i)1\right|,$$
where $`X_1,\mathrm{},X_k`$ are independent, distributed according to the a probability measure $`\mu `$, and under the assumption that every $`fF`$ satisfies that $`๐ผf^2=1`$. We assume further that $`F`$ is a bounded set with respect to the $`\psi _2`$ norm, defined for a random variable $`Y`$ by
$$Y_{\psi _2}=inf\{u>0:๐ผ\mathrm{exp}(Y^2/u^2)2\}.$$
To formulate the following result, we require the notion of the $`\gamma _p`$ functional \[Ta2\]. For a metric space $`(T,d)`$, an admissible sequence of $`T`$ is a collection of subsets of $`T`$, $`\{T_s:s0\}`$, such that for every $`s1`$, $`|T_s|=2^{2^s}`$ and $`|T_0|=1`$. for $`p=1,2`$, we define the $`\gamma _p`$ functional by
$$\gamma _p(T,d)=inf\underset{tT}{sup}\underset{s=0}{\overset{\mathrm{}}{}}2^{s/p}d(t,T_s),$$
where the infimum is taken with respect to all admissible sequences of $`T`$.
Theorem D There exist absolute constants $`c_1,c_2,c_3`$ and for which the following holds. Let $`(\mathrm{\Omega },\mu )`$ be a probability space, set $`F`$ be a subset of the unit sphere of $`L_2(\mu )`$ and assume that $`\mathrm{diam}(F,_{\psi _2})=\alpha `$. Then, for any $`\theta >0`$ and $`k1`$ satisfying
$$c_1\alpha \gamma _2(F,_{\psi _2})\theta \sqrt{k},$$
with probability at least $`1\mathrm{exp}(c_2\theta ^2k/\alpha ^4)`$,
$$\underset{fF}{sup}|Z_f|\theta .$$
Moreover, if $`F`$ is symmetric, then
$$๐ผ\underset{fF}{sup}|Z_f|c_3\alpha ^2\frac{\gamma _2(F,_{\psi _2})}{\sqrt{k}}.$$
Theorem D improves a result of a similar flavor from \[KM\] in two ways. First of all, the bound on the probability is exponential in the sample size $`k`$ which was not the case in \[KM\]. Second, we were able to bound the expectation of the supremum of the empirical process using only a $`\gamma _2`$ term. This fact is surprising because the expectation of the supremum of empirical processes is usually controlled by two terms; the first one bounds the subgaussian part of the process and is captured by the $`\gamma _2`$ functional with respect to the underlying metric. The other term is needed to control sub-exponential part of the empirical process, and is bounded by the $`\gamma _1`$ functional with respect to an appropriate metric (see \[Ta2\] for more information on the connections between the $`\gamma _p`$ functionals and empirical processes). Theorem D shows that the expectation of the supremum of $`|Z_f|`$ behaves as if $`\{Z_f:fF\}`$ were a subgaussian process with respect to the $`\psi _2`$ metric (although it is not), and this is due to the key fact that all the functions in $`F`$ have the same second moment.
We end this introduction with the organization of the article. In Section 1 we present the proof of Theorem D and some of its corollaries we require. In Section 2 we illustrate these results in the case of linear processes which corresponds to linear measurements. In Section 3 we investigate the approximate reconstruction problem for a general set, and in Section 4 we present a proof for the exact reconstruction scheme, with its application to the geometry of random $`\{1,1\}`$-polytopes.
Throughout this article we will use letters such as $`c,c_1,..`$ to denote absolute constants which may change depending on the context. We denote by $`(e_i)_{i=1}^n`$ the canonical basis of $`^n`$, by $`|x|`$ the Euclidean norm of a vector $`x`$ and by $`B_2^n`$ the Euclidean unit ball. We also denote by $`|I|`$ the cardinality of a finite set $`I`$.
Acknowledgement: A part of this work was done when the second and the third authors were visiting the Australian National University, Canberra; and when the first and the third authors were visiting Universitรฉ de Marne-la-Vallรฉe. They wish to thank these institutions for their hospitality.
## 1 Empirical Processes
In this section we present some results in empirical processes that will be central to our analysis. All the results focus on understanding the process $`Z_f=\frac{1}{k}_{i=1}^kf^2(X_i)๐ผf^2`$, where $`k1`$ and $`X_1,\mathrm{},X_k`$ are independent random variables distributed according to a probability measure $`\mu `$. In particular, we investigate the behavior of $`sup_{fF}|Z_f|`$ in terms of various metric structures on $`F`$, and under the key assumption that every $`fF`$ has the same second moment. The parameters involved are standard in generic chaining type estimates (see \[Ta2\] for a comprehensive study of this topic).
Recall that the $`\psi _p`$ norm of a random variable $`X`$ is defined as
$$X_{\psi _p}=inf\{u>0:๐ผ\mathrm{exp}\left(|X|^p/u^p\right)2\}.$$
It is standard to verify (see for example \[VW\]) that if $`X`$ has a bounded $`\psi _2`$ norm, then it is subgaussian with parameter $`cX_{\psi _2}`$ for some absolute constant $`c`$. More generally, a bounded $`\psi _p`$ norm implies that $`X`$ has a tail behavior, $`(|X|>u)`$, of the type $`\mathrm{exp}(cu^p/X_{\psi _p})`$.
Our first fundamental ingredient is the well known Bernsteinโs inequality which we shall use in the form of a $`\psi _1`$ estimates (\[VW\]).
###### Lemma 1.1
Let $`Y_1,\mathrm{},Y_k`$ be independent random variables with zero mean such that for some $`b>0`$ and every $`i`$, $`Y_i_{\psi _1}b`$. Then, for any $`u>0`$,
$$\left(\left|\frac{1}{k}\underset{i=i}{\overset{k}{}}Y_i\right|>u\right)2\mathrm{exp}\left(ck\mathrm{min}(\frac{u}{b},\frac{u^2}{b^2})\right),$$
(1.1)
where $`c>0`$ is an absolute constant.
We will be interested in classes of functions $`FL_2(\mu )`$ bounded in the $`\psi _2`$-norm; we assume without loss of generality that $`F`$ is symmetric and we let $`diam(F,_{\psi _2}):=2sup_{fF}f_{\psi _2}`$. Additionally, in many technical arguments we shall often consider classes $`FS_{L_2}`$, where $`S_{L_2}=\{f:f_{L_2}=1\}`$ is the unit sphere in $`L_2(\mu )`$.
Let $`X_1,X_2,\mathrm{}`$ be independent random variables distributed according to $`\mu `$. Fix $`k1`$ and for $`fF`$ define the random variables $`Z_f`$ and $`W_f`$ by
$$Z_f=\frac{1}{k}\underset{i=1}{\overset{k}{}}f^2(X_i)๐ผf^2\text{and}W_f=\left(\frac{1}{k}\underset{i=1}{\overset{k}{}}f^2(X_i)\right)^{1/2}.$$
The first lemma follows easily from Bernsteinโs inequality. We state it in the form convenient for further use, and give a proof of one part, for completeness.
###### Lemma 1.2
There exists an absolute constant $`c_1>0`$ for which the following holds. Let $`FS_{L_2}`$, $`\alpha =diam(F,_{\psi _2})`$ and set $`k1`$. For every $`f,gF`$ and every $`u2`$ we have
$$\left(W_{fg}ufg_{\psi _2}\right)2\mathrm{exp}(c_1ku^2).$$
Also, for every $`u>0`$,
$$\left(\left|Z_fZ_g\right|u\alpha fg_{\psi _2}\right)2\mathrm{exp}\left(c_1k\mathrm{min}(u,u^2)\right),$$
and
$$\left(\left|Z_f\right|u\alpha ^2\right)2\mathrm{exp}\left(c_1k\mathrm{min}(u,u^2)\right).$$
Proof. We show the standard proof of the first estimate. Other estimates are proved similarly (see e.g., \[KM\], Lemma 3.2).
Clearly,
$$๐ผW_{fg}^2=\frac{1}{k}๐ผ\underset{i=1}{\overset{k}{}}(fg)^2(X_i)=๐ผ(fg)^2(X_1)=fg_{L_2}^2.$$
Applying Bernsteinโs inequality it follows that for $`t>0`$,
$`\left(\left|W_{fg}^2fg_{L_2}^2\right|t\right)`$
$``$ $`2\mathrm{exp}\left(ck\mathrm{min}({\displaystyle \frac{t}{(fg)^2_{\psi _1}}},\left({\displaystyle \frac{t}{(fg)^2_{\psi _1}}}\right)^2)\right).`$
Since $`h^2_{\psi _1}=h_{\psi _2}^2`$ for every function $`h`$, then letting $`t=(u^21)fg_{\psi _2}^2`$,
$`(W_{fg}^2`$ $``$ $`u^2fg_{\psi _2}^2)`$
$``$ $`\left(W_{fg}^2fg_{L_2}^2(u^21)fg_{\psi _2}^2\right)`$
$``$ $`2\mathrm{exp}\left(ck\mathrm{min}(u^2/2,u^4/4)\right),`$
as promised.
Now we return to one of the basic notions used in this paper, that of the $`\gamma _2`$-functional. Let $`(T,d)`$ be a metric space. Recall that an admissible sequence of $`T`$ is a collection of subsets of $`T`$, $`\{T_s:s0\}`$, such that for every $`s1`$, $`|T_s|=2^{2^s}`$ and $`|T_0|=1`$.
###### Definition 1.3
For a metric space $`(T,d)`$ and $`p=1,2`$, define
$$\gamma _p(T,d)=inf\underset{tT}{sup}\underset{s=0}{\overset{\mathrm{}}{}}2^{s/p}d(t,T_s),$$
where the infimum is taken with respect to all admissible sequences of $`T`$. In cases where the metric is clear from the context, we will denote the $`\gamma _p`$ functional by $`\gamma _p(T)`$.
Set $`\pi _s:TT_s`$ to be a metric projection function onto $`T_s`$, that is, for every $`tT`$, $`\pi _s(t)`$ is a nearest element to $`t`$ in $`T_s`$ with respect to the metric $`d`$. It is easy to verify by the triangle inequality that for every admissible sequence and every $`tT`$, $`_{s=0}^{\mathrm{}}2^{s/2}d(\pi _{s+1}(t),\pi _s(t))(1+1/\sqrt{2})_{s=0}^{\mathrm{}}2^{s/2}d(t,T_s)`$ and that $`diam(T,d)2\gamma _2(T,d)`$.
We say that a set $`F`$ is star-shaped if the fact that $`fF`$ implies that $`\lambda fF`$ for every $`0\lambda 1`$.
The next Theorem shows that excluding a set with exponentially small probability, $`W_f`$ is close to being an isometry in the $`L_2(\mu )`$ sense for functions in $`F`$ that have a relatively large norm.
###### Theorem 1.4
There exist absolute constants $`c,\overline{c}>0`$ for which the following holds. Let $`FL_2(\mu )`$ be star-shaped, $`\alpha =diam(F,_{\psi _2})`$ and $`k1`$. For any $`0<\theta <1`$, with probability at least $`1\mathrm{exp}(\overline{c}\theta ^2k/\alpha ^4)`$, then for all $`fF`$ satisfying $`๐ผf^2r_k^{}(\theta )^2`$, we have
$$(1\theta )๐ผf^2\frac{W_f^2}{k}(1+\theta )๐ผf^2,$$
(1.2)
where
$$r_k^{}(\theta )=r_k^{}(\theta ,F):=inf\{\rho >0:\rho c\alpha \frac{\gamma _2(F\rho S_{L_2},_{\psi _2})}{\theta \sqrt{k}}\}.$$
(1.3)
The two-sided inequality (1.2) is intimately related to an estimate on $`sup_{fF}|Z_f|`$, which, in turn, is based on two ingredients. The first one shows, in the language of the standard chaining approach, that one can control the โend partsโ of all chains. its proof is essentially the same as Lemma 2.3 from \[KM\].
###### Lemma 1.5
There exists an absolute constant $`C`$ for which the following holds. Let $`FS_{L_2}`$, $`\alpha =diam(F,_{\psi _2})`$ and $`k1`$. There is $`F^{}F`$ such that $`|F^{}|4^k`$ and with probability at least $`1\mathrm{exp}(k)`$, we have, for every $`fF`$,
$$W_{f\pi _F^{}(f)}C\gamma _2(F,_{\psi _2})/\sqrt{k},$$
(1.4)
where $`\pi _F^{}(f)`$ is a nearest point to $`f`$ in $`F^{}`$ with respect to the $`\psi _2`$ metric.
Proof. Let $`\{F_s:s0\}`$ be an โalmost optimalโ admissible sequence of $`F`$. Then for every $`fF`$,
$$\underset{s=0}{\overset{\mathrm{}}{}}2^{s/2}\pi _{s+1}(f)\pi _s(f)_{\psi _2}2\gamma _2(F,_{\psi _2}).$$
Let $`s_0`$ be the minimal integer such that $`2^{s_0}>k`$, and let $`F^{}=F_{s_0}`$. Then $`|F^{}|2^{2k}=4^k`$. Write
$$f\pi _{s_0}(f)=\underset{s=s_0}{\overset{\mathrm{}}{}}\left(\pi _{s+1}(f)\pi _s(f)\right).$$
Since $`W`$ is sub-additive then
$$W_{f\pi _{s_0}}(f)\underset{s=s_0}{\overset{\mathrm{}}{}}W_{\pi _{s+1}(f)\pi _s(f)}.$$
For any $`fF`$, $`ss_0`$ and $`\xi 2`$, noting that $`2^s>k`$, it follows by Lemma 1.2 that
$$\left(W_{\pi _{s+1}(f)\pi _s(f)}\xi \frac{2^{s/2}}{\sqrt{k}}\pi _{s+1}(f)\pi _s(f)_{\psi _2}\right)2\mathrm{exp}(c_1\xi ^22^s).$$
(1.5)
Since $`|F_s|2^{2^s}`$, there are at most $`2^{2^{s+2}}`$ pairs of $`\pi _{s+1}(f)`$ and $`\pi _s(f)`$. Thus, for every $`ss_0`$, the probability of the event from (1.5) holding for some $`fF`$ is less than or equal to $`2^{2^{s+2}}2\mathrm{exp}(c_1\xi ^22^s)\mathrm{exp}(2^{s+3}c\xi ^22^s)`$, which, for $`\xi \xi _0:=\mathrm{max}(4/\sqrt{c_1},2)`$, does not exceed $`\mathrm{exp}(c_1\xi ^22^{s1})`$.
Combining these estimates together it follows that
$$W_{f\pi _{s_0}(f)}\xi \underset{s=s_0}{\overset{\mathrm{}}{}}\frac{2^{s/2}}{\sqrt{k}}\pi _{s+1}(f)\pi _s(f)_{\psi _2}2\xi \frac{\gamma _2(F,_{\psi _2})}{\sqrt{k}},$$
outside a set of probability
$$\underset{s=s_0}{\overset{\mathrm{}}{}}\mathrm{exp}(c_1\xi ^22^{s1})\mathrm{exp}(c_1\xi ^22^{s_0}/4).$$
We complete the proof setting, for example $`\xi =\mathrm{max}(\xi _0,2/\sqrt{c_1})`$ and recalling that $`2^{s_0}>k`$.
###### Remark 1.6
The proof of the lemma shows that there exist absolute constants $`c^{},c^{\prime \prime }>0`$ such that for every $`\xi c^{}`$,
$$\left(\underset{fF}{sup}W_{f\pi (f)}\xi \frac{\gamma _2(F,_{\psi _2})}{\sqrt{k}}\right)\mathrm{exp}(c^{\prime \prime }\xi ^2k),$$
a fact which will be used later.
The next lemma estimates the supremum $`sup_{fF^{}}|Z_f|`$, where the supremum is taken over a subset $`F^{}`$ of $`F`$ of a relatively small cardinality, or in other words, over the โbeginning partโ of a chain. However, in order to get an exponential in $`k`$ estimates on probability, we require a separate argument (generic chaining) for the โmiddle partโ of a chain while for the โvery beginningโ it is sufficient to use a standard concentration estimate.
###### Lemma 1.7
There exist absolute constants $`C`$ and $`c^{\prime \prime \prime }>0`$ for which the following holds. Let $`FS_{L_2}`$ and $`\alpha =diam(F,_{\psi _2})`$. Let $`k1`$ and $`F^{}F`$ such that $`|F^{}|4^k`$. Then for every $`w>0`$,
$$\underset{fF^{}}{sup}|Z_f|C\alpha \frac{\gamma _2(F,_{\psi _2})}{\sqrt{k}}+\alpha ^2w,$$
with probability larger than or equal to $`13\mathrm{exp}(c^{\prime \prime \prime }k\mathrm{min}(w,w^2))`$.
Proof. Let $`(F_s)_{s=0}^{\mathrm{}}`$ be an almost optimal admissible sequence of $`F^{}`$, set $`s_0`$ to be the minimal integer such that $`2^{s_0}>2k`$ and fix $`s_1s_0`$ to be determined later. Since $`|F^{}|4^k`$, it follows that $`F_s=F^{}`$ for every $`ss_0`$, and that
$$Z_fZ_{\pi _{s_1}(f)}=\underset{s=s_1+1}{\overset{s_0}{}}\left(Z_{\pi _s(f)}Z_{\pi _{s1}(f)}\right).$$
By Lemma 1.2, for every $`fF^{}`$, $`1ss_0`$ and $`u>0`$,
$`\left(\right|Z_{\pi _s(f)}Z_{\pi _{s1}(f)}|`$ $`u\alpha \sqrt{{\displaystyle \frac{2^s}{k}}}\pi _{s+1}(x)\pi _s(x)_{\psi _2})`$
$`2\mathrm{exp}(c_1\mathrm{min}((u\sqrt{2^s/k}),(u\sqrt{2^s/k})^2))`$
$`2\mathrm{exp}(c_1\mathrm{min}(u,u^2)2^{s2}).`$
(For the latter inequality observe that if $`ss_0`$ then $`2^s/k4`$, and thus $`\mathrm{min}((u\sqrt{2^s/k}),(u\sqrt{2^s/k})^2)\mathrm{min}(u,u^2)\mathrm{\hspace{0.17em}2}^s/(4k)`$.)
Taking $`u`$ large enough (for example, $`u=2^5/c_1`$ will suffice) we may ensure that
$$\underset{s=s_1+1}{\overset{s_0}{}}2^{2^{s+2}}\mathrm{exp}(c_1u2^{s2})\underset{s=s_1+1}{\overset{s_0}{}}\mathrm{exp}(2^{s+3})\mathrm{exp}(2^{s_1}).$$
Therefore, since there are at most $`2^{2^{s+2}}`$ possible pairs of $`\pi _{s+1}(f)`$ and $`\pi _s(f)`$, there is a set of probability at most $`\mathrm{exp}(2^{s_1})`$ such that outside this set we have
$$\underset{fF^{}}{sup}|Z_fZ_{\pi _{s_1}(f)}|\frac{\alpha u}{\sqrt{k}}\underset{s=s_1}{\overset{s_0}{}}2^{s/2}\pi _{s+1}(x)\pi _s(x)_{\psi _2}c^{}\alpha \frac{\gamma _2(F)}{\sqrt{k}}.$$
Denote $`F_{s_1}`$ by $`F^{\prime \prime }`$ and observe that $`|F^{\prime \prime }|2^{2^{s_1}}`$. Thus the later estimate implies
$$\underset{fF^{}}{sup}|Z_f|c^{}\alpha \frac{\gamma _2(F)}{\sqrt{k}}+\underset{gF^{\prime \prime }}{sup}|Z_g|.$$
Applying Lemma 1.2, for every $`w>0`$ we get
$$\left(|Z_g|\alpha ^2w\right)2\mathrm{exp}(c_1k\mathrm{min}(w,w^2)).$$
Thus, given $`w>0`$, choose $`s_1s_0`$ to be the largest integer such that $`2^{s_1}<c_1k\mathrm{min}(w,w^2)/2`$. Therefore, outside a set of probability less than or equal to $`|F^{\prime \prime }|\mathrm{\hspace{0.17em}2}\mathrm{exp}(c_1k\mathrm{min}(w,w^2))\mathrm{exp}(c_1k\mathrm{min}(w,w^2)/2)`$ we have $`|Z_g|\alpha ^2w`$ for all $`gF^{\prime \prime }`$. To conclude, outside a set of probability $`3\mathrm{exp}(c_1k\mathrm{min}(w,w^2)/2)`$,
$$\underset{fF^{}}{sup}|Z_f|c^{}\alpha \frac{\gamma _2(F)}{\sqrt{k}}+\alpha ^2w,$$
as required.
Proof of Theorem 1.4. Fix an arbitrary $`\rho >0`$, and for the purpose of this proof we let $`F(\rho )=F/\rho S_{L_2}`$, where $`F/\rho =\{f/\rho :fF\}`$.
Our first and main aim is to estimate $`sup_{fF(\rho )}|Z_f|`$ on a set of probability close to 1.
Fix $`u,w>0`$ to be determined later. Let $`F^{}F(\rho )`$ be as Lemma 1.5, with $`|F^{}|4^k`$. For every $`fF(\rho )`$ denote by $`\pi (f)=\pi _F^{}(f)`$ a closest point to $`f`$ with respect to the $`\psi _2`$ metric on $`F(\rho )`$. By writing $`f=(f\pi (f))+\pi (f)`$, it is evident that
$$|Z_f|W_{f\pi (f)}^2+2W_{f\pi (f)}W_{\pi (f)}+|Z_{\pi (f)}|,$$
and thus,
$$\underset{fF(\rho )}{sup}|Z_f|\underset{fF(\rho )}{sup}W_{f\pi (f)}^2+2\underset{fF(\rho )}{sup}W_{f\pi (f)}\underset{gF^{}}{sup}W_g+\underset{gF^{}}{sup}|Z_g|.$$
(1.6)
Applying Lemma 1.5, the first term in (1.6) is estimated using the fact that
$$\underset{fF(\rho )}{sup}W_{f\pi (f)}C\frac{\gamma _2(F(\rho ),_{\psi _2})}{\sqrt{k}},$$
with probability at least $`1\mathrm{exp}(k)`$.
For every $`fF(\rho )`$ and every $`u>0`$ we have
$$\left\{W_f1+u\alpha ^2\right\}\left\{W_f^21+u\alpha ^2\right\}\left\{|Z_f|u\alpha ^2\right\}$$
and, by Lemma 1.2, the latter probability is at most $`2\mathrm{exp}(ck\mathrm{min}(u,u^2))`$, where $`c>0`$ is an absolute constant.
Combining these two estimates with Lemma 1.7 and substituting into (1.6), $`sup_{fF(\rho )}|Z_f|`$ is upper bounded by
$`C^2{\displaystyle \frac{\gamma _2(F(\rho ),_{\psi _2})^2}{k}}+C{\displaystyle \frac{\gamma _2(F(\rho ),_{\psi _2})}{\sqrt{k}}}(1+u\alpha ^2)`$ (1.7)
$`+C^{\prime \prime }\alpha {\displaystyle \frac{\gamma _2(F(\rho ),_{\psi _2})}{\sqrt{k}}}+\alpha ^2w,`$
with probability at least $`12e^k2e^{ck\mathrm{min}(u,u^2)}3e^{ck\mathrm{min}(w,w^2)}`$.
Given $`0<\theta <1`$ we want the condition $`sup_{fF(\rho )}|Z_f|\theta `$ to be satisfied with probability close to 1. This can be achieved by imposing suitable conditions on the parameters involved. Namely, if $`u=1/\alpha ^2<1`$ and if $`\rho >0`$ and $`w>0`$ satisfy
$$\stackrel{~}{C}\alpha \frac{\gamma _2(F(\rho ),_{\psi _2})}{\sqrt{k}}\theta /4,C^{\prime \prime \prime }\alpha ^2w\theta /4,$$
(1.8)
where $`\stackrel{~}{C}=\mathrm{max}(2C,C^{\prime \prime })`$, then each of the last three terms in (1.7) is less than or equal to $`\theta /4`$, and the first term is less than or equal to $`(\theta /4)^2`$.
In order to ensure that (1.8) holds, we let $`w=\mathrm{min}(1,\theta /(4C^{\prime \prime \prime }\alpha ^2))`$. The above discussion shows that as long as $`\rho `$ satisfies
$$4\stackrel{~}{C}\alpha \frac{\gamma _2(F(\rho ),_{\psi _2})}{\theta \sqrt{k}}1,$$
(1.9)
then $`sup_{fF(\rho )}|Z_f|\theta `$ on a set of measure larger than or equal to $`17e^{ck\theta ^2/\alpha ^4}`$, where $`c>0`$ is an absolute constant. Hence, whenever $`\rho `$ satisfies (1.9) then (1.2) holds for all $`fF(\rho )`$. Finally, note that $`\gamma _2(F(\rho ),_{\psi _2})=(1/\rho )\gamma _2(F\rho S_{L_2},_{\psi _2})`$, and thus (1.9) is equivalent to the inequality in the definition of $`r_k^{}(\theta )`$.
To conclude the proof, for a fixed $`0<\theta <1`$ set $`r=r_k^{}(\theta )`$, with $`c=4\stackrel{~}{C}`$ being the constant from (1.9). Note that if $`X_1,\mathrm{},X_k`$ satisfy (1.2) for all $`fF(r)`$ then, since $`F`$ is star-shaped, the homogeneity of this condition implies that the same holds for all $`fF`$ with $`๐ผf^2r^2`$, as claimed.
Let us note two consequences for the supremum of the process $`Z_f`$, which is of independent interest.
###### Corollary 1.8
There exist absolute constants $`C^{},c^{}>0`$ for which the following holds. Let $`FS_{L_2}`$, $`\alpha =diam(F,_{\psi _2})`$ and $`k1`$. With probability at least $`1\mathrm{exp}(c^{}\gamma _2^2(F,_{\psi _2})/\alpha ^3)`$ one has
$$\underset{fF}{sup}|Z_f|C^{}\alpha \mathrm{max}(\frac{\gamma _2(F,_{\psi _2})}{\sqrt{k}},\frac{\gamma _2^2(F,_{\psi _2})}{k}).$$
Moreover, if $`F`$ is symmetric,
$$๐ผ\underset{fF}{sup}|Z_f|C^{}\alpha \mathrm{max}(\frac{\gamma _2(F,_{\psi _2})}{\sqrt{k}},\frac{\gamma _2^2(F,_{\psi _2})}{k}).$$
Proof. This follows from the proof of Theorem 1.4 with $`\rho =1`$. More precisely, the first part is a direct consequence of (1.7).
For the โmoreover partโ first use (1.6) for expectations, estimate the middle term by Cauchy-Schwarz inequality and note that $`W_g^21+Z_g`$ for all $`gF^{}`$ to yield that in order to estimate $`๐ผsup_{fF}|Z_f|`$ it suffice to bound
$$๐ผ\underset{fF^{}}{sup}|Z_f|\mathrm{and}๐ผ\underset{fF}{sup}W_{f\pi (f)}^2.$$
For simplicity denote $`\gamma _2(F,_{\psi _2})`$ by $`\gamma _2(F)`$ and let us begin with the second term. Applying Remark 1.6 and setting $`G=\{f\pi (f):fF\}`$ and $`u=c^{}\gamma _2(F)/\sqrt{k}`$, where $`c^{}`$ is the constant from the remark, we obtain
$`{\displaystyle _0^{\mathrm{}}}\left(\underset{gG}{sup}W_g^2t\right)๐t`$ $`u^2+{\displaystyle _{u^2}^{\mathrm{}}}\left(\underset{gG}{sup}W_g^2t\right)๐t`$
$`u^2+u^2{\displaystyle _1^{\mathrm{}}}\mathrm{exp}(c^{\prime \prime }vk).dv,`$
where the last inequality follows by changing the integration variable to $`t=u^2v`$. This implies that
$$๐ผ\underset{fF}{sup}W_{f\pi (f)}^2C^{}\frac{\gamma _2^2(F,_{\psi _2})}{k},$$
for some absolute constant $`C^{}>1`$.
Next, we have to bound $`๐ผsup_{fF^{}}|Z_f|`$, and to that end we use Lemma 1.7. Setting $`u=2C\alpha \gamma _2(F)/\sqrt{k}`$, and then changing the integration variable to $`t=u/2+\alpha ^2w`$, it is evident that
$`{\displaystyle _0^{\mathrm{}}}\left(\underset{fF^{}}{sup}|Z_f|t\right)๐t`$ $`u+{\displaystyle _u^{\mathrm{}}}\left(\underset{fF^{}}{sup}|Z_f|t\right)๐t`$
$`=`$ $`u+\alpha ^2{\displaystyle _{u/2\alpha ^2}^{\mathrm{}}}\left(\underset{fF^{}}{sup}|Z_f|{\displaystyle \frac{u}{2}}+\alpha ^2w\right)๐w`$
$``$ $`u+3\alpha ^2{\displaystyle _{u/2\alpha ^2}^{\mathrm{}}}\mathrm{exp}\left(c^{\prime \prime \prime }k\mathrm{min}(w^2,w)\right)๐w.`$
Changing variables in the last integral $`w=ru/2\alpha ^2`$ and using the fact that $`\gamma _2(F)1`$ for a symmetric set $`F`$ in the unit sphere, the last expression is bounded above by
$$u+(3/2)u_1^{\mathrm{}}\mathrm{exp}\left(c^{\prime \prime \prime }\mathrm{min}(\frac{r}{2\alpha ^2},\left(\frac{r}{2\alpha ^2}\right)^2)\right)๐r=C^{}u,$$
where $`C>0`$ is an absolute constant.
## 2 Subgaussian Operators
We now illustrate the general result of Section 1 in the case of linear processes, which was the topic that motivated our study. The processes corresponds then to random matrices with rows distributed according to measures on $`^n`$ satisfying some natural geometric conditions. Our result imply concentration estimates for related random subgaussian operators, which eventually lead to the desired reconstruction results for linear measurements for general sets.
The fundamental result that allows us to pass from the purely metric statement of the previous section to the geometric result we present below follows from Talagrandโs lower bound on the expectation of the supremum of a Gaussian process in terms of $`\gamma _2`$ of the indexing set. To present our result, the starting point is the fundamental definition of the $`\mathrm{}_{}`$-functional (which is in fact the so-called $`\mathrm{}`$-functional of a polar set).
###### Definition 2.1
Let $`T^n`$ and let $`g_1,\mathrm{},g_n`$ be independent standard Gaussian random variables. Denote by $`\mathrm{}_{}(T)=๐ผsup_{tT}\left|_{i=1}^ng_it_i\right|`$, where $`t=(t_i)_{i=1}^n^n`$.
There is a close connection between the $`\mathrm{}_{}`$\- and $`\gamma _2`$\- functionals given by the majorizing measure Theorem. Let $`\{G_t:tT\}`$ be a Gaussian process indexed by a set $`T`$, and for every $`s,tT`$, let $`d^2(s,t)=๐ผ|G_sG_t|^2`$. Then
$$c_2\gamma _2(T,d)๐ผ\underset{tT}{sup}|G_t|c_3\gamma _2(T,d),$$
where $`c_2,c_3>0`$ are absolute constants. The upper bound is due to Fernique \[F\] and the lower bound was established by Talagrand \[Ta1\]. The proof of both parts and the most recent survey on the topic can be found in \[Ta2\]
In particular, if $`T^n`$ and $`G_t=g_it_i`$, then $`d(s,t)=|st|`$, and thus
$$c_2\gamma _2(T,||)\mathrm{}_{}(T)c_3\gamma _2(T,||),$$
(2.1)
###### Definition 2.2
A probability measure $`\mu `$ on $`^n`$ is called isotropic if for every $`y^n`$, $`๐ผ|X,y|^2=|y|^2`$, where $`X`$ is distributed according to $`\mu `$.
A measure $`\mu `$ satisfies a $`\psi _2`$ condition with a constant $`\alpha `$ if for every $`y^n`$, $`X,y_{\psi _2}\alpha |y|`$.
A subgaussian or $`\psi _2`$ operator is a random operator of the form
$$\mathrm{\Gamma }=\underset{i=1}{\overset{k}{}}X_i,.e_i$$
(2.2)
where the $`X_i`$ are distributed according to an isotropic $`\psi _2`$ measure.
Perhaps the most important example of an isotropic $`\psi _2`$ probability measure on $`^n`$ with a bounded constant other than the Gaussian measure is the uniform measure on $`\{1,1\}^n`$. Naturally, if $`X`$ is distributed according to a general isotropic $`\psi _2`$ measure then the coordinates of $`X`$ need no longer be independent. For example, the normalized Lebesgue measure on an appropriate multiple of the unit ball in $`\mathrm{}_p^n`$ for $`2p\mathrm{}`$ is an isotropic $`\psi _2`$ measure with a constant independent of $`n`$ and $`p`$. For more details on such measures see \[MP\].
For a set $`T^n`$ and $`\rho >0`$ let
$$T_\rho =T\rho S^{n1}.$$
(2.3)
The next result shows that given $`T^n`$, subgaussian operators are very close to being an isometry on the subset of elements of $`T`$ which have a โlarge enoughโ norm.
###### Theorem 2.3
There exist absolute constants $`c,\overline{c}>0`$ for which the following holds. Let $`T^n`$ be a star-shaped set. Let $`\mu `$ be an isotropic $`\psi _2`$ probability measure with constant $`\alpha 1`$. Let $`k1`$, and $`X_1,\mathrm{},X_k`$ be independent, distributed according to $`\mu `$ and define $`\mathrm{\Gamma }`$ by (2.2). For $`0<\theta <1`$, with probability at least $`1\mathrm{exp}(\overline{c}\theta ^2k/\alpha ^4)`$, then for all $`xT`$ such that $`|x|r_k^{}(\theta )`$, we have
$$(1\theta )|x|^2\frac{|\mathrm{\Gamma }x|^2}{k}(1+\theta )|x|^2,$$
(2.4)
where
$$r_k^{}(\theta )=r_k^{}(\theta ,T):=inf\{\rho >0:\rho c\alpha ^2\mathrm{}_{}(T_\rho )/\left(\theta \sqrt{k}\right)\}.$$
(2.5)
In particular, with the same probability, every $`xT`$ satisfies
$$|x|^2\mathrm{max}\{(1\theta )^1|\mathrm{\Gamma }x|^2/k,r_k^{}(\theta )^2\}.$$
Proof. We use Theorem 1.4 for the set of functions $`F`$ consisting of linear functionals of the form $`f=f_x=,x`$, for $`xT`$. By the isotropicity of $`\mu `$, $`f_{L_2}=|x|`$ for $`f=f_xF`$. Also, since $`\mu `$ is $`\psi _2`$ with constant $`\alpha `$ then it follows by (2.1) that for all $`\rho >0`$,
$$\gamma _2(F\rho S_{L_2},_{\psi _2})\alpha \gamma _2(F\rho S_{L_2},_{L_2})(\alpha /c_2)\mathrm{}_{}(T_\rho ),$$
as promised.
###### Remark 2.4
It is clear from the proof of Theorem 1.4 that the upper estimates in (2.4) hold for $`\theta 1`$ as well, with appropriate probability estimates and a modified expression for $`r_k^{}`$ in (2.5). Note that of course in this case the lower estimate in (2.4) became vacuous. The same remark is valid for the estimate in (1.2) as well.
The last result immediately leads to an estimate for the diameter of random sections of a set $`T`$ in $`^n`$, given by kernels of random operators $`\mathrm{\Gamma }`$, and which is a $`\psi _2`$-counterpart of the main result from \[PT1\] (see also \[PT2\]).
###### Corollary 2.5
There exist absolute constants $`\stackrel{~}{c},\stackrel{~}{c}^{}>0`$ for which the following holds. Let $`T^n`$ be a star-shaped set and let $`\mu `$ be an isotropic $`\psi _2`$ probability measure with constant $`\alpha 1`$. Set $`k1`$, put $`X_1,\mathrm{},X_k`$ to be independent, distributed according to $`\mu `$ and define $`\mathrm{\Gamma }`$ by (2.2). Then, with probability at least $`1\mathrm{exp}(\stackrel{~}{c}k/\alpha ^4)`$,
$$diam(\mathrm{ker}\mathrm{\Gamma }T)r_k^{}(T),$$
where
$$r_k^{}=r_k^{}(T):=inf\{\rho >0:\rho \stackrel{~}{c}^{}\alpha ^2\mathrm{}_{}(T_\rho )/\sqrt{k}\}.$$
(2.6)
Moreover, with the same probability, $`diam(\mathrm{ker}\mathrm{\Gamma }T)\stackrel{~}{c}^{}\alpha ^2\mathrm{}_{}(T)/\sqrt{k}`$.
The Gaussian case (that is, when $`\mu `$ is the standard Gaussian measure on $`^n`$), although not explicitly stated in \[PT1\], follows immediately from the proof in that paper. The parameter $`inf\{\rho >0:\mathrm{}_{}(T\rho B_2^n)C\rho \sqrt{k}\}`$ was introduced in \[PT2\].
A version of Corollary 2.5 for random $`\pm 1`$-vectors follows from the result in \[A\], as observed in \[MP\].
Proof of Corollary 2.5. Applying Theorem 2.3 with $`\theta =1/2`$, say, we get that if $`xT`$ and $`|x|r_k^{}(1/2)`$ then $`\mathrm{\Gamma }x0`$. Thus for $`x\mathrm{ker}\mathrm{\Gamma }T`$ we have $`|x|r_k^{}(1/2)`$ and the first conclusion follows by adjusting the constants.
Observe that since the function $`\mathrm{}_{}(T_\rho )/\rho =\mathrm{}_{}\left((1/\rho )TS^{n1}\right)`$ is decreasing in $`\rho `$ then $`r_k^{}`$ actually satisfies the equality in the defining formula (2.6). Combining this and the obvious estimate $`\mathrm{}_{}(T_\rho )\mathrm{}_{}(T)`$, concludes the โmoreoverโ part.
Finally, let us note a special case of Theorem 2.3 for subsets of the sphere.
###### Corollary 2.6
Let $`TS^{n1}`$ and let $`\mu `$, $`\alpha `$, $`k`$, $`X_i`$, $`\mathrm{\Gamma }`$ and $`\theta `$ be the same as in Theorem 2.3. As long as $`k`$ satisfies $`k(c^{}\alpha ^4/\theta ^2)\mathrm{}_{}(T)^2`$, then with probability at least $`1\mathrm{exp}(\overline{c}\theta ^2k/\alpha ^4)`$, for all $`xT`$,
$$1\theta \frac{|\mathrm{\Gamma }x|}{\sqrt{k}}1+\theta ,$$
(2.7)
where $`c,\overline{c}>0`$ are absolute constants.
Proof. Let $`c,\overline{c}>0`$ be the constants from Theorem 2.3. Observe that the condition on $`k`$, with $`c^{}=c^2`$, is equivalent to $`r_k^{}(\theta ,\stackrel{~}{T})1`$, where $`\stackrel{~}{T}=\{\lambda x:xT,0\lambda 1\}`$. Then (2.7) immediately follows from (2.4).
## 3 Approximate reconstruction
Next, we show how one can apply Theorem 2.3 to reconstruct any fixed $`vT`$ for any set $`T^n`$, where the data at hand are linear subgaussian measurements of the form $`X_i,v`$.
The reconstruction algorithm we choose is as follows: for a fixed $`\epsilon >0`$, find some $`tT`$ such that
$$\left(\frac{1}{k}\underset{i=1}{\overset{k}{}}\left(X_i,vX_i,t\right)^2\right)^{1/2}\epsilon .$$
The fact that we only need to find $`t`$ for which $`(X_i,t)_{i=1}^k`$ is close to $`(X_i,v)_{i=1}^k`$ rather than equal to it, is very important algorithmically because it is a far simpler problem.
Let us show why such an algorithm can be used to solve the approximate reconstruction problem.
Consider $`\overline{T}=\{\lambda (ts):t,sT,0\lambda 1\}`$ and observe that by Theorem 2.3, for every $`0<\theta <1`$, with high probability, every such $`tT`$ satisfies that
$$|tv|\frac{\epsilon }{1\theta }+r_k^{}(\theta ,\overline{T}).$$
Hence, to bound the reconstruction error, one needs to estimate $`r_k^{}(\theta ,\overline{T})`$. Of course, if $`T`$ happens to be convex and symmetric then $`\overline{T}2T`$ which is star-shaped and thus
$$|tv|\frac{\epsilon }{1\theta }+r_k^{}(\theta ,2T).$$
In a more general case, when $`T`$ is symmetric and quasi-convex with constant $`a1`$, (i.e., $`T+T2aT`$ and $`T`$ is star-shaped), then
$$|tv|\frac{\epsilon }{1\theta }+r_k^{}(\theta ,aT).$$
Therefore, in the quasi-convex case, the ability to approximate any point in $`T`$ using this kind of random sampling depends on the expectation of the supremum of a Gaussian process indexed by the intersection of $`T`$ and a sphere of a radius $`\rho `$ as a function of the radius. For a general set $`T`$, the reconstruction error is controlled by the behavior of the expectation of the supremum of the Gaussian process indexed by the intersection of $`\overline{T}`$ with spheres of radius $`\rho `$, and this function of $`\rho `$ is just the modulus of continuity of the Gaussian process indexed by the set $`\{\lambda t:0\lambda 1,tT\}`$ (i.e, the expectation of the supremum of the Gaussian process indexed by the set $`\{\lambda (ts):0\lambda 1,t,sT,|ts|=\rho \}`$).
The parameters $`r_k^{}(\theta ,T)`$ can be estimated for the unit ball of classical normed or quasi-normed spaces. The two example we consider here are the unit ball in $`\mathrm{}_1^n`$, denoted by $`B_1^n`$, and the unit balls in the weak-$`\mathrm{}_p^n`$ spaces $`\mathrm{}_{p,\mathrm{}}^n`$ for $`0<p<1`$, denoted by $`B_{p,\mathrm{}}^n`$. Recall that $`B_{p,\mathrm{}}^n`$ is the set of all $`x=(x_i)_{i=1}^n^n`$ such that the cardinality $`|\{i:|x_i|s\}|s^p`$ for all $`s>0`$, and observe that $`B_{p,\mathrm{}}^n`$ is a quasi convex body with constant $`a=2^{1/p}`$. Let us mention that there is nothing โmagicalโ about the examples we consider here. Those are simply the cases considered in \[CT1, RV\].
In order to bound $`r_k^{}`$ for these sets we shall use the approach from \[GLMP\], and combine it with Theorem 2.3 to recover and extend the results from \[CT1, RV\].
###### Theorem 3.1
There is an absolute constant $`\overline{c}`$ for which the following holds. Let $`1kn`$ and $`0<\theta <1`$, and set $`\epsilon >0`$. Let $`\mu `$ be an isotropic $`\psi _2`$ probability measure on $`^n`$ with constant $`\alpha `$, and let $`X_1,\mathrm{},X_k`$ be independent, distributed according to $`\mu `$. For any $`0<p<1`$, with probability at least $`1\mathrm{exp}(\overline{c}\theta ^2k/\alpha ^4)`$, if $`v,yB_{p,\mathrm{}}^n`$ satisfy that $`(_{i=1}^kX_i,vy^2/k)^{1/2}\epsilon `$, then
$$|yv|\frac{\epsilon }{1\theta }+2^{1/p+1}\left(\frac{1}{p}1\right)^1\left(C_{\alpha ,\theta }\frac{\mathrm{log}(C_{\alpha ,\theta }n/k)}{k}\right)^{1/p1/2},$$
where $`C_{\alpha ,\theta }=c\alpha ^4/\theta ^2`$ and $`c>0`$ is an absolute constant.
If $`v,yB_1^n`$ satisfy the same assumption then with the same probability estimate,
$$|yv|\frac{\epsilon }{1\theta }+\left(C_{\alpha ,\theta }\frac{\mathrm{log}(C_{\alpha ,\theta }n/k)}{k}\right)^{1/2}.$$
To prove Theorem 3.1, we require the following elementary fact.
###### Lemma 3.2
Let $`0<p<1`$ and $`1mn`$. Then, for every $`x^n`$,
$$\underset{zB_{p,\mathrm{}}^n\rho B_2^n}{sup}x,z2\rho \left(\underset{i=1}{\overset{m}{}}x_{i}^{}{}_{}{}^{2}\right)^{1/2},$$
where $`\rho =\left(1/p1\right)^1m^{1/21/p}`$ and $`(x_i^{})_{i=1}^n`$ is a non-increasing rearrangement of $`(|x_i|)_{i=1}^n`$.
Moreover,
$$\underset{zB_1^n\rho B_2^n}{sup}x,z2\rho \left(\underset{i=1}{\overset{m}{}}x_{i}^{}{}_{}{}^{2}\right)^{1/2},$$
with $`\rho =1/\sqrt{m}`$.
Proof. We will present a proof for the case $`0<p<1`$. The case of $`B_1^n`$ is similar.
Recall a well known fact that for two sequences of positive numbers $`a_i,b_i`$ such that $`a_1a_2\mathrm{}`$, the sum $`a_ib_{\pi (i)}`$ is maximal over all permutations $`\pi `$ of the index set, if $`b_{\pi (1)}b_{\pi (2)}\mathrm{}`$. It follows that, for any $`\rho >0,m1`$ and $`zB_{p,\mathrm{}}^n\rho B_2^n`$,
$`x,z`$ $`\rho \left({\displaystyle \underset{i=1}{\overset{m}{}}}x_{i}^{}{}_{}{}^{2}\right)^{1/2}+{\displaystyle \underset{i>m}{}}{\displaystyle \frac{x_i^{}}{i^{1/p}}}`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{m}{}}}x_{i}^{}{}_{}{}^{2}\right)^{1/2}\left(\rho +{\displaystyle \frac{1}{\sqrt{m}}}{\displaystyle \underset{i>m}{}}{\displaystyle \frac{1}{i^{1/p}}}\right)`$
$``$ $`\left({\displaystyle \underset{i=1}{\overset{m}{}}}x_{i}^{}{}_{}{}^{2}\right)^{1/2}\left(\rho +\left({\displaystyle \frac{1}{p}}1\right)^1{\displaystyle \frac{1}{m^{1/p1/2}}}\right).`$
By the definition of $`\rho `$, this completes the proof.
Consider the set of elements in the unit ball with โshort supportโ, defined by
$$U_m=\{xS^{n1}:\left|\{i:x_i0\}\right|m\}.$$
Note that Lemma 3.2 combined with a duality argument implies that for every $`1mn`$ and every $`I\{1,\mathrm{},n\}`$ with $`|I|m`$,
$$\sqrt{|I|}B_1^nS^{n1}2convU_mS^{n1}.$$
(3.1)
The next step is to bound the expectation of the supremum of the Gaussian process indexed by $`U_m`$.
###### Lemma 3.3
There exist an absolute constant $`c`$ such that for every $`1mn`$,
$$\mathrm{}_{}(convU_m)c\sqrt{m\mathrm{log}(cn/m)}.$$
Proof. Recall that for every $`1mn`$, there is a set $`\mathrm{\Lambda }_m`$ of cardinality at most $`5^m`$ such that $`B_2^m2conv\mathrm{\Lambda }_m`$ (for example, a successive approximation shows that we may take as $`\mathrm{\Lambda }_m`$ an $`1/2`$-net in $`B_2^m`$). Hence there is a subset of $`B_2^n`$ of cardinality at most $`5^m\left(\genfrac{}{}{0pt}{}{n}{m}\right)`$ such that $`U_m2conv\mathrm{\Lambda }`$. It is well known (see for example\[LT\]) that for every $`TB_2^n`$,
$$\mathrm{}_{}(convT)=\mathrm{}_{}(T)c\sqrt{\mathrm{log}(|T|)},$$
and thus,
$$\mathrm{}_{}(convU_m)c\sqrt{\mathrm{log}\left(5^m\left(\genfrac{}{}{0pt}{}{n}{m}\right)\right)},$$
from which the claim follows.
Finally, we are ready to estimate $`r_k^{}(\theta ,B_{p,\mathrm{}}^n)`$ and $`r_k^{}(\theta ,B_1^n)`$.
###### Lemma 3.4
There exists an absolute constant $`c`$ such that for any $`0<p<1`$ and $`1kn`$,
$$r_k^{}(\theta ,B_{p,\mathrm{}}^n)c\left(\frac{1}{p}1\right)^1\left(\frac{\mathrm{log}(cn\alpha ^4/\theta ^2k)}{\theta ^2k/\alpha ^4}\right)^{1/p1/2}$$
and
$$r_k^{}(\theta ,B_1^n)c\left(\frac{\mathrm{log}(cn\alpha ^4/\theta ^2k)}{\theta ^2k/\alpha ^4}\right)^{1/2}.$$
Proof. Again, we present a proof for $`B_{p,\mathrm{}}^n`$, while the treatment of $`B_1^n`$ is similar and thus omitted.
Let $`0<p<1`$ and $`1kn`$, and set $`1mn`$ to be determined later. Clearly,
$$\left(\underset{i=1}{\overset{m}{}}x_{i}^{}{}_{}{}^{2}\right)^{1/2}=\underset{yU_m}{sup}x,y,$$
and thus, by Lemma 3.2, $`\mathrm{}_{}(B_{p,\mathrm{}}^n\rho B_2^n)2\rho \mathrm{}_{}(U_m)`$, where $`\rho =(1/p1)^1m^{1/21/p}`$. From the definition of $`r_k^{}(\theta )`$ in Theorem 2.3, it suffices to determine $`m`$ (and thus $`\rho `$) such that
$$c\alpha ^2\mathrm{}_{}(U_m)\theta \sqrt{k},$$
which by Lemma 3.3, comes to $`c\alpha ^2\sqrt{m\mathrm{log}(cn/m)}\theta \sqrt{k}`$ for some other numerical constant $`c`$. It is standard to verify that the last inequality holds true provided that
$$mc\frac{\theta ^2k/\alpha ^4}{\mathrm{log}\left(cn\alpha ^4/\theta ^2k\right)},$$
and thus
$$r_k^{}(\theta ,B_{p,\mathrm{}}^n)c\left(\frac{1}{p}1\right)^1\left(\frac{\mathrm{log}\left(cn\alpha ^4/\theta ^2k\right)}{\theta ^2k/\alpha ^4}\right)^{1/p1/2}.$$
Proof of Corollary 3.1. The proof follows immediately from Theorem 2.3 and Lemma 3.4.
## 4 Exact reconstruction
Let us consider the following problem from the error correcting code theory. A linear code is given by an $`n\times (nk)`$ real matrix $`A`$. Thus, a vector $`x^{nk}`$ generates the vector $`Ax^n`$. Suppose that $`Ax`$ is corrupted by a noise vector $`z^n`$ and the assumption we make is that $`z`$ is sparse, that is, has a short support, which we denote by supp$`(z)`$=$`\{i:z_i0\}`$. The problem is to reconstruct $`x`$ from the data, which is the noisy output $`y=Ax+z`$.
For this purpose, consider a $`k\times n`$ matrix $`\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }A=0`$. Thus $`\mathrm{\Gamma }z=\mathrm{\Gamma }y`$ and correcting the noise is reduced to identifying the sparse vector $`z`$ (rather than approximating it) from the data $`\mathrm{\Gamma }z`$ \- which is the problem we focus on here.
In this context, a linear programming approach called the basis pursuit algorithm, was recently shown to be relevant for this goal \[CDS\]. This method is the following minimization problem
$$(P)\mathrm{min}t_\mathrm{}_1,\mathrm{\Gamma }t=\mathrm{\Gamma }z$$
(and recall that the $`\mathrm{}_1`$-norm is defined by $`t_\mathrm{}_1=_{i=1}^n|t_i|`$ for any $`t=(t_i)_{i=1}^n^n`$).
For the analysis of the reconstruction of sparse vectors by this basis pursuit algorithm, we refer to \[CDS\] and the recent papers \[CT2, CT3\].
In this section, we show that if $`\mathrm{\Gamma }`$ is an isotropic $`\psi _2`$ matrix then with high probability, for any vector $`z`$ whose support has size less than $`\frac{Ck}{\mathrm{log}(Cn/k)}`$ (for some absolute constant $`C`$), the problem $`(P)`$ above has a unique solution that is equal to $`z`$. It means that such random matrices can be used to reconstruct any sparse vector, as long as the size of the support is not too large. This extends the recent result proved in \[CT2\] and \[RV\] for Gaussian matrices.
###### Theorem 4.1
There exist absolute constants $`c,C`$ and $`\overline{c}`$ for which the following holds. Let $`\mu `$ be an isotropic $`\psi _2`$ probability measure with constant $`\alpha 1`$. For $`1kn`$, set $`X_1,\mathrm{},X_k`$ to be independent, distributed according to $`\mu `$ and let $`\mathrm{\Gamma }=_{i=1}^kX_i,e_i`$. Then with probability at least $`1\mathrm{exp}(\overline{c}k/\alpha ^4)`$, any vector $`z`$ satisfying
$$|supp(z)|\frac{Ck}{\alpha ^4\mathrm{log}(cn\alpha ^4/k)}$$
is the unique minimizer of the problem
$$(P)\mathrm{min}t_\mathrm{}_1,\mathrm{\Gamma }t=\mathrm{\Gamma }z.$$
The proof of Theorem 4.1 is based on a scheme of the proof of \[CT2\], but is simpler and more general, as it holds for an arbitrary isotropic $`\psi _2`$ random matrix.
Let us remark that problem $`(P)`$ is equivalent to the following one
$$(P^{})\underset{t^n}{\mathrm{min}}yAt_\mathrm{}_1$$
where $`\mathrm{\Gamma }A=0`$. Thus we obtain the reconstruction result:
###### Corollary 4.2
Let $`A`$ be a $`n\times (nk)`$ matrix. Set $`\mathrm{\Gamma }`$ to be a $`k\times n`$ matrix that satisfies the conclusion of the previous Theorem with the constants $`c`$ and $`C`$, and for which $`\mathrm{\Gamma }A=0`$. For any $`x^{nk}`$ and any $`y=Ax+z`$, if $`|\mathrm{supp}(z)|\frac{Ck}{\mathrm{log}(cn/k)}`$, then $`x`$ is the unique minimizer of the problem
$$\underset{t^n}{\mathrm{min}}yAt_\mathrm{}_1.$$
The condition $`\mathrm{\Gamma }A=0`$ means that the range of $`A`$ is a subspace of the kernel of $`\mathrm{\Gamma }`$. Due to the rotation invariance of a Gaussian matrix (from both sides), the range and the kernel are random elements of the Grassmann manifold of the corresponding dimensions. Therefore, random Gaussian matrices $`A`$ and $`\mathrm{\Gamma }`$ satisfy the conclusion of Corollary 4.2.
### 4.1 Proof of Theorem 4.1
As in \[CT2\], the proof consists of finding a simple condition for a fixed matrix $`\mathrm{\Gamma }`$ to satisfy the conclusion of our Theorem. We then apply a result from the previous section to show that random matrices satisfy this condition.
The first step is to provide some criteria which ensure that the problem $`(P)`$ has a unique solution as specified in Theorem 4.1. This convex optimization problem can be represented as a linear programming problem. Indeed, let $`z^n`$ and set
$$I^+=\{i:z_i>0\},I^{}=\{i:z_i<0\},I=I^+I^{}.$$
(4.1)
Denote by $`๐`$ the cone of constraint
$$๐=\{t^n:\underset{iI^+}{}t_i\underset{iI^{}}{}t_i+\underset{iI^c}{}|t_i|0\}$$
corresponding to the $`\mathrm{}_1`$-norm.
Note that if $`|t|`$ is small enough then $`z+t_\mathrm{}_1=_{iI^+}(z_i+t_i)_{iI^{}}(z_i+t_i)+_{iI^c}|t_i|`$. Thus, the solution of $`(P)`$ is unique and equals to $`z`$ if and only if
$$\mathrm{ker}\mathrm{\Gamma }๐=\{0\}$$
(4.2)
By the Hahn-Banach separation Theorem, the latter is equivalent to the existence of a linear form $`\stackrel{~}{w}^n`$ vanishing on $`\mathrm{ker}\mathrm{\Gamma }`$ and positive on $`๐\{0\}`$.
After appropriate normalization, it is easy to check that such an $`\stackrel{~}{w}`$ satisfies that $`\stackrel{~}{w}=_{i=1}^k\alpha _iX_i`$, $`\stackrel{~}{w},e_i=1`$ for all $`iI^+`$, $`\stackrel{~}{w},e_i=1`$ for all $`iI^{}`$, and $`|\stackrel{~}{w},e_i|<1`$ for all $`iI^c`$. Setting $`w=_{i=1}^k\alpha _ie_i`$ and noticing that $`\stackrel{~}{w},e_i=w,\mathrm{\Gamma }e_i`$ we arrive at the following criterion.
###### Lemma 4.3
Let $`\mathrm{\Gamma }`$ be a $`k\times n`$ matrix and $`z^n`$. With the notation (4.1), the problem
$$(P)\mathrm{min}t_\mathrm{}_1,\mathrm{\Gamma }t=\mathrm{\Gamma }z$$
has a unique solution which equals to $`z`$, if and only if there exists $`w^k`$ such that
$$iI^+w,\mathrm{\Gamma }e_i=1,iI^{}w,\mathrm{\Gamma }e_i=1,iI^c|w,\mathrm{\Gamma }e_i|<1.$$
The second preliminary result we require follows from Corollary 2.6 and the estimates of the previous section.
###### Theorem 4.4
There exist absolute constants $`c,C`$ and $`\overline{c}`$ for which the following holds. Let $`\mu `$, $`\alpha `$, $`k`$ and $`\mathrm{\Gamma }`$ be as in Theorem 4.1. Then, for every $`0<\theta <1`$, with probability at least $`1\mathrm{exp}(\overline{c}\theta ^2k/\alpha ^4)`$, every $`x2convU_{4m}S^{n1}`$ satisfies that
$$(1\theta )|x|^2\frac{|\mathrm{\Gamma }x|^2}{k}(1+\theta )|x|^2,$$
(4.3)
provided that
$$mC\frac{\theta ^2k/\alpha ^4}{\mathrm{log}\left(cn\alpha ^4/\theta ^2k\right)}.$$
Proof. Applying Corollary 2.6 to $`T=2convU_{4m}S^{n1}`$, we only have to check that $`k(c^{}\alpha ^4/\theta ^2)\mathrm{}_{}(T)^2`$, which from Lemma 3.3 reduces to verifying that $`k(c^{}\alpha ^4/\theta ^2)cm\mathrm{log}(cn/m)`$. The conclusion now follows from the same computation as in the proof of Lemma 3.4.
Proof of Theorem 4.1. Observe that if $`t๐S^{n1}`$ then $`t_\mathrm{}_12_{iI}|t_i|2\sqrt{|I|}`$, where $`I`$ is the support of $`z`$. Hence,
$$๐S^{n1}\sqrt{4|I|}B_1^nS^{n1}.$$
This inclusion and condition (4.2) clearly imply that if $`\mathrm{\Gamma }`$ does not vanish on any point of $`\sqrt{4|I|}B_1^nS^{n1}`$, then the solution of $`(P)`$ is unique and equals to $`z`$. By (3.1) we have
$$\sqrt{4|I|}B_1^nS^{n1}2convU_{4m}S^{n1}.$$
Therefore, if $`\mathrm{\Gamma }`$ does not vanish on any point of $`2convU_{4m}S^{n1}`$ then $`z`$ is the unique solution of $`(P)`$. Applying Theorem 4.4, the lower bound in (4.3) shows that indeed, $`\mathrm{\Gamma }`$ does not vanish on any point of the required set, provided that
$$m\frac{Ck}{\alpha ^4\mathrm{log}(cn\alpha ^4/k)}$$
for some suitable constants $`c`$ and $`C`$.
### 4.2 The geometry of faces of random polytopes
Next, we investigate the geometry of random polytopes. Let $`\mathrm{\Gamma }`$ be a $`k\times n`$ isotropic $`\psi _2`$ matrix. For $`1in`$ let $`v_i=\mathrm{\Gamma }(e_i)`$ be the vector columns of the matrix $`\mathrm{\Gamma }`$ and set $`K^+(\mathrm{\Gamma })`$ (resp. $`K(\mathrm{\Gamma })`$) to be the convex hull (resp., the symmetric convex hull) of these vectors.
In this situation, the random model that makes sense is when $`X=(x_i)_{i=1}^n`$, where $`(x_i)_{i=1}^n`$ are independent, identically distributed random variables for which $`๐ผ|x_i|^2=1`$ and $`x_i_{\psi _2}\alpha `$. It is standard to verify that in this case $`X=(x_i)_{i=1}^n`$ is an isotropic $`\psi _2`$ vector with constant $`\alpha `$, and moreover, each vertex of the polytope is given by $`v_i=(x_{i,j})_{j=1}^k`$.
A polytope is called $`m`$-neighborly if any set of less than $`m`$ vertices is the vertex set of a face. In the symmetric setting, we will say that a symmetric polytope is $`m`$-symmetric-neighborly if any set of less than $`m`$ vertices containing no-opposite pairs, is the vertex set of a face.
The condition of Lemma 4.3 may be reformulated by saying that the set $`\{v_i:iI^+\}\{v_i:iI^{}\}`$ is the vertex set of a face of the polytope $`K(\mathrm{\Gamma })`$. Thus, the condition for the exact reconstruction using the basis pursuit method for any vector $`z`$ with $`|supp(z)|m`$ may be reformulated as a geometric property of the polytope $`K(\mathrm{\Gamma })`$ (see \[CT2, RV\]); namely, that for all disjoint subsets $`I^+`$ and $`I^{}`$ of $`\{1,\mathrm{},n\}`$ such that $`|I^+|+|I^{}|m`$, the set $`\{v_i:iI^+\}\{v_i:iI^{}\}`$ is the vertex set of a face of the polytope $`K(\mathrm{\Gamma })`$. That is, $`K(\mathrm{\Gamma })`$ is $`m`$-symmetric-neighborly. A similar analysis may be done in the non-symmetric case, for $`K^+(\mathrm{\Gamma })`$, where now $`I^{}`$ is empty.
###### Lemma 4.5
Let $`\mathrm{\Gamma }`$, $`K(\mathrm{\Gamma })`$ and $`K^+(\mathrm{\Gamma })`$ be as above. Then the problem
$$(P)\mathrm{min}t_\mathrm{}_1,\mathrm{\Gamma }t=\mathrm{\Gamma }z$$
has a unique solution which equals to $`z`$ for any vector $`z`$ (resp., $`z0`$) such that $`|\mathrm{supp}(z)|m`$, if and only if $`K(\mathrm{\Gamma })`$ (resp., $`K^+(\mathrm{\Gamma })`$) is $`m`$-symmetric-neighborly (resp., $`m`$-neighborly).
Applying Theorem 4.1, we obtain
###### Theorem 4.6
There exist absolute constants $`c,C`$ and $`\overline{c}`$ for which the following holds. Let $`\mu `$ be an isotropic $`\psi _2`$ probability measure with constant $`\alpha 1`$ and let $`k`$ and $`\mathrm{\Gamma }`$ be as above. Then, with probability at least $`1\mathrm{exp}(\overline{c}k/\alpha ^4)`$, the polytopes $`K^+(\mathrm{\Gamma })`$ and $`K(\mathrm{\Gamma })`$ are $`m`$-neighborly and $`m`$-symmetric-neighborly, respectively, for every $`m`$ satisfying
$$m\frac{Ck}{\alpha ^4\mathrm{log}(cn\alpha ^4/k)}$$
The statement of Theorem 4.6 for $`K(\mathrm{\Gamma })`$ and for a Gaussian matrix $`\mathrm{\Gamma }`$ is the main result of \[RV\]. However, a striking fact is that the same results holds for a random $`\{1,1\}`$-matrix. In such a case, $`K^+(\mathrm{\Gamma })`$ is the convex hull of $`n`$ random vertices of the discrete cube $`\{1,+1\}^k`$, also known as a random $`\{1,1\}`$-polytope. With high probability, every $`(m1)`$-dimensional face of $`K^+(\mathrm{\Gamma })`$ is a simplex and there are $`\left(\genfrac{}{}{0pt}{}{n}{m}\right)`$ such faces, for $`mCk/\mathrm{log}(cn/k)`$.
###### Remark 4.7
Let us mention some related results about random $`\{1,1\}`$-polytopes. A result of \[BP\] states that for such polytopes, the number of facets , which are the $`k1`$-dimensional faces, may be super-exponential in the dimension $`k`$, for an appropriate choice of the number $`n`$ of vertices. Denote by $`f_q(K^+(\mathrm{\Gamma }))`$ the number of $`q`$-dimensional faces of the polytope $`K^+(\mathrm{\Gamma })`$. The quantitative estimate in \[BP\] was recently improved in \[GGM\] where it is shown that there are positive constants $`a,b`$ such that for $`k^an\mathrm{exp}(bk)`$, one has $`๐ผf_{k1}(K^+(\mathrm{\Gamma }))(\mathrm{ln}n/\mathrm{ln}k^a)^{k/2}.`$ For lower dimensional faces a threshold of $`f_q(K^+(\mathrm{\Gamma }))/\left(\genfrac{}{}{0pt}{}{n}{q+1}\right)`$ was established in \[K\].
S. Mendelson Centre for Mathematics and its Applications, The Australian National University, Canberra, ACT 0200, Australia
shahar.mendelson@anu.edu.au
A. Pajor Laboratoire dโAnalyse et Mathรฉmatiques Appliquรฉes, Universitรฉ de Marne-la-Vallรฉe, 5 boulevard Descartes, Champs sur Marne, 77454 Marne-la-Vallee, Cedex 2, France
alain.pajor@univ-mlv.fr
N. Tomczak-Jaegermann Department of Mathematical and Statistical Sciences,
University of Alberta, Edmonton, Alberta, Canada T6G 2G1
nicole@ellpspace.math.ualberta.ca |
warning/0506/hep-ph0506153.html | ar5iv | text | # CKM Overview and Determinations from ๐ต Decays
## 1 Introduction
Flavour physics is at the moment the most active field of elementary particle physics. In particular the dedicated $`B`$-physics experiments BaBar and Belle produce a wealth of data which allows for a determination of fundamental parameters of the Standard Model with unprecedented accuracy. Flavour physics and CP violation are governed by the Cabibbo-Kobayashi-Maskawa (CKM) matrix$`^\mathrm{?}`$ which relates the flavour and mass eigenstates of quarks. The CKM-matrix elements appear as coupling constants in the charged-current transitions.
As a unitary complex $`3\times 3`$ matrix it has in principle nine real parameters, five of which can be eliminated due to phase redefinitions of the quarks. The three-generation CKM matrix therefore has three angles and one complex phase. The latter one is the only source of CP violation within the Standard Model (SM). There are many different ways of parametrizing the CKM matrix. For practical purposes most useful is the so called Wolfenstein parametrization$`^\mathrm{?}`$
$$V_{\mathrm{CKM}}=\left(\begin{array}{ccc}1\frac{\lambda ^2}{2}& \lambda & A\lambda ^3(\rho i\eta )\\ \lambda & 1\frac{\lambda ^2}{2}& A\lambda ^2\\ A\lambda ^3(1\rho i\eta )& A\lambda ^2& 1\end{array}\right)$$
(1)
which is an expansion to $`๐ช(\lambda ^3)`$ in the small parameter $`\lambda =|V_{us}|0.22`$. It is possible to improve the Wolfenstein parametrization to include higher orders of $`\lambda `$.$`^\mathrm{?}`$ The Wolfenstein parametrization makes the nearly diagonal structure of the CKM matrix obvious.
For phenomenological studies of CP-violating effects, the so called standard unitarity triangle (UT) plays a special role. It is a graphical representation of the orthogonality of the first and third column of the CKM matrix, namely
$$V_{ud}V_{ub}^{}+V_{cd}V_{cb}^{}+V_{td}V_{tb}^{}=0$$
(2)
in the $`(\rho ,\eta )`$ plane. This unitarity relation involves simultaneously the elements $`V_{ub}`$, $`V_{cb}`$, and $`V_{td}`$ which are under extensive discussion at present. The area of this and all other unitarity triangles equals half the absolute value of $`J_{CP}=\mathrm{Im}(V_{us}V_{cb}V_{ub}^{}V_{cs}^{})`$, the Jarlskog measure of CP violation.$`^\mathrm{?}`$ Usually, one chooses a phase convention where $`V_{cd}V_{cb}^{}`$ is real and rescales the above equation with $`|V_{cd}V_{cb}^{}|=A\lambda ^3`$. This leads to the triangle in Figure 1 with a base of unit length and the apex $`(\overline{\rho },\overline{\eta })`$.
A phase transformation in (2) only rotates the triangle, but leaves its form unchanged. Therefore, the angles and sides of the unitarity triangle are physical observables and can be measured. The angles $`\gamma `$ and $`\beta `$ are given by the arguments of $`V_{ub}^{}`$ and $`V_{td}^{}`$, respectively, whereas $`\alpha =\pi \beta \gamma `$ as a consequence of unitarity. Much effort is put into the determination of all the UT parameters. One tries to measure as many parameters as possible. The consistency of the various measurements tests the consequences of unitarity in the three-generation Standard Model. Any discrepancy with the SM expectations would imply the presence of new channels or particles contributing to the decay under consideration.
All angles and sides of the standard unitarity triangle can be determined from $`B`$ meson decays. The most prominent determination methods are presented in the following.
## 2 Unitarity triangle angles from $`B`$ decays
### 2.1 $`\mathrm{sin}2\beta `$ from $`BJ/\psi K_s`$
When a $`B`$ meson decays into a final state $`f_{CP}`$ which is an eigenstate of $`CP`$, the $`B`$ meson can decay either directly into $`f_{CP}`$ or mix into a $`\overline{B}`$ before decaying. The time-dependent CP asymmetry of such a decay takes a particularly simple form:
$`๐_{CP}(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\overline{B}^0f_{CP})\mathrm{\Gamma }(B^0f_{CP})}{\mathrm{\Gamma }(\overline{B}^0f_{CP})+\mathrm{\Gamma }(B^0f_{CP})}}`$ (3)
$`=`$ $`S_f^{\mathrm{mix}}\mathrm{sin}(\mathrm{\Delta }Mt)+C_f^{\mathrm{decay}}\mathrm{cos}(\mathrm{\Delta }Mt)`$
where $`\mathrm{\Delta }M`$ is the mass difference between the mass eigenstates in the $`B^0\overline{B}^0`$ system. The $`\mathrm{sin}(\mathrm{\Delta }Mt)`$ term with coefficient $`S_f^{\mathrm{mix}}`$ is due to interference of CP violation in mixing and decay, whereas the $`\mathrm{cos}(\mathrm{\Delta }Mt)`$ term with coefficient $`C_f^{\mathrm{decay}}`$ is non-vanishing when the decay directly violates CP. The latter is possible if at least two amplitudes with different weak and strong phases contribute to the decay under consideration. If, on the other hand, a single amplitude dominates, $`C_f=0`$. The prototype example of such a decay is the mode $`BJ/\psi K_s`$ whose time-dependent CP asymmetry is given by $`\mathrm{sin}(2\beta )\mathrm{sin}(\mathrm{\Delta }Mt)`$ where $`2\beta `$ is the weak phase of $`B^0\overline{B}^0`$ mixing.
The cleanest way to produce $`B`$ mesons is to operate an $`e^+e^{}`$ collider at the $`\mathrm{{\rm Y}}(4s)`$ resonance. The decay of the $`\mathrm{{\rm Y}}(4s)`$ produces with equal probability a $`B^+B^{}`$ or $`B^0\overline{B}^0`$ pair, where the latter is produced in a coherent quantum state. To measure a time-dependent CP asymmetry in neutral $`B`$-meson decays, it is necessary to determine the time difference $`\mathrm{\Delta }t`$ between the decays of the two mesons inside the pair. This can be achieved by measuring the distance between the two decay vertices which requires that the two $`B`$ mesons have a non-vanishing relative velocity. In the $`\mathrm{{\rm Y}}(4s)`$ rest frame the pair of $`B`$ mesons is produced almost at rest. In the asymmetric $`B`$ factories PEP-II$`^\mathrm{?}`$ and KEK-B$`^\mathrm{?}`$ the electron and positron beams have unequal energies which produces the $`B`$ meson pairs with a boost in the laboratory frame and therefore enables the measurement of time-dependent CP asymmetries.
The current world average of $`\mathrm{sin}2\beta `$ from $`BJ/\psi K_s`$ and other charmonium modes is$`^\mathrm{?}`$
$$\mathrm{sin}2\beta =0.725\pm 0.037$$
(4)
which is in perfect agreement with expectations from indirect measurements. The Standard-Model CKM picture is therefore established and the task is to look for small deviations from it.
### 2.2 $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ from $`bs`$ penguins
Decays like $`B\varphi K_s`$ are also dominated by a single decay amplitude: the $`bs`$ penguin transition which in the Standard Model has no weak phase. Therefore one expects that$`^\mathrm{?}`$
$$|S_{\psi K_s}S_{\varphi K_s}|_{\mathrm{SM}}0.04$$
(5)
such that the time-dependent CP asymmetry in $`B\varphi K_s`$ also measures $`\mathrm{sin}2\beta `$ to good approximation. The world average of $`\mathrm{sin}2\beta `$ from $`B\varphi K_s`$ and other s-penguin modes is$`^\mathrm{?}`$
$$\mathrm{sin}2\beta =0.43\pm 0.07$$
(6)
which is a $`3.7\sigma `$ deviation from the charmonium-mode value in (4). A persisting discrepancy with more data would be a rather clean hint towards new physics. Because the charmonium modes are governed by tree level decays whereas the $`bs`$ transitions are penguin decays, new physics could easily enter $`S_{\psi K_s}`$ and $`S_{\varphi K_s}`$ differently.
### 2.3 $`\alpha `$ from $`B\{\rho /\pi \}\{\rho /\pi \}`$
The quark-level transition $`\overline{b}\overline{u}`$ is accompanied by the CKM element $`V_{ub}=|V_{ub}|e^{i\gamma }`$ which carries the weak phase $`\gamma `$ allowing for direct CP violation. In the interference with $`B\overline{B}`$ mixing, the resulting weak phase is $`\beta +\gamma =\pi \alpha `$. The $`\overline{b}\overline{u}`$ transition appears for example in $`B\pi \pi `$, $`B\rho \pi `$, and $`B\rho \rho `$ decays. If this were the only contribution to these modes, their time-dependent CP asymmetry would directly measure the UT angle $`\alpha `$. However, the same final states can also be obtained via the $`bd`$ QCD-penguin diagram which carries the weak phase $`\beta `$. Hence, due to the interference of tree and penguin amplitudes, the time-dependent CP asymmetry in $`B\{\rho /\pi \}\{\rho /\pi \}`$ measures some $`\alpha _{\mathrm{eff}}`$ whose deviation from $`\alpha `$ depends on the relative size of tree and penguin contributions.
Using an isospin analysis proposed first by Gronau and London$`^\mathrm{?}`$, it is possible to extract $`\alpha `$ up to discrete ambiguities constructing the triangles
$`A_++\sqrt{2}A_{00}`$ $`=`$ $`\sqrt{2}A_{+0}`$ (7)
$`\overline{A}_++\sqrt{2}\overline{A}_{00}`$ $`=`$ $`\sqrt{2}\overline{A}_{+0}`$ (8)
Here the subscripts denote the charge of the final state mesons and the bar represents that a $`\overline{B}`$ instead of a $`B`$ meson is decaying. The Gronau-London method leads to a theoretically very clean determination of $`\alpha `$ but is experimentally extremely challenging because the small $`B\pi ^0\pi ^0`$ rates and CP asymmetries have to be measured.
So far, only upper bounds on the CP-averaged branching ratios $`B(B\pi ^0\pi ^0)`$ and $`B(B\rho ^0\rho ^0)`$ exist. Grossman and Quinn showed how these measurements together with the other $`B\{\rho /\pi \}\{\rho /\pi \}`$ branching ratios can be used to get an upper bound on the error on $`\mathrm{sin}2\alpha `$ due to penguin diagram effects.$`^\mathrm{?}`$ Because the bound on $`B(B\pi ^0\pi ^0)`$ is rather large, only the weak bound $`\alpha _{\mathrm{eff}}\alpha <35^{}`$ can be inferred. However, the bound on $`B(B\rho ^0\rho ^0)`$ is much smaller which indicates that penguin effects are small. Furthermore, it is dominated by the longitudinal polarization states such that the final state is a CP eigenstate to good approximation.
Even though $`\rho ^\pm \pi ^{}`$ are not CP eigenstates, one can extract $`\alpha `$ from $`B\rho \pi `$ decays using a Dalitz plot analysis.$`^\mathrm{?}`$ The currently best results for $`\alpha `$ come from this method and from $`S_{\rho ^+\rho ^{}}`$ leading to the world average$`^\mathrm{?}`$
$$\alpha =\left(100_{10}^{+12}\right)^{}$$
(9)
### 2.4 The angle $`\gamma `$
The angle $`\gamma `$ is the relative phase between $`V_{ub}`$ and $`V_{cb}`$ and can therefore be measured when $`bu`$ and $`bc`$ transitions can interfere in one decay. An analysis of the time-dependent CP asymmetry of $`B^0D^\pm \pi ^{}`$ for instance gives $`2\beta +\gamma `$. Although the final states are not CP eigenstates, the CKM-favored amplitude in $`B^0D^{}\pi ^+`$ can interfere with the doubly CKM-suppressed amplitude in $`B^0D^+\pi ^{}`$.$`^\mathrm{?}`$ The angle $`\gamma `$ can be extracted largely independent of new-physics effects from triangle relations between the tree-level decays $`B^\pm \{D^0,\overline{D}^0\}K^\pm fK^\pm `$.$`^{\mathrm{?},\mathrm{?}}`$ The classic proposal is by Gronau and Wyler where the $`D`$ mesons have to be reconstructed as CP eigenstates.$`^\mathrm{?}`$ This is experimentally hard because one has to measure small interference terms and is therefore not very sensitive to $`\gamma `$. Belle overcomes this difficulties using a Dalitz plot analysis of $`B^+D^0K^+`$ and fits to the interference pattern of $`D^0`$ and $`\overline{D}^0`$.$`^\mathrm{?}`$ Not accessible at the current $`B`$-meson factories are $`B_s`$ mesons. These are the โeldoradoโ for $`\gamma `$ determinations, offering many more possibilities, for example via the decay $`B_sK^{}K^{}`$.
### 2.5 CKM constraints using theory input
Probably the best strategy to extract information on the UT in future will be to use theory input. Fleischer, Gronau, Rosner, and many others suggested to use $`SU(2)`$ or $`SU(3)`$ flavour symmetry in addition to some plausible dynamical assumptions for an extraction of $`\gamma `$.$`^\mathrm{?}`$ Even better control over theoretical uncertainties can be obtained within the framework of QCD factorization where strong phases or the ratio of penguin to tree amplitudes can be calculated.$`^\mathrm{?}`$ Figure 2
shows the potential of using QCD-factorization predictions for non-leptonic charmless two-body $`B`$ decays to determine $`\gamma `$. Lacking a recent analysis, the value $`\gamma =\left(62_9^{+6}\right)^{}`$ extracted from data available in winter 2004 is quoted.
## 3 Unitarity triangle sides from $`B`$ decays
### 3.1 The absolute normalization: $`|V_{cb}|`$
The CKM matrix element $`|V_{cb}|`$ can be measured most precisely with inclusive or exclusive semileptonic $`bcl\nu `$ decays. To extract $`|V_{cb}|`$ from the differential decay spectrum in exclusive $`BD^{()}l\overline{\nu }`$ decays, knowledge of the form factor $`_{D^{()}}(1)`$ at zero recoil is needed. These form factors are genuinely non-perturbative quantities which have to be calculated using e.g. lattice QCD or QCD sum rules. The semileptonic width $`\mathrm{\Gamma }(BX_cl\nu )`$ of the inclusive decay, on the other hand, can be calculated using an operator product expansion (OPE) which is a simultaneous expansion in $`\alpha _s(m_b)`$ and $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$. The necessary non-perturbative HQET parameters can be extracted measuring moments of the lepton energy or hadronic invariant mass. A recent global fit to data gives$`^\mathrm{?}`$
$$|V_{cb}|=(41.3\pm 0.6\pm 0.1)10^3$$
(10)
### 3.2 $`|V_{ub}|`$ from $`BX_ul\nu `$
$`|V_{ub}|`$ can similarly be extracted from either exclusive or inclusive $`bul\nu `$ decays. For the inclusive semileptonic $`\overline{B}X_ul^{}\overline{\nu }`$ decays tight experimental cuts are necessary to discriminate against the large charm background. These cuts restrict the hadronic final state to the shape-function region with large energy $`E_Xm_B`$ but only moderate invariant mass $`M_X\sqrt{m_B\mathrm{\Lambda }_{\mathrm{QCD}}}`$. The phase space can be depicted easiest in the hadronic variables $`P_\pm =E_H\left|\stackrel{}{P}_H\right|`$, the energy of the hadronic final state minus or plus the absolute value of its three momentum.$`^\mathrm{?}`$ Figure 3
shows the distribution of events in the phase space for the variables $`P_{}`$ and $`P_+`$. The vast majority of events is located in the shape-function region of small $`P_+`$ and large $`P_{}`$.
In order to discriminate against the large charm background, one has to apply tight experimental cuts. So far, cuts on the charged-lepton energy $`E_l>(M_B^2M_D^2)/(2M_B)`$, the hadronic invariant-mass squared $`s_H<M_D^2`$, and the dilepton mass squared $`q^2>(M_BM_D)^2`$ have been employed. Only the $`E_l`$ cut can be applied without neutrino reconstruction. Unfortunately, it has a very low efficiency and is therefore theoretically disfavored. The hadronic-mass cut is in principle the ideal separator between $`\overline{B}X_ul^{}\overline{\nu }`$ and $`\overline{B}X_cl^{}\overline{\nu }`$ events. However, in practice one has to lower the cut due to the experimental resolution on $`s_H`$, thereby reducing the efficiency. Cutting on the hadronic variable $`P_+`$ provides a new method for a precision measurement of $`|V_{ub}|`$ which combines good theoretical control with high efficiency and a powerful discrimination against charm background.$`^{\mathrm{?},\mathrm{?}}`$ The fraction $`F_P`$ of all $`\overline{B}X_ul\overline{\nu }`$ events with hadronic light-cone momentum $`P_+\mathrm{\Delta }_P`$ can be calculated systematically using a two-step matching of QCD current correlators onto soft-collinear and heavy-quark effective theory. The prediction for the fraction of events with the optimal cut $`P_+M_D^2/M_B`$ is$`^\mathrm{?}`$
$$F_P=(79.6\pm 10.8\pm 6.2\pm 8.0)\%$$
(11)
where the errors represent the sensitivity to the shape function, an estimate of $`๐ช(\alpha _s^2)`$ contributions, and power corrections, respectively. The CKM-matrix element $`|V_{ub}|`$ can be extracted by comparing a measurement of the partial rate $`\mathrm{\Gamma }_u(P_+\mathrm{\Delta }_P)`$ with a theoretical prediction for the product of the event fraction $`F_P`$ and the total inclusive $`\overline{B}X_ul\overline{\nu }`$ rate. The resulting theoretical uncertainty on $`|V_{ub}|`$ is$`^\mathrm{?}`$
$$\frac{\delta |V_{ub}|}{|V_{ub}|}=(\pm 7\pm 4\pm 5\pm 4)\%$$
(12)
where the last error comes from the uncertainty in the total rate. At leading power in $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ the shape-function uncertainty could be eliminated by relating the $`P_+`$ spectrum to the $`\overline{B}X_s\gamma `$ photon spectrum, both given at tree level directly by the shape function. With a 5% relative theoretical error on $`|V_{ub}|`$, corrections suppressed by a power of $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ are considered the second largest source of uncertainty. Subleading shape functions have been investigated first by Bauer, Luke, and Mannel$`^\mathrm{?}`$ and more recently in $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ using the two-step matching procedure developed in $`^{\mathrm{?},\mathrm{?}}`$. At tree level, the results of this SCET analysis can be expressed in terms of three subleading shape functions defined via the Fourier transforms of forward matrix elements of bi-local light-cone operators in heavy-quark effective theory. Numerically, power corrections indeed have the estimated 10% effect on the value of $`F_P(\mathrm{\Delta }_P)`$ near the charm threshold .
### 3.3 $`|V_{td}|`$ and $`|V_{ts}|`$
Finally, we want to discuss the extraction of the CKM matrix elements $`|V_{td}|`$ and $`|V_{ts}|`$. The standard way to determine these quantities is via the mass difference $`\mathrm{\Delta }M_{d/s}`$ between the mass eigenstates in the $`B_d^0`$-$`\overline{B}_d^0`$ and $`B_s^0`$-$`\overline{B}_s^0`$ systems, respectively. The quantities
$$\mathrm{\Delta }M_q=\frac{G_F^2}{6\pi ^2}\eta _Bm_{B_q}B_{B_q}F_{B_q}^2M_W^2S_0(x_t)|V_{tq}|^2$$
(13)
are directly proportional to $`|V_{tq}|^2`$. In the above formula, $`G_F`$ is the Fermi constant, $`\eta _B=0.55\pm 0.01`$ is a QCD factor, $`B_{B_q}`$ are the so-called bag parameters, $`F_{B_q}`$ is the $`B_q`$-meson decay constant, and $`S_0(x_t)`$ is the Inami-Lim function for the box diagram with a top-quark exchange. Bag parameter and decay constant are non-perturbative input parameters which have to be extracted for instance from lattice QCD data.
$`|V_{td}|`$ and $`|V_{ts}|`$ are also accessible in the radiative decays $`BX_{d/s}\gamma `$ which are mediated predominantly via the electromagnetic penguin operator $`Q_7`$. Whereas the inclusive $`BX_s\gamma `$ branching ratio was measured experimentally to very high accuracy, it is nearly impossible to achieve a corresponding measurement of the inclusive $`bd\gamma `$ branching ratio. The exclusive branching ratios $`B\rho \gamma `$ are experimentally easier accessible. However, they are theoretically more involved than the inclusive modes, because bound state effects have to be taken into account. Yet, a systematic and model-independent analysis of exclusive radiative decays is possible in the heavy-quark limit $`m_b\mathrm{\Lambda }_{\mathrm{QCD}}`$. The relevant hadronic matrix elements of local operators in the weak Hamiltonian simplify in this limit because perturbatively calculable hard-scattering kernels can be separated from non-perturbative form factors and universal light-cone distribution amplitudes. This approach to exclusive radiative decays$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ is similar in spirit to the treatment of hadronic matrix elements in two-body non-leptonic $`B`$ decays formulated by Beneke, Buchalla, Neubert, and Sachrajda.$`^\mathrm{?}`$ The non-perturbative form factors represent the biggest source of theoretical uncertainty in the prediction of the exclusive branching ratios. A very interesting strategy to reduce the form-factor uncertainties is to use ratios like the one of CP-averaged branching ratios
$$R_0=\frac{B(B^0\rho ^0\gamma )+B(\overline{B}^0\rho ^0\gamma )}{B(B^0K^0\gamma )+B(\overline{B}^0\overline{K}^0\gamma )}$$
(14)
the isospin breaking ratio
$$\mathrm{\Delta }(\rho \gamma )=\frac{2\mathrm{\Gamma }(B^0\rho ^0\gamma )\mathrm{\Gamma }(B^\pm \rho ^\pm \gamma )}{2\mathrm{\Gamma }(B^0\rho ^0\gamma )+\mathrm{\Gamma }(B^\pm \rho ^\pm \gamma )}$$
(15)
or the ratio of $`B\rho l\nu `$ and $`B\rho \gamma `$ decay rates. In these ratios only the ratios of the form factors enter. These form-factor ratios are known in certain limiting cases, like flavour $`SU(3)`$ or isospin symmetry, or in the large-energy limit.
The form-factor ratio $`\xi =F_K^{}/F_\rho `$ which enters $`R_0`$ for example differs from unity only because of $`SU(3)`$-breaking effects. Within the framework of QCD factorization the ratio $`R_0`$ can be calculated at next-to-leading order in $`\alpha _s`$ and to leading order in the heavy-quark limit. This ratio measures to very good approximation the side $`R_t=\sqrt{(1\overline{\rho })^2+\overline{\eta }^2}`$ of the standard unitarity triangle and the theoretical uncertainty in the relation to $`R_t`$ comes in essence solely from $`\xi `$.$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ From light-cone sum rule and lattice QCD calculations we choose $`\xi =1.2\pm 0.1`$ as averaged input value for the form-factor ratio $`\xi `$.$`^\mathrm{?}`$ The current experimental status of the branching ratio measurements is$`^\mathrm{?}`$
$`B(B^0K^0\gamma )`$ $`=`$ $`(4.01\pm 0.20)10^5`$ (16)
$`B(B^0\rho ^0\gamma )`$ $`<`$ $`0.410^6`$ (17)
which leads to $`R_0<0.01`$. This implies
$$R_t<0.76\frac{\xi }{1.2},\left|\frac{V_{td}}{V_{ts}}\right|<0.17\frac{\xi }{1.2},|V_{td}|<6.710^3\frac{\xi }{1.2}$$
(18)
The implication of this bound in the $`(\overline{\rho },\overline{\eta })`$ plane is shown in the left-hand plot of Figure 4.
We see that the current upper bound for $`B(B^0\rho ^0\gamma )`$ already cuts into the standard fit region for the apex of the unitarity triangle. Values of $`\overline{\rho }`$ and $`\overline{\eta }`$ to the left of the curved band are excluded at 90 % confidence level. The right-hand plot shows in addition the constraints coming from an assumed measurement of the isospin breaking ratio $`\mathrm{\Delta }(\rho \gamma )_{\mathrm{exp}}=0`$, which would correspond to the Standard Model prediction for a CKM angle $`\gamma =60^{}`$.
Another way to reduce hadronic uncertainties is to use the ratio of $`B\rho l\nu `$ and $`B\rho \gamma `$ decay rates. The simplification occurs because relations exist between the corresponding form factors in the large-energy limit. Since only $`B\rho `$ transitions are involved, the problems with $`SU(3)`$ breaking are avoided. The ratio of $`B\rho l\nu `$ events in the part of phase space where the large-energy limit is valid, divided by the $`B\rho \gamma `$ branching ratio then measures the CKM quantity
$$\left|\frac{V_{ud}V_{ub}}{V_{td}V_{tb}}\right|^2=\frac{\overline{\rho }^2+\overline{\eta }^2}{(1\overline{\rho })^2+\overline{\eta }^2}$$
(19)
With $`10^9`$ $`B\overline{B}`$ pairs available at the end of the Belle and BaBar experiments, a measurement of $`|V_{ud}V_{ub}/V_{td}^2|`$ with an error of $`\pm 25\%`$ could be achieved. The corresponding constraint in the $`(\overline{\rho },\overline{\eta })`$ plane is shown in Figure 5.$`^\mathrm{?}`$
We observe that the constraint is quite stringent, in particular in the important region corresponding to the standard fit results.
## 4 Conclusions
The Standard-Model CKM picture survived its first major test with flying colours: The constraints from the lengths of the sides and those from $`\mathrm{sin}2\beta `$ (and $`\epsilon _K`$) point to the same region for the apex of the unitarity triangle. This can be seen nicely in Figure 6
which shows the most recent global CKM fit in the $`(\overline{\rho },\overline{\eta })`$ plane.$`^\mathrm{?}`$ Now the task is to look for corrections to the standard picture, rather than for alternatives. The only noteworthy discrepancy from theory expectations is the extraction of $`\mathrm{sin}2\beta `$ from $`bs`$ penguin modes. If the currently large deviation of $`|S_{\psi K_s}S_{bs}|`$ from $`0`$ persists with more data, this would be a rather clean hint towards new physics. The angle $`\gamma `$ and $`|V_{ub}|`$ are at the moment the most important theory targets. New extraction techniques for all unitarity triangle parameters are still emerging. During the last years, great progress was made in $`B`$ physics but no big surprises emerged. The discovery of new physics in the flavour sector seems to be very challenging, both theoretically and experimentally.
## Acknowledgments
I would like to thank the organizers of the XXXXth Rencontres de Moriond for inviting me to such an enriching conference. Thank you to all the participants for interesting talks, good discussions and great skiing. I acknowledge financial support from the EU Marie Curie Conferences programme. This work was supported in part by the DFG Sonderforschungsbereich/Transregio 9 โComputer-gestรผtzte Theoretische Teilchenphysikโ. I would like to express my warmest thanks to the Benedictine Archabbey of St. Ottilien for hospitality while these proceedings were completed.
โ U. โ I. โ O. โ G. โ D. โ
## References |
warning/0506/hep-ph0506001.html | ar5iv | text | # Chromoelectric fields and quarkonium-hadron interactions at high energies
## I Introduction
Quarkonium-hadron cross sections ($`\sigma _{\mathrm{\Phi }h}`$) are a necessary tool to understand the forthcoming data on quarkonium production, which will become available at RHIC. In the last six years many efforts have been devoted to this problem nos1 and real progress has been achieved, especially in what concerns the cross sections at low energies, close to the dissociation threshold. In the energy region far from threshold the situation is less clear and even the energy dependence is still subject of debate. Extrapolation from calculations valid at low energies points to different directions. Results obtained with the non-relativistic quark model wongs indicate a rapidly falling cross section. This behavior is due to the gaussian tail of the quark wave functions used in the quark exchange model. This same behavior could be found within chiral meson Lagrangian approaches with the introduction of $`\sqrt{s}`$ dependent form factors osl . In QCD sum rules regras the cross section was found to be monotonically increasing with energy.
The calculations of $`\sigma _{\mathrm{\Phi }h}`$ designed to be valid at high energies ($`\sqrt{s}20`$ GeV) are quite few: the Bhanot-Peskin (BP) approach pes ; kar ; kar2 ; arleo ; oh ; song , perturbative QCD plus geometrical extrapolation gerland , the model of the stochastic vacuum (MSV) dnnr and the light-cone dipole formalism huf2000 . During the last years the leading order BP approach has been used very often. However, the recent next to leading order calculations presented in song show that, for the charmonium, the formalism breaks down because this system is not heavy enough. Most of the calculations mentioned above predict a rising cross section. In Ref. oh , $`\sigma _{\mathrm{\Phi }h}`$ falls with the energy and in Ref. dnnr it stays constant.
If the quarkonium is treated as an ordinary hadron, its cross section for interaction with any other ordinary hadron must increase smoothly at higher energies, in much the same way as the proton-proton or pion-proton cross sections. The underlying reason is the increasing role played by perturbative QCD dynamics and the manifestation of the partonic nature of all hadrons. However this partonic picture starts to be dominant only at much higher energies ($`\sqrt{s}\mathrm{\hspace{0.17em}100}`$ GeV). In the energy region relevant for RHIC physics non-perturbative aspects are still very important. In the high energy calculations mentioned above, different non-perturbative ingredients were employed: moments of the gluon distribution in the hadron pes ; kar ; kar2 ; arleo ; oh ; song ; hadron and quarkonium wave functions gerland and QCD vacuum expectation values (condensates) dnnr .
Since there are still discrepancies concerning numbers (which may vary by one order of magnitude for different estimates) and the energy behavior, we think that it is interesting to calculate $`\sigma _{\mathrm{\Phi }h}`$ with a non-perturbative approach, putting emphasis on the role played by the chromoelectric fields. In benzahra a similar treatment was adopted to study the quarkonium dissociation inside a QCD plasma. The color electric fields appearing in the transition matrix element were related to the color charge density of the medium, which, in turn, was computed in a specific model of the QGP. Here we start with a similar expression for the transition amplitude but, since we are in a purely hadronic phase, we must know the chromoelectric field inside nucleons and pions. There has been progress in the study of these fields, coming from models of the QCD vacuum michael , from lattice QCD suga , from the Field Correlator Method (FCM) simonov and from Coulomb gauge QCD adam . We hope that we can benefit from these advances and use the profiles of the chromoelectric fields estimated in these works in our problem. For this purpose, we treat the interaction between the quarkonium and hadron as being analogous to the interaction of a small dipole traversing a large capacitor and interacting with the color electric field but not with its sources. In the final part of this work we discuss the validity of this last assumption. Using a contact interaction between a heavy quark (or antiquark) and a quark (or antiquark) we compute the corresponding cross section and find that it is indeed much smaller than the heavy quark-external field cross section. The model developed here bears some resemblance to the Bhanot-Peskin picture, but is much simpler. Some simplifying assumptions are used to render the calculations quasi analytic and preserve the understanding of the basic physics.
## II The model
### II.1 The interaction Hamiltonian
The starting point is the assumption that the quarkonium (dipole) is small compared with the hadron (capacitor). As a consequence, the $`\overline{Q}Q`$ pair will interact mostly with the external color field but not with the (quark) sources. Moreover, the external color field is considered to have only low momentum components (โsoft gluonsโ) and thus is able to transfer only a small amount of energy, which will be barely enough to dissociate the bound state. In the case of the charmonium, the typical binding energy is $`ฯต0.6`$ GeV. Therefore, in a first approximation
$$ฯตM_\mathrm{\Phi }$$
(1)
where $`M_\mathrm{\Phi }`$ is the mass of the bound state ($`M_\mathrm{\Phi }3`$ GeV). In the case of the bottonium this approximation is even better. The binding energy is also small compared to the collision energy
$$ฯต\sqrt{s}$$
(2)
Inequality (1) justifies the use of quantum mechanical perturbation theory (the Born approximation) and inequality (2) justifies the use of the eikonal approximation, which, in this case, implies that the hadron follows a straight line trajectory and remains essentially undisturbed during the interaction. In Figure 1 we present our picture of the scattering and our choice of coordinates, in the quarkonium rest frame: $`\stackrel{}{r}_1`$ and $`\stackrel{}{r}_2`$ are the quark and antiquark coordinates and $`\stackrel{}{E}^a`$ is the chromoelectric field in the projectile, which will be a proton or a pion, moving with constant velocity $`\stackrel{}{v}`$ at impact parameter $`\stackrel{}{b}`$.
With these assumptions we can write the interaction Hamiltonian as:
$$H_{int}=g(T_1^a\stackrel{}{E}_1^a\stackrel{}{r}_1+\overline{T}_2^b\stackrel{}{E}_2^b\stackrel{}{r}_2)$$
(3)
where $`T^a`$ ($`\overline{T}^b`$) are the generators of color group SU(3) in the fundamental (conjugate) representation. $`\stackrel{}{E}_1^a`$ and $`\stackrel{}{E}_2^b`$ are the chromoelectric fields generated by the hadron in motion (capacitor) and โfeltโ by quark and antiquark in the bound state respectively. They have to be Lorentz transformed to the quarkonium rest frame, bringing to our calculation a Lorentz gamma factor, which is the source of the energy dependence ($`\sqrt{s}`$) of our results. We shall for the moment neglect the magnetic component, since it does not do any work on the charges and thus is not effective in the energy transfer. Besides, the magnetic interaction is inversely proportional to the quark mass, being thus suppressed.
We can represent this external field by:
$`\stackrel{}{E}^a(r_e,t)`$ $`=`$ $`\gamma \stackrel{}{E}_0^aexp\left({\displaystyle \frac{(Xx_e)^2}{d^2}}{\displaystyle \frac{(Yy_e)^2}{d^2}}\right)exp\left(\gamma ^2{\displaystyle \frac{\left[vtz_e\right]^2}{d^2}}\right)`$ (4)
with $`e=1,2`$ . X, Y and Z are the hadron coordinates and $`\gamma `$ is the usual Lorentz factor. $`Z=vt`$, because the hadron moves with velocity $`\stackrel{}{v}`$ along the z axis. $`\stackrel{}{E}_0^a`$, which will be abreviated by $`E`$, is the color electric field at the center of the projectile. The projectile mean square radius is related to the parameter $`d`$ through:
$$\sqrt{r_h^2}=\mathrm{\hspace{0.17em}0.86}d$$
We neglect the deflection of the hadron trajectory, because we are studying reactions in the high energy and non-perturbative regime, i.e., with low momentum transfer. X and Y are related with the impact parameter $`b`$ by: $`b^2=X^2+Y^2`$. Notice that, by simplicity, we choose one preferencial direction for the field, in this case, the $`x`$-axis.
Neglecting the CM motion, (3) can be rewritten as
$$H_{int}=g(\frac{\lambda ^a}{2}E_1^a+\frac{\lambda ^{b^T}}{2}E_2^b)(\frac{x_1x_2}{2})$$
(5)
Also for the sake of simplicity, when working with (5), we will take $`x_1x_2a`$, where $`a`$ is the typical separation between quark and antiquark. Initially the quark-antiquark pair is in a localized region of the space.
### II.2 The initial state
The initial wave function of system has spatial and color parts defined by:
$$\mathrm{\Psi }_i=f(r_1,r_2)c_nd_n$$
(6)
where $`c_n`$ and $`d_n`$, with $`n=1,2,3`$, are the initial color vectors griff for quark and antiquark respectively, taken in a color singlet state. We choose
$$f(r_1,r_2)=N_iexp[\frac{\stackrel{}{r}_1^2}{a^2}]exp[\frac{\stackrel{}{r}_2^2}{a^2}]exp(i\epsilon _it)$$
(7)
where $`\epsilon _i`$ ($`\epsilon _i=M_\mathrm{\Phi }`$) is the quarkonium initial energy and $`N_i`$ is a normalization constant given by:
$$N_i^2=(\frac{2}{\pi })^3\frac{1}{a^6}$$
The initial wave function $`\mathrm{\Psi }_i`$ describes the confinement of quarks and also asymptotic freedom, as it allows the quarks to be independent inside the bag. It is easy to see that the connection between the quarkonium mean square radius and the parameter $`a`$ is
$$\sqrt{r_{Q\overline{Q}}^2}=\mathrm{\hspace{0.17em}1.09}a$$
### II.3 The final state
Under the action of the external field the initial wave function $`\mathrm{\Psi }_i`$ evolves to a final state $`\mathrm{\Psi }_f`$:
$$\mathrm{\Psi }_f=t(r_1,r_2)c_jd_k$$
(8)
where $`c_j`$ and $`d_k`$, with $`j,k=1,2,3`$, are the quark and antiquark final color vectors and $`t(r_1,r_2)`$ the spatial part of the wave function. In the final state of this reaction we have to deal with the transition of a pair of an excited quark and an antiquark to a pair of mesons $`D`$ \- $`\overline{D}`$ (or $`B\overline{B}`$). This transition is highly non-perturbative and has to be modelled. We shall use here two approaches.
model A
We first assume that the quark and antiquark are converted into two free mesons (a $`M`$ and a $`\overline{M}`$) which are thus described by plane waves:
$`t_A(r_1,r_2)`$ $`=`$ $`N_Aexp(i\stackrel{}{p}_1.\stackrel{}{r}_1)exp(i\stackrel{}{p}_2.\stackrel{}{r}_2)exp(i\epsilon _ft)`$ (9)
where $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}_2`$ are the meson momenta and $`N_A`$ is a normalization constant given by:
$$N_A^2=\frac{1}{V^2}$$
with $`V`$ being an arbitrary normalization volume, which will be cancelled in the calculation of the cross section. In the above expression $`\epsilon _f`$ is the final energy of the $`Q\overline{Q}`$ pair. The energy transferred during the reaction must be sufficient to dissociate the bound state into a pair of mesons with open charm ($`D\overline{D}`$) or beauty ($`B\overline{B}`$) and therefore:
$$\epsilon _f=\sqrt{(\stackrel{}{p}_1)^2+m_M^2}+\sqrt{(\stackrel{}{p}_2)^2+m_{\overline{M}}^2}$$
(10)
where $`m_M`$ ($`m_{\overline{M}}`$) is the mass of the meson coming out from the fragmentation of the quark (antiquark). With this definition of $`\epsilon _f`$ we implicitly account for the conversion of quarks into hadrons, a process which cannot be better described in this simple model.
The assumptions (9) and (10) are reasonable but they represent a case of โextreme freedomโ : they do not take into account the energy loss from a parent quark when it is converted to a (less energetic) final meson. This process is described, in certain situations, by the fragmentation functions. Morevover, the final mesons can have any momentum and even though higher momenta will be naturally suppressed in the calculation, we are overestimating the phase space of the reaction.
model B
Given these weak points of (9) and (10) we shall also use a second approach for the final state which is more conservative. We shall assume that the energy transferred to the heavy quarkonium $`\mathrm{\Phi }`$ will transform it into an excited (but still bound) state $`\mathrm{\Phi }^{}`$. The mass of this excited state will be taken to be slightly higher than the first charmonium and bottonium excitations $`\mathrm{\Psi }^{}`$ and $`\mathrm{{\rm Y}}^{}`$ respectively. It is known that these excitations are very weakly bound. Therefore, by choosing slightly higher masses for them, which are above the $`D`$-$`\overline{D}`$ and $`B`$-$`\overline{B}`$ decay thresholds, we are simulating a fragmentation process to a pair of nearly at rest mesons. This assumption is complementary to (9) and (10) since here we give to the heavy quarks only the โminimal freedomโ. The ground state wave function was chosen to be the Gaussian (7). Taking the harmonic oscillator as inspiration, we choose the wave function of the first excited state as a function which is odd in the $`x`$ direction (this is the direction of the chromoelectric field) and symmetric in $`x_1`$ and $`x_2`$:
$`t_B(r_1,r_2)`$ $`=`$ $`N_B{\displaystyle \frac{x_1+x_2}{2}}exp[{\displaystyle \frac{\stackrel{}{r}_1^2}{a_{}^{}{}_{}{}^{2}}}]exp[{\displaystyle \frac{\stackrel{}{r}_2^2}{a_{}^{}{}_{}{}^{2}}}]exp(i\epsilon _ft)`$ (11)
where the normalization constant is:
$$N_B^2=(\frac{4}{\pi })^3\frac{1}{a_{}^{}{}_{}{}^{8}}$$
and $`a^{}`$ is related to the size of the state $`\mathrm{\Psi }^{}`$ or $`\mathrm{{\rm Y}}^{}`$. Using the wave function (11) has some advantages. In first place, it avoids the definition of a fragmentation mechanism with the introduction of new parameters. In second place, as it can be seen, (11) is orthogonal to (9), so that the matrix element $`<\mathrm{\Psi }_f|H_{int}|\mathrm{\Psi }_i>`$ is zero if the Hamiltonian is a constant. Notice that this does not happen when we use (9) and therefore the approach A might contain spurious contributions. The same comment is valid for the calculations made in Ref. benzahra . This makes the contrast between approaches A and B even more necessary. Finally, in what follows we shall use the Hamiltonian (3), without the approximation $`x_1x_2a`$ made in model A.
Transition amplitudes and cross sections
The transition amplitude for model A can be easily computed from (5), (4), (6), (8) and (9):
$$T_{fi}=\mathrm{\Psi }_f|H_{int}|\mathrm{\Psi }_i=๐td^3\stackrel{}{r}_1d^3\stackrel{}{r}_2\mathrm{\Psi }_f^{}(\stackrel{}{r}_1,\stackrel{}{r}_2)H_{int}(\stackrel{}{r}_1,\stackrel{}{r}_2)\mathrm{\Psi }_i(\stackrel{}{r}_1,\stackrel{}{r}_2)$$
(12)
An analogous expression holds for model B with the use of (3), (4), (6), (8) and (11). We next take the amplitude squared $`|T_{fi}|^2=T_{fi}^{}T_{fi}`$ and since color is not observed, we take the average of the all initial color states and the sum of all final states:
$$|T_{fi}|^2\overline{|T_{fi}|^2}\frac{1}{3}\underset{n}{}\frac{1}{8}\underset{a}{}\underset{j}{}\underset{k}{}|T_{fi}|^2$$
(13)
The cross section with model A is given by:
$$\sigma _A=\frac{V}{(2\pi ^3)}d^3p_1\frac{V}{(2\pi ^3)}d^3p_2\mathrm{\hspace{0.17em}2}\pi _0^{\mathrm{}}๐bb\overline{|T_{fi}|^2}$$
(14)
The above expression is very simple and can be calculated almost analytically. Because of the gaussian Ansatz (4) and (6) we can easily integrate (12) over the coordinates and over the impact parameter. In the last step of (14), the integration over the phase space had to be done numerically. In fkn we made the additional assumption that the outgoing mesons are nearly at rest and we could thus simplify (10) and perform the integration over $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}_2`$ analytically. Here we prefer to be more โexactโ and perform the last integrations numerically.
The cross section with model B is simply given by:
$$\sigma _B=2\pi _0^{\mathrm{}}๐bb\overline{|T_{fi}|^2}$$
(15)
which, after the proper substitutions and integrations yields:
$$\sigma _B=\frac{32}{3}\pi ^5gE_0^2\frac{\gamma ^2}{\gamma ^21}\frac{d^{10}a^8a^{10}}{(a^2+a^2)^5[a^2a^2+d^2(a^2+a^2)]^3}exp\left(\omega ^2\frac{\frac{\gamma ^2a^2a^2}{(a^2+a^2)}+d^2}{2(\gamma ^21)}\right)$$
(16)
where:
$$\omega =\epsilon _f\epsilon _i=M_\mathrm{\Phi }^{}M_\mathrm{\Phi }$$
(17)
From the above expression we can observe that the cross section rises with the energy ($`\gamma `$) and saturates at a constant value. The enhancement of the chromoelectric field is tamed by the Lorentz contraction of the projectile. As for the size parameters, $`a`$, $`a^{}`$ and $`d`$, the cross section first rises and then falls with increasing values of the parameters. The values of the maxima strongly depend on the model and might change for a different choice of wave functions. However, the physical picture is very simple. Expression (16) tells us that the probability of converting a quarkonium of given initial size $`a`$ to a final state with size $`a^{}`$ tends to zero if $`a^{}=0`$ or if $`a^{}\mathrm{}`$ because the overlap between these very different states and the initial state is zero. For the same reason the cross section vanishes for $`a=0`$ and for $`a\mathrm{}`$. The parameter $`d`$ is associated with the extension of the capacitor. When it goes to infinity the spatial dependence of the potential disappears, it becomes a constant and then $`<\mathrm{\Psi }_f|H_{int}|\mathrm{\Psi }_i><\mathrm{\Psi }_f|\mathrm{\Psi }_i>=0`$.
### II.4 The interaction with the sources
In the introduction it was assumed that the quarkonium is well represented by a small dipole, which traverses a large capacitor. However this may be a too strong assumption because the dipole is not always so small. For example, comparing the size of the charmonium with the size of pion we have tipically $`\frac{a}{d}\frac{0.4}{0.6}0.67`$. Therefore it is necessay to include the interaction between the quark and antiquark in the quarkonium with the sources (the โplatesโ of the capacitor) which may be either a quark and an antiquark in the case of the pion or quark and a diquark in the case of the proton.
In order to take these interactions into account we shall assume that the interaction between a quark (or diquark) in the capacitor and a charm quark (or antiquark) in the dipole can be divided into a short distance and a long distance part. The later was already included before in the interaction with the chromoelectric fields produced by the sources. The former will be modelled as follows.
model C
The short distance interaction can be approximated by the contact interaction part (the one with the delta function) of the one-gluon exchange potential brac :
$$H_{int}=V_{OGE}=\underset{i=a,b}{}\underset{j=1,2}{}\frac{\alpha _s}{4}\stackrel{}{\lambda _i}\stackrel{}{\lambda _j}\left(\frac{1}{r_{ij}}\frac{2\pi }{3m_im_j}\stackrel{}{\sigma _i}\stackrel{}{\sigma _j}\delta ^3(\stackrel{}{r_{ij}})\right)$$
(18)
where $`\lambda `$ and $`\sigma `$ are the Gell-Mann and Pauli matrices respectively, which are responsible for color and spin interactions. The Coulomb term in the above expression will be neglected because it is of long range. The labels $`i=a,b`$ and $`j=1,2`$ refer to particles in the capacitor and dipole respectively. With this notation, in the interaction between particle $`a`$ and $`1`$ the delta function above takes the form:
$$\delta ^3(\stackrel{}{r_a}\stackrel{}{r_1})=\delta (x_ax_1)\times \delta (y_ay_1)\times \delta (z_az_1)$$
(19)
where $`\stackrel{}{r_1}=(x_1,y_1,z_1)`$ is the same as before and $`\stackrel{}{r_a}=(x_a,y_a,z_a)`$ is the coordinate of the particle $`a`$ in the quarkonium rest frame. In order to compute the transition amplitude we need to know the new wave functions, which now include both the quarkonium and the capacitor. They are:
$$\mathrm{\Psi }_i=f(\stackrel{}{r_1},\stackrel{}{r_2})g(\stackrel{}{r_a},\stackrel{}{r_b})c_nd_ne_mh_m$$
(20)
and
$$\mathrm{\Psi }_f=t_C(\stackrel{}{r_1},\stackrel{}{r_2})g(\stackrel{}{r_a},\stackrel{}{r_b})c_id_je_lh_k$$
(21)
In the above expression the function $`f`$ is the same as before and given by (7). The function $`t_C`$ represents the spatial distribution of the heavy quarks in the final state, which is assumed to be an excited but still bound state, very much like in model B. However, if we would choose $`t_C=t_B`$, the transition amplitude $`<\mathrm{\Psi }_f|H_{int}|\mathrm{\Psi }_i>`$ would vanish because the contact interaction does not depend on the coordinates and hence $`<\mathrm{\Psi }_f|\mathrm{\Psi }_i>`$ is the product of an odd by an even function of $`x`$, being thus zero. Since we are mostly interested in knowing the order of magnitude of this contact interaction we shall approximate the final state wave function by a gaussian, given by:
$$t_C(\stackrel{}{r_1},\stackrel{}{r_2})=N_Cexp(\frac{r_1^2}{a^{}_{}{}^{}2})exp(\frac{r_2^2}{a^{}_{}{}^{}2})e^{i\epsilon _ft}$$
(22)
with the normalization constant given by:
$$N_C^2=\left(\frac{2}{\pi }\right)^3\frac{1}{a^6}$$
(23)
The computation of the contact interaction requires the knowledge of the positions of the quarks in the capacitor, which is given by the function $`g`$
$`g(\stackrel{}{r_b},\stackrel{}{r_b})`$ $`=`$ $`N_Pexp\left[{\displaystyle \frac{(x_aX)^2}{d^2}}\right]exp\left[{\displaystyle \frac{(y_aY)^2}{d^2}}\right]`$ (24)
$`\times `$ $`exp\left[{\displaystyle \frac{(x_bX)^2}{d^2}}\right]exp\left[{\displaystyle \frac{(y_bY)^2}{d^2}}\right]`$
$`\times `$ $`exp\left[{\displaystyle \frac{\gamma ^2(z_aZ)^2}{d^2}}\right]exp\left[{\displaystyle \frac{\gamma ^2(z_bZ)^2}{d^2}}\right]`$
where $`Z=vt`$, $`d`$ and $`\gamma `$ have the same meaning as before and $`N_P`$ is the normalization constant of the projectile wave function, given by:
$$N_P^2=\frac{8\gamma ^2}{\pi ^3d^6}$$
(25)
Notice that $`g`$ is the same in the initial and in the final state. This assumption is consistent with the eikonal approximation introduced above and avoids the introduction of new parameters.
With these ingredients we can evaluate the transition amplitude:
$`T_{fi}`$ $`=`$ $`\mathrm{\Psi }_f|H_{int}|\mathrm{\Psi }_i`$ (26)
$`=`$ $`{\displaystyle ๐td^3\stackrel{}{r}_1d^3\stackrel{}{r}_2d^3\stackrel{}{r_a}d^3\stackrel{}{r_b}\mathrm{\Psi }_f^{}(\stackrel{}{r}_1,\stackrel{}{r}_2,\stackrel{}{r_a},\stackrel{}{r_b})H_{int}(\stackrel{}{r}_1,\stackrel{}{r}_2,\stackrel{}{r_a},\stackrel{}{r_b})\mathrm{\Psi }_i(\stackrel{}{r}_1,\stackrel{}{r}_2,\stackrel{}{r_a},\stackrel{}{r_b})}`$
and the cross section:
$`\sigma _C`$ $`=`$ $`{\displaystyle \frac{2^{10}}{3^4}}\pi \alpha _s^2\left({\displaystyle \frac{1}{m_am_1}}+{\displaystyle \frac{1}{m_am_2}}+{\displaystyle \frac{1}{m_bm_1}}+{\displaystyle \frac{1}{m_bm_2}}\right)^2`$ (27)
$`\times `$ $`{\displaystyle \frac{\gamma ^2}{\gamma ^21}}{\displaystyle \frac{a^6a^6}{(a^2+a^2)^5[d^2(a^2+a^2)+2a^2a^2]}}exp\left(\omega ^2{\displaystyle \frac{\frac{\gamma ^2a^2a^2}{(a^2+a^2)}+d^2}{4(\gamma ^21)}}\right)`$
where we have used (13) and the analogous expression for the sum and average over spins. Apart from a numerical factor, (16) and (27) have the same energy dependence. This is so because the same Lorentz contraction in the exponent of the Hamiltonian (3) and (4) leading to (16) is now present in the capacitor wave function (24). Moreover, the same Lorentz $`\gamma `$ factor, previously multiplying the $`\stackrel{}{E}^a`$ field in (4) reappears now in the normalization constant (25). The dependence of (27) on $`a`$ and $`a^{}`$ is qualitatively the same as the one found in (16) and has the same physical origin. Finally, the cross section above is now a monotonically decreasing function of $`d`$. The observed behavior with $`d`$ means that, in a larger capacitor the quarks are spread across a larger transverse area and it becomes more difficult for them to find the charm quarks in the target and suffer a contact interaction.
## III Results and discussion
In the numerical estimates presented below, we shall adopt $`d=0.8`$ and $`0.6`$ fm for the proton and the pion respectively. We shall also take $`a=0.4`$ and $`0.2`$ fm for the $`J/\psi `$ and $`\mathrm{{\rm Y}}`$ respectively and $`a^{}=0.8`$ and $`0.45`$ fm for the $`\mathrm{\Psi }^{}`$ and $`\mathrm{{\rm Y}}^{}`$. The bound states $`\mathrm{\Psi }`$ ($`m_\mathrm{\Psi }=3.07`$ GeV) and $`\mathrm{{\rm Y}}`$ ($`m_\mathrm{{\rm Y}}=9.46`$ GeV) will be, in model A, dissociated into pairs of mesons $`D`$ ($`m_D=1.87`$ GeV) and $`B`$ ($`m_B=5.27`$ GeV). The excited states used in models B and C have masses $`m_\mathrm{\Phi }^{}=3.8`$ GeV and $`m_\mathrm{\Phi }^{}=11`$ GeV in the case of charmonium and bottonium respectively. The value of the strong coupling constant and the constituent quark masses are the same used in brac , i.e., $`\alpha _s=0.64`$, $`m_q=0.3`$ GeV, $`m_c=1.2`$ GeV, $`m_b=4.74`$ GeV and the diquark mass is $`m_d=0.60`$ GeV.
As it is clear from (5) and (4), we need to know the average value of the color electric field in the projectile $`gE=h|gE|h`$. In a first approximation this number might be identified with the string tension $`\kappa 0.18`$ $`GeV^2`$ or $`\kappa 0.9`$ $`GeV/fm`$. The string tension calculated in adam is somewhat larger. In simonov the transverse profile of the string was studied. The strength of $`h|gE|h`$ depends on the quark-antiquark (or quark-diquark) separation, being larger for larger systems and so far it has been calculated only for large systems. Therefore $`h|gE|h`$ is another source of differences between a proton and a pion projectile. Taking an average of the values found in simonov we choose $`h|gE|h=1`$ GeV/fm.
As mentioned in the introduction, our model has common aspects with the BP approach. Therefore we shall, in what follows, compare our results for $`\sigma _{\mathrm{\Phi }h}`$ with those obtained by Kharzeev in kar2 :
$$\sigma _{\mathrm{\Phi }h}=2.5\left(1\frac{\lambda _0}{\lambda }\right)^{6.5}\text{mb}$$
(28)
with $`\lambda `$ given by
$$\lambda \frac{(sM_\mathrm{\Phi }^2)}{2M_\mathrm{\Phi }}$$
(29)
and $`\lambda _0(M_h+\epsilon )`$, where $`M_h`$ is the projectile mass and
$$\epsilon =2m_MM_\mathrm{\Phi }.$$
(30)
In Figure 2 we show the cross sections for the proton-charmonium dissociation obtained with model A (dotted lines) and model B (dashed lines) and compare them with the BP cross section (solid line with stars) given by (28). The two upper curves are obtained with $`h|gE|h=1`$ GeV/fm and the two lower curves with $`h|gE|h=0.57`$ GeV/fm (model A) and $`h|gE|h=0.53`$ GeV/fm (model B). With these smaller values of the chromoelectric field our curves come close to (28). Figure 3 shows the corresponding cross sections for the proton-bottonium dissociation. Again, the two upper curves are obtained with $`h|gE|h=1`$ GeV/fm and the two lower curves with $`h|gE|h=0.69`$ GeV/fm (model A) and $`h|gE|h=0.49`$ GeV/fm (model B). As in the previous figure, reducing the value of $`h|gE|h`$ leads to some agreement with (28). Given the conceptual resemblance between our model and the BP one, it is reassuring to find a certain similarity between the results, both in magnitude and energy behavior, once an appropriate value of $`h|gE|h`$ is chosen.
In Figure 4 we show the cross section for $`J/\psi `$ dissociation by pions compared with results obtained the meson exchange model osl (thin dotted line), with the quark exchange model wongs (thin long dashed line), with short distance QCD (the BP approach) Eq. (28) (thick solid line) and QCD sum rules nos1 (thin solid line). In spite of the fact that, at such low energies our approach looses validity, it is, nevertheless, interesting to observe that our curve is in the center of the region covered by the other calculations. In Figure 5 we compare the cross sections $`pJ/\psi `$ (upper curves) and $`\pi J/\psi `$ (lower curves) calculated with models A (dotted lines) and B (dashed lines). In the high energy limit, where both cross sections are nearly constant, we observe that the relation between the cross sections is:
$$\sigma _{p\mathrm{\Phi }}\mathrm{\hspace{0.17em}3}\sigma _{\pi \mathrm{\Phi }}\text{model A}$$
(31)
$$\sigma _{p\mathrm{\Phi }}\mathrm{\hspace{0.17em}4.2}\sigma _{\pi \mathrm{\Phi }}\text{model B}$$
(32)
which in both cases is much larger than the one expected from the additive quark model:
$$\sigma _{p\mathrm{\Phi }}\frac{3}{2}\sigma _{\pi \mathrm{\Phi }}$$
(33)
This is remarkable since the additive quark model relation holds for other high energy scattering processes like $`\pi p`$ and $`pp`$. Since $`h|gE|h`$ was kept the same for both cases, this unexpected relation between the cross sections must come from differences in the wave functions. In Figure 6 we repeat this comparison for the reactions $`p\mathrm{{\rm Y}}`$ and $`\pi \mathrm{{\rm Y}}`$, finding (31) for both models. We have kept $`h|gE|h=1`$ GeV/fm for both projectiles. Taking $`p|gE|p>\pi |gE|\pi `$ would increase the deviation from (33).
In the high energy limit ordinary hadrons are expected to have a geometrical total cross section. Since the quarkonium dissociation discussed here is a more specific reaction it is not obvious that its cross section follows a geometrical behavior. Such a behavior was found in arleo : $`\sigma _{\mathrm{\Phi }h}\alpha _sa_0^2`$ where $`a_0`$ is the Bohr radius of the quarkonium. In our case, as it can be seen from (16), (27) and from the numerical evaluation of (14), we have a very non-trivial dependence on $`a`$. Since the initial state (containing the variable $`a`$) is the same, the difference between models comes from the spatial dependence of the final state. The plane waves in model A have no spatial scale. Therefore they are more โinclusiveโ and so $`\sigma _A`$ should be closer to the quarkonium-hadron total cross section than $`\sigma _B`$. In model B the quarkonium ground state is converted into a resonance-like state, which wave function contains the size parameter of the resonance and distorts the final geometrical behavior. Therefore model A is closer to a geometrical behavior than model B.
In order to see how far we are from the geometrical behavior, we show in Figures 7 and 8 the dependence of $`\sigma _A`$ (dotted line) and $`\sigma _B`$ (dashed line) on $`a`$ for charmonium (Fig. 7) and bottonium (Fig. 8) dissociation. The cross sections are divided by $`a^2`$ so that geometrical behavior translates into a horizontal line. We see that, whereas model A tends to this behavior, model B is far from a geometrical behavior. This indicates again that our model is very sensitive to the choice of the final state wave functions.
In Figure 9 we show the cross section $`\sigma _C`$ (27) for charmonium dissociation by protons (solid line) and by pions (dashed line). In Figure 10 we show the same quantitity for bottonium dissociation. We use the central values for $`a`$, $`a^{}`$, $`d`$ and $`\alpha _s`$. We can see that, in all processes, the cross sections are more than two orders of magnitude smaller than the corresponding cross sections computed with model A or model B. No possible change in parameters could make these cross sections comparable. Another feature of these curves is that the cross sections for $`J/\psi `$ dissociation by pions are larger than those for protons by a factor close to $`4`$. This might be guessed looking at (27). The pion is light quark-antiquark system and the proton is light quark-diquark dipole. The diquark is twice heavier than a constituent quark. Whereas for the pion $`m_b=m_a`$, for the proton we have $`m_b=2m_a`$.
Before concluding we would like to make a remark concerning medium effects on the cross sections calculated above. We are primarily studying reactions which happen before thermalization (in nucleus-nucleus collisions) or with no thermalization at all (in proton-nucleus collisions). The formation time of the heavy quark pair is of the order of $`0.2`$ fm. The thermalization time of hadronic matter formed in heavy ion collisions is a model dependent quantity. Early estimates pointed to $`1`$ fm. Recent estimates heinz point to $`0.6`$ fm. Even taking seriously this last number, it is fair to say that heavy quark pair production (and collision with a hadron at high energies) precedes the formation of an equilibrated medium. After thermalization, the energy is completely redistributed and collisions occur at energies of the order of the temperature ($`<1`$ GeV). In this regime we do not expect our approach to be valid. The effects of a thermal medium on the heavy quarkonium are known wong : the string tension becomes weaker, the quarkonium size increases and its mass decreases. These effects are, all of them, very small except if we get close to the deconfinement transition temperature. In view of these considerations we have neglected medium effects in our calculations.
## IV Conclusions
We have developed a simple model for the non-perturbative quarkonium-hadron interaction. At the present stage of the field, this sort of model is still useful to organize the ideas. We tried to make simple and yet realistic choices for the interaction Hamiltonian and for the wave functions. In particular we have treated the final state in two very different and complementary ways. Simple models are not appropriate to provide very precise results but they can help in determining the order of magnitude of the cross sections and their behavior with the reaction energy. Having said that, we can summarize our conclusions as follows.
i) The charmonium-hadron cross section is of a few milibarns. The bottonium-hadron cross section is about four times smaller. This is in agreement with most of the previous calculations.
ii) All cross sections grow with the reaction energy and reach a plateau in the high energy limit. This is in agreement with the BP approach.
iii) In this limit they do not obey the simple relations derived from the additive quark model.
iv) Also in this limit our cross sections deviate significantly from the geometrical behavior ($`\sigma a^2`$).
v) The contact interactions between the heavy quarks and the light quarks in the light hadrons is negligible compared to the long-distance quark - $`\stackrel{}{E}^a`$ field interaction. This is surprising since sometimes the dipole and the capacitor have similar sizes. This finding gives a posteriori support to our model and also to the BP approach.
Conclusion i) may be relevant for RHIC and LHC physics. Conclusions ii, iii and iv suggest that the heavy quarkonium has interaction properties which are very different from light hadrons. This has been conjectured before. In particular, in dnw this difference was attributed to the fact that in heavy quarkonia the energy is mainly stored in the masses whereas in light hadrons the energy (mass of the hadron) comes mostly from the gluonic fields.
Acknowledgements: This work has been supported by CNPq and FAPESP. We are indebted to M. Nielsen, D.A. Fogaรงa and F. Durรฃes for fruitful discussions. |
warning/0506/hep-th0506230.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Recently, a quantum foam picture of topological string theory has been discovered . According to this duality, A-model topological string amplitudes on $`^3`$ and on more general toric Calabi-Yau manifolds can be computed by a statistical model of a melting crystal. The crystal is a physical picture of the A-model target space Kahler gravity, and as a quantum foam description, it captures the geometry up to very short distances. The mathematical side of the correspondence is the Donaldson-Thomas theory reformulation of Gromov-Witten invariants .
In particular the partition function of the melting crystal computes closed string A-model amplitudes. This was explicitly verified for $`^3`$ , and for more general non-compact geometries . This duality can be further extended by introducing non-compact brane probes in the geometry. Such brane probes were found to correspond to defects in the $`^3`$ crystal .
It is also interesting to study brane probes in the crystal model of conifold geometry. There are two ways to build a crystal for the conifold: in the first way one glues together two pieces of $`^3`$ geometries, in the second method of one makes use of the slicing suggested by the open string description. It is the latter construction we use in this paper. Here the crystal is like the $`^3`$ model, but ending in a wall in one direction.
The conifold crystal is particularly interesting because it is a clear example of open-closed duality. As is well-known, the closed topological A-model on the resolved conifold is dual to Chern-Simons theory on $`S^3`$ . Further, the natural observables of Chern-Simons theory are Wilson loop operators, related to knot and link invariants in the 3-manifold $`S^3`$. As is described in , to each knot intersecting the $`S^3`$ we can associate a Lagrangian cycle, over which probe branes can be wrapped. Adding a Wilson loop observable along a knot in the Chern-Simons side thus corresponds to inserting non-compact Lagrangian brane probes in the closed string geometry.
It is then an interesting question how various Chern-Simons knot and link invariants are encoded in the crystal model of conifold. The crystal model is a simple statistical model of an infinite crystal with a wall in one direction. Non-compact branes correspond to fermionic operators in the transfer matrix formulation of the crystal, in agreement with the general picture that non-compact D-branes in the topological B-model are fermions. Operations in the crystal, such as the computation of amplitudes of non-compact branes are therefore easy. These amplitudes are then natural generating functions of certain knot invariants.
In this paper we investigate the crystal picture of non-compact brane insertions as generating knot expansions. Earlier work in this direction includes , where it was computed that insertion of a single brane corresponds to the unknot invariant in Chern-Simons theory, and related observations about the connection of the crystal and topological vertex formalism in .
In section 2 we analyze the $`^3`$ non-compact brane amplitudes. As $`^3`$ is the limit of the conifold when the Kahler parameter is sent to infinity, these amplitudes are generating functions of the leading part for certain knot invariants. In particular, inserting two branes, one on each leg of the crystal generates the leading part of Hopf link invariants. The more general case of many brane insertions generates the leading part of Hopf link invariants in arbitrary representations. In fact, we find that the many brane amplitudes can be viewed alternatively as generating Hopf link tensor product representations, corresponding to Young diagrams with a single row. This latter point of view relates the crystal expansion to the topological vertex formulation of the brane amplitudes explicitly.
In section 3 we consider the conifold model of the crystal. In particular, we discuss how to generate the full unknot invariant in the conifold, and derive the Ooguri-Vafa generating function. In section 4 we introduce a crystal with two walls and compute the partition function of a single brane insertion.
These non-compact brane amplitudes can also be derived in the topological vertex formulation. We compute and compare the same brane amplitudes in the A-model vertex formulation, where they are naturally expressed as knot expansions. We verify the crystal and vertex results agree. While the crystal framework is schematically simple to use, the summation of vertex amplitudes in many cases is complicated. The crystal then gives a simple and natural closed expression for the vertex results. The comparison of crystal amplitudes with A-model topological vertex results is discussed in section 5. In section 6 we compare one nontrivial crystal amplitude with B-model topological vertex, also finding agreement.
Finally, section 7 contains a summary and discussion., where we consider the connection of the crystal brane amplitudes (open Donaldson-Thomas invariants), Chern-Simons invariants and Gopakumar-Vafa invariants. In particular we conjecture that free energy associated to the crystal amplitudes can be simply expressed as a Gopakumar-Vafa expansion. Thus D-brane degeneracies are simply encoded in the crystal free energy.
## 2 Knot invariants from the crystal
The Calabi-Yau crystal is defined by a statistical sum over three dimensional partitions , where partitions are weighted by $`q^{\mathrm{\#}boxes}`$, and $`q=e^{g_s}`$.
We will first consider the geometry $`^3`$. In this case the crystal is understood as filling the positive octant of $`^3`$, which is a toric base of $`^3`$. One way to imagine the 3d crystal is to build from diagonal slices of two dimensional partitions. To assemble to 3d partitions, the diagonal slices have to satisfy the interlacing condition . A simple way to compute the crystal partition function is the transfer matrix formalism of . In this formalism we assign a fermionic Fock space to each two dimensional diagonal slice. To construct the crystal in operator language, we use bosonization of the chiral fermion $`\psi (z)=:e^{\varphi (z)}:`$, and the creation/annihilation part of the bosonic vertex operator, $`\mathrm{\Gamma }_\pm (z)`$. In this way the crystal partition function is built as
$$Z(q)=\underset{3dpartitions}{}q^{\mathrm{\#}boxes}=0|\underset{m=1}{\overset{\mathrm{}}{}}\mathrm{\Gamma }_+(q^{m1/2})\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{\Gamma }_{}(q^{n+\frac{1}{2}})|0$$
From the commutation relations
$$\mathrm{\Gamma }_+(z)\mathrm{\Gamma }_{}(z^{})=(1z/z^{})^1\mathrm{\Gamma }_{}(z^{})\mathrm{\Gamma }_+(z)$$
it is straightforward to see that the partition function is the McMahon function $`M(q)`$
$$Z(q)=M(q)=\underset{n=1}{\overset{\mathrm{}}{}}(1q^n)^n.$$
In the B-model picture non-compact Lagrangian probe branes can be thought of as fermions inserted in the geometry. These fermions are virtually free chiral fermions except they transform between different patches with Fourier transformation . The corresponding crystal description of probe branes are fermionic operators <sup>1</sup><sup>1</sup>1For the crystal we use p=0 framing.
$$\mathrm{\Psi }_D(z)=\mathrm{\Gamma }_{}^1(z)\mathrm{\Gamma }_+(z).$$
Similarly, anti-branes are represented by
$$\mathrm{\Psi }_{\overline{D}}(z)=\mathrm{\Gamma }_{}(z)\mathrm{\Gamma }_+^1(z).$$
A Lagrangian probe brane with geometry $`S^1\times ^2`$ (with crystal axis $`(x,y,z)`$)
$$y=x+u=z+u,u=g_s(N+1/2)>0$$
ending on the $`y`$ axis at distance $`u`$ is described inserting a fermionic operator $`\mathrm{\Psi }_{D,y}(e^u)`$ at the slice $`t=N+1`$. Similarly, a brane at distance $`v`$ on the $`x`$-axis is described by $`\mathrm{\Psi }_{D,x}(e^v)`$ at the negative side of the diagonal $`t=(N+1)`$ (Fig. 1).
Inserting $`m`$ Lagrangian branes at distances $`g_s(N_i+\frac{1}{2})`$ ($`i=1\mathrm{}m`$) on the $`y`$ axis, and $`n`$ Lagrangian anti-branes<sup>2</sup><sup>2</sup>2We could of course have inserted branes, which would cause a change of framing difference in the end result. on the x-axis at distances $`g_s(M_j+\frac{1}{2})`$, ($`j=1\mathrm{}n`$), as first derived in , gives
$$Z(a_1,\mathrm{}a_n;b_1\mathrm{}b_m;q)=\mathrm{\Psi }_{\overline{D},x}(b_1)\mathrm{}\mathrm{\Psi }_{\overline{D},x}(b_n)\mathrm{\Psi }_{D,y}(a_1)\mathrm{}\mathrm{\Psi }_{D,y}(a_m)$$
$$=M(q)\left(\underset{i=1}{\overset{m}{}}\underset{j=1}{\overset{n}{}}\frac{L(a_i,q)L(b_j,q)}{(1a_ib_j)}\right)\left(\underset{i>j}{}(1a_i/a_j)(1b_i/b_j)\right).$$
(1)
Here $`L(a_i,q)`$ for $`a_i=q^{N_i+\frac{1}{2}}`$ denotes the quantum dilogarithm
$$L(a_i,q)=\underset{i=1}{\overset{\mathrm{}}{}}(1q^{n+N_i})=\underset{n}{}a_i^nh_n(q^\rho ),$$
(2)
which can also be expressed in terms of the complete symmetric polynomials (defined in Appendix A). The quantum dilogarithm is the brane wavefunction, as also can be seen from direct disk amplitude computation , as well as from the insertion of a fermionic operator to the B-model geometry corresponding to the limit shape of the crystal . The additional factors of type $`(1a_i/a_j)`$ and $`(1a_ib_j)`$ correspond to stretched strings between branes. We will explicitly see later how these arise in the A-model topological vertex picture.
In the following we will re-interpret this expression as a generating function of certain knot invariants in arbitrary representations. As $`^3`$ can be thought of as a limit of the conifold when its Kahler parameter $`t=g_sN\mathrm{}`$, by geometric transition we expect to see the leading part of knot invariants. We are then probing the invariants of $`U(\mathrm{})`$ Chern-Simons theory. More precisely, from the geometric picture of the Lagrangian branes with topology $`S^1\times ^2`$ we expect to find unknot and Hopf link invariants. As the crystal result is written entirely in terms of dilogarithms and simple prefactors from the stretched strings, it is not immediately obvious that these would provide the generating functions for more complicated link invariants, for example for Hopf link invariants in tensor product representations. It will turn out that the simplicity of crystal results is partly due to a particularly natural framing choice.
### 2.1 Single unknot
Consider first a single brane on the $`y`$ axis. In this case, we have
$$Z(a,q)=M(q)L(a,q).$$
Normalized by $`M(q)`$, it is indeed the leading part of the generating function for unknot invariants as computed in Chern-Simons theory after the geometric transition . It is also explicitly seen as a generating function by expanding
$`L(a,q)`$ $`=`$ $`e^{_{n=1}^{\mathrm{}}\frac{a^n}{n[n]}}=1+{\displaystyle \frac{a}{(q^{\frac{1}{2}}q^{\frac{1}{2}})}}+a^2{\displaystyle \frac{q^2}{(q^21)(q1)}}+\mathrm{}`$ (3)
$`=`$ $`{\displaystyle \underset{Ronerow}{}}W_Ra^{|R|}`$
Here the notation is $`[n]=q^{n/2}q^{n/2}`$, and the sum is rewritten as a sum over representations $`R`$. $`W_R`$ are the unknot invariants in zero framing, and the summation runs over one row representations only, so that $`1=\mathrm{}`$, $`2=\mathrm{}\mathrm{}`$, etc. So a single brane inserted in the $`C^3`$ crystal computes the generating function of the leading part of unknot invariants, for one row representations.
### 2.2 Hopf link
Inserting an antibrane on $`x`$-axis and a brane on $`y`$-axis gives the generating function of Hopf link invariants for single row representations. Expanding the normalized part of the partition function
$$\stackrel{~}{Z}(a,b,q)=\frac{Z(a,b,q)}{M(q)}=\frac{L(a,q)L(b,q)}{(1ab)}$$
(4)
gives
$$\stackrel{~}{Z}(a,b,q)=\underset{R,Ponerow}{}q^{\frac{\kappa _R+\kappa _P}{2}}W_{R^tP^t}a^{|R|}b^{|P|}$$
(5)
where $`\kappa _R=|R|+_iR_i(R_i2i)`$, for a general representation. In the summation we only have one row representations. We will later prove this expansion by comparing with the topological vertex, and the q-dependent prefactors will be seen as vertex framing factors $`(1,0)`$ : i.e. a brane at framing $`1`$ and an antibrane at framing $`0`$. Alternatively, when expressed in terms of $`q^1`$ this expansion gives Hopf Link coefficients with knot framing $`(1,1)`$.
### 2.3 Hopf link with many rows
In the general case, for n branes on the y-axis and m antibranes on the x-axis the normalized partition function (1) generates the leading part of Hopf link coefficients with $`(n,m)`$ rows
$`\stackrel{~}{Z}(a_1,\mathrm{}a_n;b_1,\mathrm{}b_n;q)=`$
$`{\displaystyle \underset{R_1,\mathrm{}R_n}{}}{\displaystyle \underset{P_1,\mathrm{}P_m}{}}q^{\frac{\kappa _R+\kappa _P}{2}}W_{P^tR^t}a_1^{|R_1|}\mathrm{}a_n^{|R_n|}b_1^{|P_1|}\mathrm{}b_m^{|P_m|}`$ (6)
where $`R=(R_n,\mathrm{},R_1)`$ and $`P=(P_m,\mathrm{}P_1)`$ are $`n`$ and $`m`$ row representations respectively. The last summation also contains โimproperโ Young- diagrams. For a proper Young diagram $`(R_n,\mathrm{},R_1)`$, we must have $`R_1R_2\mathrm{}R_n`$. Our summation contains also a finite number of terms where this condition is not satisfied, but these improper contributions can still be formally written using the definitions of Schur functions and Casimir $`\kappa _R`$. Appendix B contains a partial proof of this formula (done for a simplified case) as well as details on the improper contributions.
While here it appears that the crystal generates Hopf link invariants in arbitrary representations, when comparing with the topological vertex, we will find the same crystal partition function with a number of branes inserted on each leg can be viewed as generating more complicated link invariants, corresponding to the tensor product representations of Hopf link in one-row representations. We will return to this point in section 5, where the crystal partition function as a knot generating function will be re-examined.
## 3 Conifold crystal
In the following we will examine how to obtain knot invariants from the crystal model of resolved conifold $`O(1)O(1)^1`$. The crystal melting model describing topological A-model on this geometry was obtained using the large $`N`$ dual Chern-Simons theory in . The geometry of the crystal reflects the toric diagram of the resolved conifold, and it is obtained by inserting a wall in one direction. We will insert the wall at the positive slice $`N`$, constructing the conifold geometry with Kahler parameter $`t=g_sN`$, which we often refer to as $`Q=e^t=q^N`$ (Fig. 2).
The partition function of this crystal model is thus obtained as
$$Z^{P^1}(q,N)=0|\underset{m=1}{\overset{\mathrm{}}{}}\mathrm{\Gamma }_+(q^{m1/2})\underset{n=1}{\overset{N}{}}\mathrm{\Gamma }_{}(q^{(n1/2)})|0=M(q)e^{_k\frac{Q^k}{k[k]^2}}$$
(7)
in agreement with the topological vertex result (68) . Taking the Kahler parameter $`t\mathrm{}`$ gives back the partition function of the $`^3`$ crystal. Non-compact Lagrangian branes in the crystal are again defects described by fermionic operators.
### 3.1 Full unknot invariant
Unknot invariants with many row representation can be generated by inserting a number of branes on the non-compact leg of the conifold crystal. This is analogous to the topological vertex picture as will be seen in section 5. Since now we have the full conifold geometry, we will get the full unknot invariants, unlike in the $`^3`$ geometry which could only see the leading part of knot invariants (with $`t\mathrm{}`$).
Including $`m`$ antibranes at positions $`a_i=q^{N_i+1/2}`$, $`i=1\mathrm{}m`$, at the non-compact leg their normalized partition function can be written as <sup>3</sup><sup>3</sup>3Here normalization is with the conifold partition function and additional $`\xi (q)=_{i=1}^{\mathrm{}}1/(1q^i)`$ factors which has to be dropped in comparison with topological string amplitudes.
$$\stackrel{~}{Z}_D^{P^1}(a_1,\mathrm{}a_n)=\left[\underset{i<j}{\overset{m}{}}(1\frac{a_i}{a_j})\right]\underset{i=1}{\overset{m}{}}\frac{L(a_i,q)}{L(a_iQ,q)}.$$
(8)
Taking a single brane first at $`a=q^{N_1+1/2}`$ gives the full unknot generating function for single row representations by the rearrangement
$`\stackrel{~}{Z}_D^{P^1}(a)`$ $`=`$ $`{\displaystyle \frac{L(a,q)}{L(aQ,q)}}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}a^n\left({\displaystyle \underset{i=0}{\overset{n}{}}}h_i(q^\rho )h_{ni}(Qq^\rho )\right)=`$ (9)
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}a^nh_n(q^\rho ,Qq^\rho ).`$
In the first equality the expression of dilogarithm in terms of symmetric polynomials is used given (2), in the second equality (48) was used. The final coefficient $`h_n(q^\rho ,Qq^\rho )`$ is precisely the full unknot invariant (quantum dimension) for one-row representation, $`R=(n,0\mathrm{}0)`$. Taking $`m`$ branes and expanding in their positions $`(a_1,\mathrm{}a_m)`$ gives similarly unknot invariants with $`m`$-row representation. The proof of this is completely analogous to the induction included in Appendix B.
We note that the full unknot invariants were extracted before in , following a different prescription based on branes inserted in the compact leg of $`^1`$. Our procedure is different and is motivated by the topological vertex picture as will be discussed in more detail below.
### 3.2 Ooguri-Vafa generating function
Chern-Simons theory on $`S^3`$ is the large $`N`$-dual to closed topological string theory on the resolved conifold. The duality can be seen as a geometric transition \- wrapping a large number of branes on the base $`S^3`$ of deformed conifold, in the large $`N`$ limit the geometry transits to the resolved conifold without branes.
The geometry can be probed by non-compact branes . Wrapping the probe branes on a Lagrangian cycle, intersecting the $`S^3`$ in a given knot, the worldvolume theory on the probe branes will also be a Chern-Simons theory. In addition, there are open string stretched between the probe branes and the original large number of branes wrapping the $`S^3`$ and making the geometric transition. Integrating out these degrees of freedom gives an effective theory on the probes branes, which is Chern-Simons theory with additional corrections - the Ooguri-Vafa generating function. For a single unknot it is given as
$$Z_{\mathrm{OV}}=\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{(e^{\frac{nt}{2}}e^{\frac{nt}{2}})}{n[n]}a_{\mathrm{OV}}^n\right]$$
where $`[n]=q^{n/2}q^{n/2}`$ as before, and $`a_{OV}`$ is the parameter of the one-dimensional holonomy matrix, that is the integral of holonomy of the Chern-Simons gauge field around the circle loop (corresponding to the unknot) intersecting the $`S^3`$. After analytic continuation<sup>4</sup><sup>4</sup>4Upper index $`a`$ denotes analytic.,
$$Z_{\mathrm{OV}}^a=\mathrm{exp}\left[\frac{(a_{\mathrm{OV}}^n+a_{\mathrm{OV}}^n)}{n[n]}e^{nt/2}\right]$$
(10)
the Ooguri-Vafa generating function agrees with the topological string amplitude of a probe brane inserted in the conifold geometry. The topological amplitude can also be derived considering the relevant open topological string amplitude from the M-theory point of view of . Alternatively, it can be computed in the topological vertex formulation. We will consider the latter computation in section 5.
Here we show that inserting a brane in the compact leg of the conifold crystal reproduces the Ooguri-Vafa generating function. Inserting an antibrane<sup>5</sup><sup>5</sup>5An antibrane is chosen for convenience here. When matching the crystal to the topological vertex result, we will choose the convention $`q_{vertex}=1/q_{crystal}`$, which turns an antibrane in the crystal to a brane in the vertex. on the compact leg of the crystal at the positive slice at $`a=q^{N_0+1/2}`$ we obtain
$`Z_{D,y}^{P^1}(q,N_0,N)=\xi (q)M(q)e^{_{n>0}\frac{q^{n(N+1)}}{n[n]^2}}e^{_{n>0}\frac{q^{n(N_0+1/2)}+q^{n(NN_0+1/2)}}{n[n]}}`$ (11)
$`=`$ $`\xi (q)Z^{P^1}L(a,q)L(Q/a,q),`$
where $`Z^{P^1}`$ is given in (7), and now the Kahler parameter gets shifted to $`t=g_s(N+1)`$ due to brane insertion, so that $`Q=q^{N+1}`$. This is indeed the Ooguri-Vafa generating function with the identifications
$$a_{\mathrm{OV}}=q^{N_{\mathrm{OV}}+\frac{1}{2}}N_{\mathrm{OV}}=N_0\frac{N}{2}$$
i.e. the position of brane in the geometry is measured from the middle point of the compact leg. We note that the crystal provides a straightforward way to compute this result.
Inserting more branes on the compact leg would correspond to inserting more stacks of branes in the geometry. The generating function can be easily computed on the crystal side. On the other hand, in the crystal geometry it is not clear how to incorporate increasing the number of branes in a single stack (thus increasing the holonomy matrix of probe).
It is a natural question to ask if inserting a number of branes on each leg of the conifold crystal would provide complete Hopf link invariants with many rows, similarly to the leading part of Hopf link invariants obtained from $`^3`$. For example, inserting a brane on the compact leg at position $`a`$ and an antibrane on the non-compact leg at position $`b`$ in the conifold crystal one gets
$$Z_{D;\overline{D}}^{P^1}(q,N)=Z^{P^1}\frac{L(a,q)L(b,q)L(Q/a,q)}{L(bQ,q)(1ab)},$$
where again the Kahler parameter gets shifted to $`t=g_s(N+1)`$ due to brane insertion on compact leg. Expanding in $`a`$ and $`b`$ does not naturally give many row Hopf link invariants. The reason is seen better in the language of topological vertex, where Hopf link invariants are associated to having two branes inserted, each on a non-compact leg of the conifold geometry . In the conifold crystal model a Hopf link would naturally arise from placing branes on the non-compact $`x`$-axis and another on the non-compact $`z`$-axis. In the diagonal slicing of the crystal we work in, the latter branes are not natural to insert. Working out the operators for insertion of such branes, and generating full Hopf link invariants from the crystal is left for future work.
## 4 Calabi-Yau crystal with two walls
The local conifold model for the crystal of can be easily generalized to represent the geometry with two neighbouring $`^1`$. This is naturally described by a crystal with two walls, on both the positive and negative slice, at distance $`t_1=N_1g_s`$ and and $`t_2=N_2g_s`$ respectively, which are the two Kahler parameters of the geometry (Fig. 3).
The partition function is computed as
$`Z^{2walls}(q,N_1,N_2)`$ $`=`$ $`0|{\displaystyle \underset{n=1}{\overset{N_1}{}}}\mathrm{\Gamma }_+(q^{n1/2}){\displaystyle \underset{m=1}{\overset{N_2}{}}}\mathrm{\Gamma }_{}(q^{(m1/2)})|0=`$ (12)
$`=`$ $`\mathrm{exp}{\displaystyle \underset{k>0}{}}{\displaystyle \frac{(1q^{kN_1})(1q^{kN_2})}{k[k]^2}},`$
The factors in the exponent represent (apart from the unity giving McMahon function) worldsheets wrapping each of the spheres independently, and then both of them simultaneously.
Let us now insert a brane on the right compact leg at a position given by $`a=q^{N_0+1/2}`$, which gives
$$Z_{D,y}^{2walls}=Z^{2walls}\xi (q)L(Q_1/a,q)\frac{L(a,q)}{L(aQ_2,q)},$$
(13)
where again there is a shift of the Kahler parameter corresponding to the leg the brane is put on. It is also interesting to compare this result with a brane in the resolved conifold (11). The essential difference is the dilogarithm in the denominator, which represents worldsheet wrapping a part of right $`^1`$ (of length $`N_0`$) and the whole left $`^1`$ (of length $`N_2`$).
It would be very interesting to generalize this construction by gluing together pieces of crystals to get an arbitrary toric geometry.
## 5 Comparison with the topological vertex
In this section we show that the amplitudes computed by crystal models with one or two walls, and multiple brane insertions, are indeed consistent with the topological string results. We will perform the topological vertex calculations and appropriately match vertex and crystal moduli, and find a perfect dictionary between these two points of view. Most vertex calculations are performed in A-model language , but we also provide one nontrivial example of a B-model amplitude .
Thus, let us focus on the A-model topological vertex. In this formulation the target space is a non-compact toric Calabi-Yau 3-fold, and topological amplitudes can be computed from a planar โFeynman diagramโ which encodes the geometry of the 3-fold. Each edge of such a diagram corresponds to a shrinking cycle of a toric fibration, and compact intervals represent local $`^1`$โs in Calabi-Yau geometry. The A-model vertex is a trivalent vertex for such a diagram and it encodes the structure of topological string in a single $`^3`$ patch. The full toric 3-fold can be built from $`^3`$ patches, and the amplitudes can be found by gluing vertexes according to the relevant gluing rules. The gluing process is implemented by a careful analysis of open strings ending on stacks of Lagrangian branes put on the axes. The vertex is most conveniently expressed in a representation basis as $`C_{R_1R_2R_3}`$, with each representation corresponding to a stack of branes on a single axis of $`^3`$ patch.
Apart from gluing string amplitudes, the vertex allows also to compute open string amplitudes in presence of a particular class of special Lagrangian branes of topology $`\times ๐^1`$. The projection of these branes onto the plane of a toric diagram is a semi-infinite line with its endpoint attached to one edge of the diagram. For example, the partition function for inserting 3 non-compact branes on each leg of $`^3`$ is computed as
$$Z(V_1,V_2,V_3)=\underset{R_1,R_2,R_3}{}C_{R_1,R_2,R_3}\mathrm{Tr}_{R_1}V_1\mathrm{Tr}_{R_2}V_2\mathrm{Tr}_{R_3}V_3,$$
where $`V_i`$ are sources (holonomy matrices) corresponding to inserted branes, and in general can be given by infinite matrices. This amplitude is written in the so-called canonical framing. In general, the vertex exhibits a framing ambiguity, which is a statement that one needs to specify one integer for each stack of branes to fully determine the amplitude. This is intimately connected with framing ambiguity in knot theory, and can be traced by a derivation of the vertex from Chern-Simons theory. All necessary details about computational framework for A-model, including framing ambiguity and other subtleties, are given in appendix C.
To match crystal and vertex results, a few important issues have to be taken into account. Firstly, the topological vertex is normalized in such a way that the McMahon function $`M(q)`$ of $`^3`$ does not arise from calculations. Secondly, we need to choose some particular framing; usually this is $`(1)`$ framing on one leg and the canonical one for other legs. Then, we have to take holonomy matrices $`V_i`$ to be one dimensional
$$V_i=a_i=q^{N_i+1/2}$$
so in this sense the crystal can see only a fraction of what the full vertex computes. On the other hand, the crystal calculations are much simpler, so this is quite an advantage of using it.
If we have a single brane on one leg, then the above $`a_i`$ become simply moduli seen in the crystal. For more branes on one leg, we have to introduce parameters which give their positions, which must be combined with holonomy matrices appropriately. We will see examples of this in what follows.
Finally, we perform substitution
$$q\frac{1}{q}=q_{crystal}$$
(14)
in vertex result to map crystal-branes to vertex-branes. In fact, in topological string such an operation exchanges branes to antibranes , and what we call branes and antibranes can be regarded just as a convention. Not performing the $`q`$ inversion would result in mapping crystal branes to vertex antibranes. We choose the former point of view. In fact, the $`q`$ inversion is important only for configurations with branes; the partition functions without any branes is invariant under $`q1/q`$.
Let us note, that while some of the topological vertex amplitudes we consider here were already written in the literature; it is not at all obvious that these amplitudes given by topological vertex rules in terms of sums over representations can be resummed into compact expressions, involving just dilogarithms and simple polynomials (as we have seen from crystal point of view). This fact was also noticed in . Nonetheless, with the proper vertex framing chosen we rederive all these crystal results (which are in crystal canonical framing).
By construction, the topological vertex includes the correct worldsheet instantons which can appear in any toric construction, with or without Lagrangian probe branes. The contributions from specific instantons which stretch between probe branes can be read off from the form of the free energy. Specifically, the $`\mathrm{Li}_1`$ function in the factor
$$(1ab)=\mathrm{exp}\left(\mathrm{Li}_1(ab)\right),$$
(15)
appearing in all calculations involving more than one probe brane, shows that this is a contribution from annuli instantons and not of any higher genera instantons.
### 5.1 Resolved conifold results
We shall first test the conifold crystal results, the $`^3`$ crystal results will naturally follow then from taking the Kahler parameter to infinity. The resolved conifold partition function $`Z^{P^1}`$ (68) has been computed before in several places , and it is in agreement with the crystal result.
Let us then compute brane configurations corresponding to those found in the crystal language. We start with a single brane on the external leg of the conifold in canonical framing
$$Z_{Dext}^{P^1}(V)=\underset{P,R}{}C_{R^tP}(Q)^{|R|}C_RTr_PV.$$
(16)
Using identities on Schur functions we get (see also )
$$Z_{Dext}^{P^1}(V)=Z^{P^1}\underset{P}{}s_P(Qq^\rho ,q^\rho )Tr_PV.$$
(17)
Taking the matrix $`V`$ to be one dimensional $`V=a=q^{N_0+1/2}`$, and using $`Tr_R(a)=s_R(a)`$ and formula (Appendix A), we obtain
$$Z_{Dext}^{P^1}=Z^{P^1}\underset{R}{}s_R(Qq^\rho ,q^\rho )s_R(a)=Z^{P^1}\frac{L(aQ,q_{crystal})}{L(a,q_{crystal})}.$$
(18)
It is important to note, that the sum is in fact performed over representations corresponding to diagrams with only one row (for a single number $`a`$ and for any representation given by a diagram with more than one row $`s_{rep.with>1rows}(a)=0`$). Taking into account the mapping (14) we find perfect agreement with the normalized crystal result for antibranes (8).
A single brane can also be situated on the compact leg of the conifold at position $`g_sD`$
$$Z_{Dint}^{P^1}=\underset{R,Q^L,Q^R}{}C_{RQ^L}(1)^sq^fe^LC_{R^tQ^R}Tr_{Q^L}VTr_{Q^R}V^1.$$
It is possible to perform resummation for $`V=a=q^{N_0+1/2}`$ and if $`(1)`$ framing is chosen. If we follow the crystal convention and set the size of the compact leg to be $`N+1`$ (the shift is responsible for brane insertion), and absorb the brane position into its modulus by defining
$$N_0^{}=D+N_0,a^{}=q^{N_0^{}+\frac{1}{2}},$$
(19)
we get after substitution (14)
$$Z_{D,y}^{P^1}=Z^{P^1}(N+1)L(a^{},q_{crystal})L(Q/a^{},q_{crystal}),$$
which is the same result as (11).
It is also possible to insert several branes on the external or internal leg. For example, for $`M`$ branes on the compact leg in $`(1)`$ framing we take $`V_i=a_i=q^{N_i+1/2}`$ and then get analogous factors as above. The Kahler parameter gets modified to $`N+M`$, and brane positions $`D_i`$ get absorbed similarly as above into $`N_i^{}`$ and modified moduli $`a_i^{}`$. We also have to take the stretched strings between the branes into account (see (72)). All these factors combine to
$$Z_{Mbranes}^{P^1}=Z^{P^1}(N^{})\left[\underset{i<j}{}(1\frac{a_i^{}}{a_j^{}})\right]\left[\underset{i=1}{\overset{M}{}}L(a_i^{},q_{crystal})L(Q/a_i^{},q_{crystal})\right],$$
(20)
which is the same as the crystal result (4.42) in .
#### 5.1.1 Brane and antibrane on two legs
Let us put one brane on the compact leg of the resolved conifold at distance $`D`$ with holonomy matrix $`V_1`$, and the second brane on non-compact leg with holonomy $`V_2`$ (Fig. 4).
We also take one-dimensional holonomy matrices $`V_i=q^{N_i+1/2}`$, and absorb the position on the compact leg into $`V_1`$
$$a=q^{D+N_1+1/2}=q^{N_1^{}+1/2},b=q^{N_2+1/2}.$$
(21)
The partition function in $`(1,0)`$ framing is
$`Z_{Dy;\overline{D}x}^{P^1}`$ $`=`$ $`{\displaystyle C_{RQ_L,P^t,}(1)^{|P|}(Q)^{|R|}C_{R^tQ_R}s_{Q_L}(a)s_{Q_R}(Q/a)s_P(b)}`$
$`\left[(1)^{|Q_LR|+|Q_RR^t|}q^{\frac{\kappa _{Q_LR}+\kappa _{Q_RR^t}}{2}}\right]=`$
$`=`$ $`L(Q/a,q^1){\displaystyle s_{R^t}(Qq^\rho )s_{Q_L}(a)s_P(b)}`$
$`c_{RQ_L}^\alpha s_{\alpha ^t/\eta }(q^\rho )s_{P^t/\eta }(q^\rho )q^{\frac{\kappa _{R^t}}{2}},`$
where the first dilog arises from $`Q_R`$ summation. Now summation over $`P`$ produces another dilog, and we can also sum over $`Q_L`$ and use (Appendix A) to get
$$Z_{Dy;\overline{D}x}^{P^1}=L(Q/a,q^1)L(b,q^1)L(a,q^1)s_R(Qq^\rho )s_{\eta /\alpha }(a)s_\eta (b)s_{R/\alpha }(q^\rho ),$$
Performing the remaining sums over $`R`$, $`\eta `$ and finally $`\alpha `$ gives the crystal result (3.2) (after (14) transformation)
$$Z_{Dy;\overline{D}x}^{P^1}=\frac{Z^{P^1}}{1ab}\frac{L(Q/a,q_{crystal})L(b,q_{crystal})L(a,q_{crystal})}{L(bQ,q_{crystal})}.$$
(22)
Here we contrast the simplicity of crystal computation of the compact form final result to the extensive use of summation formulas and Schur identities in the above vertex computation.
### 5.2 Double $`^1`$
In this section we rederive two-wall crystal amplitudes from the topological vertex perspective. At first we compute the partition function. Denoting the sizes of the right and the left leg by $`t_i`$ (and $`Q_i=e^{t_i}=q^{N_i}`$), respectively for $`i=1,2`$, the vertex rules and some rearrangements give
$`Z^{P^1P^1}`$ $`=`$ $`{\displaystyle \underset{P,R}{}}C_{P^t}(Q_2)^{|R|}C_{PR}(Q_1)^{|P|}C_{R^t}=`$ (23)
$`=`$ $`{\displaystyle \underset{\eta }{}}\left[{\displaystyle \underset{\mu }{}}s_\mu (q^\rho )s_\eta (Q_1q^\rho )s_\mu (Q_1q^\rho )\right]`$
$`\left[{\displaystyle \underset{\nu }{}}s_\nu (q^\rho )s_\eta (Q_2q^\rho )s_\nu (Q_2q^\rho )\right]=`$
$`=`$ $`\mathrm{exp}{\displaystyle \underset{k>0}{}}{\displaystyle \frac{Q_1^kQ_2^k+(Q_1Q_2)^k}{k[k]^2}},`$
This result is the same as the crystal expression (12), up to McMahon function invisible for the vertex. Since this is a partition function without any brane insertions, it is also unaffected by $`q`$ inversion.
The vertex computation with a brane on the right compact leg of double $`^1`$, in (-1) framing also agrees with crystal result. Inserting this brane at position $`D`$ from the middle vertex (Fig. 5), the topological vertex rules lead to the amplitude.
$`Z_D^{P^1P^1}`$ $`=`$ $`{\displaystyle C_{P^t}(Q_2)^{|P|}C_{RQ_L,P,}(Q_1)^{|R|}C_{R^tQ_R}}`$ (24)
$`q^{D|Q_L|}q^{(N_1D)|Q_R|}Tr_{Q_L}VTr_{Q_R}V^1`$
$`\left[(1)^{|Q_LR|+|Q_RR^t|}q^{(\kappa _{Q_LR}+\kappa _{Q_RR^t})/2}\right].`$
As before, we take one-dimensional $`V=q^{N_0+1/2}`$ and absorb the position into $`V`$ as
$$a=q^{D+N_0+1/2}=q^{N_0^{}+1/2}.$$
Performing sums over all representation in the appropriate order leads (after a little effort) to the result
$$Z_D^{P^1P^1}=Z^{P^1P^1}\frac{L(Q_1/a,q_{crystal})L(a,q_{crystal})}{L(bQ_2,q_{crystal})},$$
(25)
which is the same as the crystal answer (13) after enlarging the size of the right leg to $`N_1+1`$ due to the brane insertion.
### 5.3 $`^3`$ amplitudes
The amplitude for several branes on one axis of $`^3`$ can be computed directly from the vertex rules, but since we already have the conifold result it is easiest to take the $`N\mathrm{}`$ limit in (20). This also gives result in $`(1)`$ framing, and substituting (14) we get
$$Z_{Mbranes}^{C^3}=\left[\underset{i<j}{}(1\frac{a_i}{a_j})\right]\underset{i=1}{\overset{M}{}}L(a_i,q_{crystal}),$$
(26)
which is the result for the $`^3`$ crystal, see (1). For one brane it reduces to a single dilogarithm.
For a brane on one leg at position $`a`$ and antibrane on the other at position $`b`$, and in framing $`(1,0)`$, the vertex gives
$$Z_{D;\overline{D}}^{C^3vertex}(a,b)=\frac{1}{1ab}L(a,q_{crystal})L(b,q_{crystal}).$$
which reproduces the crystal answer (4). In this case the vertex rules can be expressed in terms of Hopf link invariants (62)
$$Z_{D;\overline{D}}^{C^3vertex}(a,b)=\underset{P,R}{}W_{PR}(1)^{|P|+|R|}q^{\frac{\kappa _P+\kappa _R}{2}}s_P(a)s_R(b),$$
so that inversing $`q`$ (14) according to our conventions and using (65) proves that this is the same Hopf link generating function as in the crystal case (5).
The calculation for two branes, one in each leg, is similar and also gives the crystal result in $`(1,0)`$ framing<sup>6</sup><sup>6</sup>6 This is also an example of a situation, which can be resummed in canonical framing, with the final result $`L(a,q)L(b,q)\frac{1a\sqrt{q}+ab}{1a\sqrt{q}}`$. This result does not agree with the crystal one (in canonical crystal framing), thus a proper choice of framing is indeed crucial.
$$Z_{D;D}^{C^3vertex}=(1ab)\frac{L(a,q_{crystal})}{L(b,q_{crystal})}.$$
The configuration with two branes on one leg and antibrane on the other is slightly more complicated. The stretched string factors between the two branes on the same leg, at positions $`a_i=q^{M_i+1/2}`$ (for $`i=1,2`$) give an $`(1a_1/a_2)`$ factor. The full amplitude, with antibrane at $`b=q^{N_1+1/2}`$, and in $`(1,0)`$ framing can be written as
$`Z_{2Dy,\overline{D}x}^{vertex}=(1{\displaystyle \frac{a_1}{a_2}}){\displaystyle }C_{P_1P_2,R^t,}s_{P^1}(a_1)s_{P_2}(a_2)s_R(b)(1)^{|R|}\times `$ (27)
$`\times `$ $`\left[(1)^{|P_1P_2|}q^{\frac{\kappa _{P_1P_2}}{2}}\right]`$
$`=`$ $`(1{\displaystyle \frac{a_1}{a_2}})L(b,q^1){\displaystyle c_{P_1P_2}^\alpha s_{\alpha ^t/\eta }(q^\rho )s_{P_1}(a_1)s_{P_2}(a_2)s_\eta (b)}.`$
After performing summations in several steps and substitution (14) we recover the crystal result (1)
$$Z_{2Dy,\overline{D}x}^{vertex}=\frac{1\frac{a_1}{a_2}}{(1a_1b)(1a_2b)}L(a_1,q_{crystal})L(a_2,q_{crystal})L(b,q_{crystal}).$$
(28)
Thus another way to look at the crystal result (28) is provided by the first line in the expansion of (27), which due to (62) can be written in terms of Hopf link invariants (with all components in knot $`(1)`$-framing) as
$`Z_{2Dy,\overline{D}x}^{vertex}`$ $`=`$ $`(1{\displaystyle \frac{a_1}{a_2}}){\displaystyle W_{P_1P_2,R,}s_{P^1}(a_1)s_{P_2}(a_2)s_R(b)}`$
$`\left[(1)^{|P_1P_2|+|R|}q^{\frac{\kappa _{P_1P_2}+\kappa _R}{2}}\right].`$
Taking out the stretched string factors $`(1a_1/a_2)`$, the crystal result is seen as a generating function for $`2+1`$ โnecklaceโ knot invariants. This knot is shown in Figure 6, arising from the tensor product representation of the Hopf link. Because of one-dimensional sources $`V_i=a_i`$, this is a generating function for representations with one row only.
#### 5.3.1 Two legs of $`^3`$ \- general situation
Finally we consider $`m`$ branes on one leg at positions $`a_i=q^{M_i+1/2}`$, and $`n`$ antibranes on the next leg at $`b_i=q^{N_i+1/2}`$. As usual we take all branes in framing $`(1)`$, which makes resummation doable. Using properties of tensor product, the part of the partition function without factors from strings stretching between branes on the same leg (72) (which is denoted by โ) can be written as
$`Z_{m,\overline{n}}^{}`$ $`=`$ $`{\displaystyle }C_{P_1\mathrm{}P_m,R_1^t\mathrm{}R_n^t,}(1)^{_i|R_i|}\left[(1)^{|_jP_j|}q^{\frac{1}{2}\kappa _{_jP_j}}\right]`$ (29)
$`s_{P_1}(a_1)\mathrm{}s_{P_m}(a_m)s_{R_1}(b_1)\mathrm{}s_{R_n}(b_n)=`$
$`=`$ $`{\displaystyle }s_{(P_1^t\mathrm{}P_m^t)/\eta }(q^\rho )s_{P_1}(a_1)\mathrm{}s_{P_m}(a_m)`$
$`s_{(R_1^t\mathrm{}R_m^t)/\eta }(q^\rho )s_{R_1}(b_1)\mathrm{}s_{R_n}(b_n),`$
where it is understood that
$$s_{(P_1^t\mathrm{}P_m^t)/\eta }=\underset{\alpha }{}c_{P_1^t\mathrm{}P_m^t}^\alpha s_{\alpha /\eta }.$$
The antibrane part takes the form (here we write the partial result for the $`R`$ summation only), according to (42)
$$c_{R_1^tR_2^t}^{\beta _1}c_{\beta _1R_3^t}^{\beta _2}\mathrm{}c_{\beta _{n2}R_n^t}^{\beta _{n1}}s_{\beta _{n1}/\eta }(q^\rho )s_{R_1}(b_1)\mathrm{}s_{R_n}(b_n)=$$
$$=s_{\beta _1^t/R_1}(b_2)s_{\beta _2^t/\beta _1^t}(b_3)\mathrm{}s_{\beta _{n1}^t/\beta _{n2}^t}(b_n)s_{\beta _{n1}/\eta }(q^\rho )s_{R_1}(b_1)=$$
$$=s_{\beta _{n1}^t}(b_1,\mathrm{},b_n)s_{\beta _{n1}/\eta }(q^\rho )=$$
$$=L(b_1,q^1)\mathrm{}L(b_n,q^1)s_\eta (b_1,\mathrm{},b_n)$$
(30)
In the same way, the brane part ($`P`$ summation separated) contributes
$$L(a_1,q^1)\mathrm{}L(a_m,q^1)s_\eta (a_1,\mathrm{},a_m)$$
(31)
The remaining summation over $`\eta `$ in (30) and (31) gives factors for strings stretched between all brane/antibrane pairs; also taking into account (72) for each pair of branes (antibranes) on the same leg finally we get (after the q-inversion)
$`Z_{m,\overline{n}}`$ $`=`$ $`\left[\left(1{\displaystyle \frac{a_1}{a_2}}\right)\mathrm{}\left(1{\displaystyle \frac{a_{m1}}{a_m}}\right)\right]\left[\left(1{\displaystyle \frac{b_1}{b_2}}\right)\mathrm{}\left(1{\displaystyle \frac{b_{n1}}{b_n}}\right)\right]`$
$`{\displaystyle \frac{1}{1a_1b_1}}\mathrm{}{\displaystyle \frac{1}{1a_mb_n}}{\displaystyle \underset{i}{}}L(a_i,q_{crystal}){\displaystyle \underset{j}{}}L(b_j,q_{crystal}),`$
and this is the same answer as we found from the crystal (1).
This more general case also can be understood as a generating function of Hopf link invariants corresponding to tensor products of one-row representations, as the first line of (29) can be written using (62) as
$`Z_{m,\overline{n}}^{}`$ $`=`$ $`{\displaystyle }W_{P_1\mathrm{}P_m,R_1\mathrm{}R_n}\left[(1)^{|_jP_j|+|_kR_k|}q^{\frac{1}{2}(\kappa _{_jP_j}+\kappa _{_kR_k})}\right]`$
$`s_{P_1}(a_1)\mathrm{}s_{P_m}(a_m)s_{R_1}(b_1)\mathrm{}s_{R_n}(b_n),`$
where factors from strings stretched between branes on the same leg (72) are taken out. The corresponding knots are shown in Figure 7, for the case of four branes and three antibranes inserted in the geometry.
Thus the crystal generating function can be interpreted in two distinct ways, in the first way described in section 2 it is the generating function of Hopf link invariants for representations with several rows. In the second way (as shown here from the topological vertex point of view) expanded without the stretched string factors it generates necklace (or tensor product) knot invariants with a single row in knot framing $`(1,1)`$.
## 6 B-model example
In the B-model topological amplitudes are computed on the mirror Calabi-Yau geometries. The mirror geometry is described by the general equation $`xyF(u,v)=0`$. To compute the B-model amplitudes we follow the formalism of closely, where the B-model amplitudes are computed as
$$vac|(branes/antibranes)|V,$$
where $`vac|`$ is a vacuum state chosen in a way which ensures that overall fermion number is zero, and
$$|V=\mathrm{exp}\underset{k,l0}{}\left(a_{kl}\psi _{k1/2}\psi _{l1/2}^{}+\stackrel{~}{a}_{kl}\psi _{k1/2}\stackrel{~}{\psi }_{l1/2}^{}\right)|0$$
is a state representing the Riemann surface $`F(u,v)=0`$ branes live on. This Riemann surface might have several asymptotic ends, with branes in each of them; we restrict ourselves putting branes in two of the patches. The quantities in these two patches are denoted without and with tilde respectively, and positions of branes are given by $`e^{u_i}=a_i`$ and respectively $`b_i`$. In B-model picture branes are represented by fermions with standard mode expansions, thus in two patches we have
$`\psi (a)`$ $`=`$ $`{\displaystyle \underset{k}{}}\psi _{k+1/2}e^{(k+1)u_i}={\displaystyle \underset{k}{}}\psi _{k+1/2}a^{k+1},`$
$`\stackrel{~}{\psi }(b)`$ $`=`$ $`{\displaystyle \underset{k}{}}\stackrel{~}{\psi }_{k+1/2}b^{k+1}.`$
Only fermions from the same patch anticommute
$$\{\psi _{k1/2},\psi _{l+1/2}^{}\}=\delta _{k,l},$$
and the bare vacuum is annihilated by all positive modes
$$\psi _{k+1/2}|0=\psi _{k+1/2}^{}|0=\stackrel{~}{\psi }_{k+1/2}|0=\stackrel{~}{\psi }_{k+1/2}^{}|0=0\text{for}k0.$$
In the case of $`^3`$ the state $`|V`$ is determined (up to $`q^{1/6}`$ factors) by
$`a_{kl}`$ $`=`$ $`(1)^ls_{hook(k+1,l+1)}(q^\rho ),`$
$`\stackrel{~}{a}_{kl}`$ $`=`$ $`(1)^lq^{\frac{\kappa _{(l+1)}}{2}}(W_{k+1,l+1}W_{k+1}W_{l+1}),`$
where $`hook(m+1,n+1)`$ is a hook representation with the relevant number of boxes in its row and column, and $`W_{k+1,l+1}`$ is Hopf link invariant for two symmetric representations with relevant number of boxes. For symmetric representation, the value of Casimir is $`\kappa _n=n^2n`$.
Now we put two antibranes in one patch (these in framing $`(1)`$) and a single brane in the other one (Fig. 8).<sup>7</sup><sup>7</sup>7We could have started with two branes and an antibrane, two antibranes and a brane is just slightly more convenient for the B-model computations. The vacuum should be chosen as $`vac|=0|\stackrel{~}{\psi }_{1/2}`$, and in this case the only contribution comes from the third coefficient (with $`1/2`$ factor) in the exponent expansion of $`|V`$. Manipulations with fermion operators lead to
$`vac|\stackrel{~}{\psi }(b)\psi ^{}(a_1)\psi ^{}(a_2)|V^{(1,0)}=`$
$`={\displaystyle \underset{p,t,r0}{}}\stackrel{~}{a}_{pt}\stackrel{~}{a}_{r0}b^{t+1}\left(a_1^{r+1}a_2^{p+1}a_1^{p+1}a_2^{r+1}\right)(1)^{pr}q^{\frac{\kappa _{p+1}+\kappa _{r+1}}{2}}.`$
Performing the summation gives
$$vac|\stackrel{~}{\psi }(b)\psi ^{}(a_1)\psi ^{}(a_2)|V^{(1,0)}=\frac{a_1a_2b}{L(a_1,q)L(a_2,q)L(b,q)}\left(\frac{1}{1a_2b}\frac{1}{1a_1b}\right)$$
$$=\frac{a_1a_2^2b^2}{L(a_1,q)L(a_2,q)L(b,q)}\frac{1\frac{a_1}{a_2}}{(1a_1b)(1a_2b)}.$$
We already know that inversing $`q`$ exchanges branes with antibranes. So if we started with two branes in the first patch and antibrane in the second, we would get dilogs in numerator. This agrees with the crystal result (1), up to the irrelevant overall $`a_1a_2^2b^2`$ factor.
## 7 Summary and discussion
In this paper we investigated the appearance of knot invariants in the construction of Calabi-Yau crystals. Inserting Lagrangian branes, the $`^3`$ crystal naturally generates the leading part of unknot and Hopf link invariants, with arbitrary number of rows. Comparison with the topological vertex gives an alternative view of the crystal generating invariants for Hopf link for tensor product representations (Fig. 7) with a single row.
### 7.1 Connection to Gopakumar-Vafa invariants
The connection to knot invariants is entirely expected from the topological vertex point of view, which is itself constructed from Chern-Simons knot invariants, using open-closed duality. However, the crystal is interesting for it simplicity summing the vertex knot expansions in natural generating functions. These generating functions are always dilogarithms and simple prefactors. We can phrase this as the statement that inserting branes in the crystal generates (open) Donaldson-Thomas invariants (related to open topological string amplitudes) and here we express these Donaldson-Thomas invariants in terms of Chern-Simons invariants. We stress the simplicity of the crystal computing these DT knot generating functions, as compared to other methods.
Open string topological amplitudes can also be derived from Gromov-Witten theory, by counting holomorphic maps with boundaries in Lagrangian submanifolds. These amplitudes can alternatively be computed from the target space point of view using the M-theory perspective , where they contain information about counting of BPS states. Based on the geometric transition picture of the conifold, the Ooguri-Vafa generating function constructed from Chern-Simons invariants can be reformulated in terms of BPS degeneracies counting D2-branes ending on D4-branes as<sup>8</sup><sup>8</sup>8The OV conjecture is naturally formulated in the free energy rather than the partition function.
$$F_{OV}=\underset{n=1}{\overset{\mathrm{}}{}}\underset{R,Q,s}{}\frac{N_{R,Q,s}}{n[n]}e^{n(t_Q+sg_s)}Tr_RV^n$$
(32)
where $`N_{R,Q.s}`$ are the BPS degeneracies labeled by representation, charge and spin content, $`[n]=q^{n/2}q^{n/2}`$ as before; $`t_Q=_Qk`$ is the area of the corresponding cycle, and $`V`$ is the holonomy matrix. For the case of the unknot in $`S^3`$ this precisely gives the Ooguri-Vafa unknot generating function
$$F_{OV}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{TrV^n+TrV^n}{n[n]}e^{nt/2}$$
The connection between Chern-Simons and Gopakumar-Vafa invariants is elaborated in the series of works .
In the case of closed topological strings, Donaldson-Thomas invariants are new invariants which reformulate the Gromov-Witten theory physically in target space language.
Given that the Calabi-Yau crystal naturally computes the closed and open topological string amplitudes in target space language, we certainly expect a natural relation to the partition function and D-brane amplitudes computed by the Gopakumar-Vafa formulation. In fact, this connection can be explicitly seen in our results already. The crystal brane amplitudes are naturally given in terms of dilogarithms, and we can use the exponential expansion of the dilogarithm
$$L(a,q)=e^{_{n=1}^{\mathrm{}}\frac{a^n}{n[n]}}$$
to extract the free energy. Recalling that our holonomy matrix is one dimensional, thus $`TrV`$ is related to $`a_{OV}`$, the free energies of the crystal brane amplitudes clearly are of the Gopakumar-Vafa form (32). This is explicitly checked by the computation of the Ooguri-Vafa generating function inserting the brane in the conifold crystal in section 3.2.
Since all of the brane amplitudes written similarly in terms of dilogarithms, we conjecture the free energies obtained from the crystal brane partition functions (extracted with the exponentiation formula) are natural expansions in the Gopakumar-Vafa invariants. Thus according to this the crystal amplitudes also naturally compute D-brane degeneracies (they can be simply read off from the expression of crystal free energy). This would all fit in the point of view that the crystal Donaldson-Thomas theory is really a target space theory, and as such it simply encodes the target space point of view of D-brane amplitudes. It would be very interesting to explore the connection between the DT, Gopakumar-Vafa and Chern-Simons invariants in more detail.
### 7.2 Open questions
There are several open questions related to our work. In particular, we only investigated the simplest unknot and Hopf link invariants from the crystal point of view. It would be very interesting to find how more complicated knots are generated from the crystal. One way to realize this would be to investigate how skein relations are represented in the crystal language. One can possibly also use the formalism of knot operators to find the crystal representation of more complicated knots, like for example torus knots. Since Lagrangian brane insertions only produce Hopf link and torus knot invariants, it would be very interesting to understand if there are natural geometric objects (like combination of branes, or new classes of branes) which would compute more complicated knots.
Another question is how to represent the full topological open A-model amplitudes in the crystal. While in the topological vertex one inserts stacks of D-branes, in the crystal we only used a single D-brane probe in each stack. That is the holonomy matrix seen in the crystal is one dimensional only, while in the topological vertex it can be arbitrarily large. Finding the representation of holonomy matrix in the crystal would allow to compute the full structure of A-model amplitudes, in particular that would give also multiple row tensor product Hopf link invariants. Introducing the holonomy matrix may have to do with the generalized fermionic operators found in the crystal in .
Another important open problem is how to extend the crystal amplitude computations for more complicated toric geometries. Clearly, one has to glue pieces of crystal geometries to study more complicated toric amplitudes than the double $`^1`$ geometry we studied in this paper. A gluing prescription for toric diagrams involving a partition function only is given in . It would be important to understand how to glue pieces of crystals with D-branes inserted. One way to proceed in this direction is to take guidance from the topological vertex gluing prescriptions, and the clear-cut relations we found between certain class of vertex and crystal brane amplitudes in this paper.
Finally, we note that the Chern-Simons model of the crystal may have a natural connection to the Brownian motion picture of . It would be interesting to investigate this direction further to find a string theory realization of this picture.
### Acknowledgments
We are grateful to Robbert Dijkgraaf and Marcos Marino for enlightening discussions. We also thank Nick Jones, Albrecht Klemm, Asad Naqvi, Takuya Okuda and Jacek Paweลczyk for valuable conversations. P.S. would like to especially thank Robbert Dijkgraaf for all the support, Amsterdam String Theory group for great hospitality and NWO Spinoza Grant for assistance. N.H. would like to thank the YITP at Stony Brook and the Research School for Theoretical Physics at ANU for hospitality. In addition N.H. is supported by a Fletcher Jones graduate fellowship from USC. The work of A.S. is partially supported by the Stichting FOM.
## Appendix A
Many of our formulas make use of properties of symmetric functions. Here we summarize the basic properties and some identities for symmetric functions: Schur polynomials $`s_R`$, elementary $`e_R`$ and complete $`h_R`$ symmetric polynomials, Newton polynomials $`P_R`$.
A symmetric polynomial $`S`$ depends on a partition $`R`$, and its argument is a string of variables $`x=(x_1,x_2,\mathrm{})`$, what we denote by
$$S_R(x)=S_R(x_1,x_2,\mathrm{}).$$
(33)
By $`q^{R+\rho }`$ we understand a string such that $`x_i=q^{R_ii+1/2}`$ for $`i=1,2,\mathrm{}`$, thus
$$S_R(q^{R+\rho })=S_R(q^{R_11/2},q^{R_23/2},\mathrm{}).$$
In particular
$$S_R(q^\rho )=S_R(q^{1/2},q^{3/2},\mathrm{}).$$
(34)
One can concatenate two strings of variables, $`x=(x_1,x_2,\mathrm{})`$ and $`y=(y_1,y_2,\mathrm{})`$, and then use it as an argument of a symmetric polynomial, which is denoted by
$$S_Q(x,y)=S_Q(x_1,x_2,\mathrm{},y_1,y_2,\mathrm{}).$$
One of the simplest examples of symmetric functions are *Newton polynomials*
$$P_R(x)=\underset{n}{}P_{R_i}(x),\text{where}P_n(x)=\underset{i=1}{}x_i^n.$$
(35)
Let us next introduce *elementary* $`e_n(x)`$ and *complete symmetric functions* $`h_n(x)`$, for $`n=0,1,2,\mathrm{}`$, in terms of a generating functions
$`E(t)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e_nt^n={\displaystyle \underset{i}{}}(1+x_it),`$ (36)
$`H(t)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}h_nt^n={\displaystyle \underset{i}{}}{\displaystyle \frac{1}{1x_it}},`$ (37)
and $`h_1=e_1=h_2=e_2=\mathrm{}=0.`$ Then, for a partition $`R=(R_1,R_2,\mathrm{})`$,
$`e_R`$ $`=`$ $`e_{R_1}e_{R_2}\mathrm{}`$
$`h_R`$ $`=`$ $`h_{R_1}h_{R_2}\mathrm{}.`$
For a partition $`R`$, the *Schur function* is defined as
$$s_R(x)=det(h_{R_ii+j})=det(e_{R_i^ti+j}).$$
(38)
Let us introduce Littlewood-Richardson coefficients $`c_{QR}^P`$ as
$$s_{QR}=s_Qs_R=\underset{P}{}c_{QR}^Ps_P,$$
(39)
which have properties
$$c_{QR}^P=c_{Q^tR^t}^{P^t}=c_{RQ}^P,c_R^P=\delta _R^P,$$
(40)
$$c_{QR}^P=0\text{for}|P||Q|+|R|.$$
(41)
It is also convenient to define multiple coefficient
$$c_{R_1\mathrm{}R_n}^P=\underset{\alpha _i}{}c_{R_1R_2}^{\alpha _1}c_{\alpha _1R_3}^{\alpha _2}c_{\alpha _2R_4}^{\alpha _3}\mathrm{}c_{\alpha _{n2}R_n}^P,$$
(42)
in terms of which a multiple tensor product takes the form
$$R_1\mathrm{}R_n=\underset{P}{}c_{R_1\mathrm{}R_n}^PP,$$
(43)
Finally we define *skew Schur functions*
$$s_{Q/R}=\underset{P}{}c_{RP}^Qs_P.$$
(44)
For trivial representation $`R=`$, we have
$$s_{Q/}=s_Q.$$
and
$$\text{If not}QRs_{R/Q}=0.$$
For Schur functions, we have the following identities
$`s_R(cx)`$ $`=`$ $`c^{|R|}s_R(x)`$
$`s_R(q^\rho )`$ $`=`$ $`q^{\kappa _R/2}s_{R^t}(q^\rho )`$
$`s_R(q^\rho )`$ $`=`$ $`(1)^{|R|}s_{R^t}(q^\rho )`$
$`s_Q(q^\rho )s_R(q^{Q+\rho })`$ $`=`$ $`s_R(q^\rho )s_Q(q^{R+\rho }).`$ (45)
Skew Schur functions satisfy
$`s_{Q/R}(cx)`$ $`=`$ $`c^{|Q||R|}s_{Q/R}(x)`$
$`s_{Q/R}(q^\rho )`$ $`=`$ $`(1)^{|Q||R|}s_{Q^t/R^t}(q^\rho ).`$ (46)
In addition, we have the summation formulas for Schur functions
$`{\displaystyle \underset{R}{}}s_R(x)s_R(y)`$ $`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle \frac{1}{1x_iy_j}}`$
$`{\displaystyle \underset{R}{}}s_R(x)s_{R^t}(y)`$ $`=`$ $`{\displaystyle \underset{i,j}{}}(1+x_iy_j)`$ (47)
and for skew Schur functions
$`{\displaystyle \underset{\eta }{}}s_{Q/\eta }(x)s_{R/\eta }(y)`$ $`=`$ $`{\displaystyle \underset{i,j}{}}(1x_iy_j){\displaystyle \underset{\eta }{}}s_{\eta /R}(x)s_{\eta /Q}(y)`$
$`{\displaystyle \underset{\eta }{}}s_{Q/\eta }(x)s_{R/\eta }(y)`$ $`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle \frac{1}{1+x_iy_j}}{\displaystyle \underset{\eta }{}}s_{\eta ^t/R}(x)s_{\eta /Q}(y)`$
$`{\displaystyle \underset{\eta }{}}s_{\eta /R}(x)s_\eta (y)`$ $`=`$ $`s_R(y){\displaystyle \underset{\mu }{}}s_\mu (x)s_\mu (y)`$
$`{\displaystyle \underset{\eta }{}}s_{\eta ^t/R}(x)s_\eta (y)`$ $`=`$ $`s_R(y){\displaystyle \underset{\mu }{}}s_\mu (x)s_{\mu ^t}(y)`$
$`{\displaystyle \underset{\eta }{}}s_{R/\eta }(x)s_{\eta /Q}(y)`$ $`=`$ $`s_{R/Q}(x,y),`$
$`{\displaystyle \underset{\eta }{}}s_{R/\eta }(x)s_\eta (y)`$ $`=`$ $`s_R(x,y).`$ (48)
the last two sums being over partitions $`\eta `$ such that $`Q\eta R`$.
For the special case a partition with a single row $`R=(R_1,0,0,\mathrm{})`$, the Schur function is related to the quantum dilogarithm as
$$s_{R=(R_1,0,\mathrm{})}(q^\rho )=(1)^{R_1}q^{R_1^2/2}\xi (q)L((R_1+\frac{1}{2})g_s,q)$$
(49)
where
$$\xi (q)=\underset{i=1}{\overset{\mathrm{}}{}}\frac{1}{1q^i}.$$
## Appendix B
Here we prove the many-row Hopf-link expansion formula (6) for the simplified case of $`n`$ branes on the positive slice only, whose positions determine the values of $`a_1,\mathrm{},a_n`$. Let us recall, that the normalized crystal partition function in the present case is
$`\stackrel{~}{Z}(a_1,\mathrm{},a_n)`$ $`=`$ $`L(a_1,q)\mathrm{}L(a_n,q)`$
$`\left(1{\displaystyle \frac{a_1}{a_2}}\right)\left(1{\displaystyle \frac{a_1}{a_3}}\right)\mathrm{}\left(1{\displaystyle \frac{a_1}{a_n}}\right)\mathrm{}\left(1{\displaystyle \frac{a_{n1}}{a_n}}\right).`$
In this case the statement (6) takes the form
$$\stackrel{~}{Z}(a_1,\mathrm{},a_n)=\underset{R_1,\mathrm{},R_n}{}a_1^{R_1}\mathrm{}a_n^{R_n}s_{(R_n,R_{n1},\mathrm{},R_1)}(q^\rho ).$$
(51)
We should note, that expansion contains Schur functions corresponding to โimproperโ partitions (with negative number of boxes, or not decreasing in length). But these are taken into account in the proof below automatically, due to structure of Schur functions.
We prove (51) by induction on number of branes $`n`$. The first step in the induction is the expression for the dilogarithm
$$L(a,q)=\underset{R=0}{\overset{\mathrm{}}{}}a^Rh_R(q^\rho ),$$
(52)
as a single variable $`a`$ the sum is over one-row partitions of length $`R`$, and $`s_{(R,0,\mathrm{})}=h_R`$.
In the second induction step, let us assume that $`\stackrel{~}{Z}(a_1,\mathrm{},a_n)`$ is given by (51), and we add one more brane at $`a_0`$. Then
$`\stackrel{~}{Z}(a_1,\mathrm{},a_n,a_0)`$ $`=`$ $`\stackrel{~}{Z}(a_1,\mathrm{},a_n)L(a_0,q)`$ (53)
$`\left(1{\displaystyle \frac{a_1}{a_0}}\right)\mathrm{}\left(1{\displaystyle \frac{a_n}{a_0}}\right).`$
If we expand w.r.t. all $`a_i`$ and use (51) and (52), the coefficient at $`a_0^{R_0}\mathrm{}a_n^{R_n}`$ is equal to (for now we skip arguments) $`(q^\rho )`$)
$$h_{R_0}s_{(R_m,\mathrm{},R_1)}h_{R_0+1}\underset{i=1}{\overset{m}{}}s_{(\widehat{i})}+h_{R_0+2}\underset{ij}{\overset{m}{}}s_{(\widehat{i},\widehat{j})}\mathrm{}h_{R_0+n}s_{(R_n1,\mathrm{},R_11)},$$
(54)
where $`\widehat{i}`$ means, that $`i`$โth variable $`R_i`$ is replaced by $`(R_i1)`$, for example
$$s_{(\widehat{i},\widehat{j})}=s_{(R_n,R_{n1},\mathrm{},R_i1,\mathrm{},R_j1,\mathrm{},R_1)}.$$
(55)
In the first term in this expression no variable is reduced by 1, and in the last term all $`m`$ variables are reduced. In other terms several variables are reduced, and the sums are over all possible combinations of choosing this number of variables from the set $`(R_1,\mathrm{},R_n)`$.
The final observation is that (54) is Laplace expansion of the determinant defining $`s_{(R_0,R_n,\mathrm{},R_1)}`$ along the first row (38)
$$s_{(R_0,R_n,\mathrm{},R_1)}=\begin{array}{cccc}h_{R_0}& h_{R_0+1}& \mathrm{}& h_{R_0+n}\\ h_{R_n1}& h_{R_n}& \mathrm{}& h_{R_n+n1}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ h_{R_1m}& h_{R_1n+1}& \mathrm{}& h_{R_1}\end{array}$$
(56)
where double lines denote determinant. This completes the induction and proves (51). In the more general case for branes in both legs a completely analogous proof can be constructed.
Improper partitions
In the above expansion, it should be stressed that in general not all $`W_{PR}`$ correspond to Hopf link invariants. They correspond only in the case when $`P`$ and $`R`$ are proper partitions, i.e. if $`R_nR_{n1}\mathrm{}R_10`$, and similarly for the representation $`P`$. Otherwise $`W_{PR}`$ are just coefficients resulting from the expansion, but generally these cannot be thought of as Hopf link invariants. Nonetheless, functions $`s_P`$ involved in $`W_{PR}`$ are still given by the determinant $`(\text{38})`$, and thus we will call them *improper Schur functions*.
Moreover, the summations over $`R_i`$ and $`P_i`$ in (6) donโt start from 0, because in the crystal partition function expansions there are also terms with negative powers of $`a_i,b_i`$. These negative powers arise only from prefactors for strings stretched between branes on the same leg, which are of the form $`(1a_i/a_j)`$, and there is always finite number of such terms.
In fact, the easiest way to take care of them is to understand the summations in (6) as running over all integers, positive and negative. The very structure of Schurโs functions, together with the fact that $`h_i=0`$ for $`i<0`$, will assure that only relevant terms will be non-zero, and we get the correct result. In particular, this means that there will be partitions $`R`$ with โnegative number of boxesโ in some rows, $`R_i<0`$. So if we expand determinant (38) for the corresponding *improper Schur functions* $`s_R`$, and use $`h_{i<0}=0`$, we are left with are Schur functions for partitions with lower number of rows, now only of positive length. These new functions can also be proper or not, according to whether lengths of their rows are properly decreasing.
Thus, if we put $`n`$ branes on one leg, the crystal expansion in fact contains information about all proper knot invariants, and finite number of improper knot invariants for partitions with all number of rows $`1,\mathrm{},n`$.
## Appendix C
In this appendix we introduce A-model topological vertex calculational framework. The most convenient form of the vertex is representation basis, in which vertex amplitudes can be expressed in terms of Schur functions. The general formula for topological vertex in the canonical framing is
$$C_{R_1R_2R_3}=q^{\frac{1}{2}(\kappa _{R_2}+\kappa _{R_3})}s_{R_2^t}(q^\rho )\underset{P}{}s_{R_1/P}(q^{R_2^t+\rho })s_{R_3^t/P}(q^{R_2+\rho }).$$
(57)
The crucial property of $`C_{R_1R_2R_3}`$ in the canonical framing is cyclicity w.r.t. representations $`R_i`$. The above formula also immediately implies
$$C_{R_1R_2R_3}=q^{\frac{1}{2}_i\kappa _{R_i}}C_{R_1^tR_3^tR_2^t}.$$
(58)
The identities from appendix Appendix A lead to the following special cases, with some representations involved being trivial $``$
$`C_R`$ $`=`$ $`q^{\kappa _R/2}s_{R^t}(q^\rho )=s_R(q^\rho ),`$ (59)
$`C_{PR}`$ $`=`$ $`q^{\frac{1}{2}\kappa _R}s_P(q^\rho )s_{R^t}(q^{\rho +P})=`$ (60)
$`=`$ $`q^{\frac{\kappa _P}{2}}{\displaystyle \underset{\eta }{}}s_{R/\eta }(q^\rho )s_{P^t/\eta }(q^\rho ).`$ (61)
The vertex with one trivial representation is closely related to the leading term of the Hopf Link invariant $`W_{PR}`$, which also can be expressed in terms of Schur functions
$`W_{PR}`$ $`=`$ $`q^{\kappa _R/2}C_{PR^t}=`$ (62)
$`=`$ $`s_P(q^\rho )s_R(q^{\rho +P})=`$ (63)
$`=`$ $`q^{\frac{1}{2}(\kappa _P+\kappa _R)}{\displaystyle \underset{\eta }{}}s_{R^t/\eta }(q^\rho )s_{P^t/\eta }(q^\rho ),`$ (64)
and it is not difficult to show that
$$W_{PR}(q)=(1)^{|P|+|R|}W_{P^tR^t}(q^1).$$
(65)
The important feature of the vertex is a framing ambiguity, which arises as a need to specify an integer number for each stack of branes on a leg of $`^3`$. The vertex in a particular framing specified by numbers $`f_1,f_2,f_3`$ corresponding to representations $`R_i`$ on different axes is given as
$$C_{R_1R_2R_3}^{f_1,f_2,f_3}=(1)^{_if_i|R_i|}q^{_if_i\kappa _{R_i}/2}C_{R_1R_2R_3},$$
(66)
where $`|R_i|`$ denotes number of boxes in the Young diagram for a given representation. The canonical framing (57) corresponds to $`f_i=0`$.
It is also possible to reverse orientation of the branes on one leg, what can be interpreted as changing branes to antibranes. To obtain vertex amplitude with an antibrane on the first axis one should substitute
$$C_{PQR}(1)^{|P|}C_{P^tQR},$$
(67)
and similarly for any other leg.
To construct the full toric diagram, one has to glue together $`^3`$ patches. Gluing together just two patches gives a resolved conifold with Kahler parameter $`Q=q^N=e^t`$, and the propagator is given by $`(Q)^{|R|}`$. The orientations of two glued axes must be consistent, and sum over representations performed, what leads to
$`Z^{P^1}`$ $`=`$ $`{\displaystyle \underset{R}{}}C_{R^t}(Q)^{|R|}C_R={\displaystyle \underset{R}{}}s_R(q^\rho )s_R(Qq^\rho )=`$ (68)
$`=`$ $`{\displaystyle \underset{i,j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{1Qq^{ij}}}=\mathrm{exp}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Q^k}{k[k]^2}}\right).`$
It is also possible, though a bit more complicated, to consider branes on internal legs of toric diagram and configurations of several stacks of branes on one leg. In the case of one stack of branes on a compact leg one more parameter $`d=g_sD`$ should be introduced, which denotes the position of the brane along the compact leg, as measured from the left vertex. To properly glue two vertexes with additional brane between them, two additional summations must be introduced representing strings ending on the brane from the left and from the right, so the relevant vertex factor takes the form
$$\underset{R,Q^L,Q^R}{}C_{RQ^L}(1)^sq^fe^LC_{R^tQ^R}Tr_{Q^L}VTr_{Q^R}V^1$$
(69)
where a framing of the brane $`p`$ has also been taken into account, so that
$`L`$ $`=`$ $`|R|t+|Q^L|d+|Q^R|(td),`$ (70)
$`f`$ $`=`$ $`{\displaystyle \frac{p}{2}}\kappa _{RQ^L}+{\displaystyle \frac{n+p}{2}}\kappa _{R^tQ^R},`$
$`s`$ $`=`$ $`|R|+p|RQ^L|+(n+p)|R^tQ^L|.`$ (71)
The additional number $`n=|v^{}\times v|`$ is determined by planar directions of two axes $`v`$ and $`v^{}`$ of glued vertexes, and in all cases we consider it equals zero.
For more stacks of branes we need to specify a position of each stack as $`d_i=g_sD_i`$. The holonomy matrix corresponding to branes at $`d_i`$ is denoted $`V_i`$. In fact, to get agreement with crystal results we will need to absorb $`d_i`$ into $`V_i`$. Moreover, in this case we have to choose different representations $`R_i`$ for each stack of branes, and for a given leg of the vertex consider the tensor product of representations $`C_{P,Q,_iR_i}`$. In addition, for each pair of branes at $`d_i,d_j`$ we have to introduce an additional factor from strings stretched between them
$$\underset{P}{}(1)^{|P|}Tr_PV_iTr_{P^t}V_j^1.$$
(72)
If there are several branes on an internal leg, also summations from the left and right vertexes for each brane must be introduced, as well as an overall summation over $`|R|`$ as in (69).
It is also important to make clear how summations over tensor products should be understood. The Hopf Link with a single factor of $`|P_1P_2|`$ can be obtained from a fusion rule
$$W_{P_1P_2,R}=\underset{\alpha }{}c_{P_1P_2}^\alpha W_{\alpha R}=\underset{\alpha }{}q^{\frac{\kappa _\alpha +\kappa _R}{2}}c_{P_1P_2}^\alpha s_{\alpha ^t/\eta }(q^\rho )s_{R^t/\eta }(q^\rho ),$$
(73)
and then related to topological vertex by (62). When a few factors of $`|P_1P_2|`$ appear, the internal summation over $`\alpha `$ should also be introduced<sup>9</sup><sup>9</sup>9We thank Marcos Marino for explaining this point., with each such factor replaced by $`\alpha `$. For example, for two stacks of branes on one leg of $`^3`$ in $`(1)`$ framing we have
$$C_{P_1P_2,R,}(1)^{|P_1P_2|}q^{\frac{\kappa _{P_1P_2}}{2}}=\underset{\alpha }{}c_{P_1P_2}^\alpha s_{\alpha ^t/\eta }(q^\rho )s_{R/\eta }(q^\rho )(1)^{|\alpha |}=$$
$$=(1)^{|P_1|+|P_2|}\underset{\alpha }{}c_{P_1P_2}^\alpha s_{\alpha ^t/\eta }(q^\rho )s_{R/\eta }(q^\rho ),$$
(74)
where two factors of $`q^{\kappa _\alpha /2}`$ (from vertex expression and $`(1)`$ framing) canceled each other, and formulas (40) and (41) have been used.
The internal summation arising in quantities with tensor product involved is a crucial and subtle issue. In particular, knot invariants in different framings but without tensor product differ just by an overall sign and factors of $`q`$. On the other hand, the summation implicit in tensor product formulae changes the structure of polynomials representing knot invariants. For example, in canonical framing we have
$$C_{11,1,}=W_{11,1,}=\frac{(q^2q+1)^2}{(q1)^3\sqrt{q}}.$$
(75)
On the other hand, in framing $`(1,0)`$
$$C_{11,1,}(1)^{|11|}q^{\frac{\kappa _{11}}{2}}=W_{11,1,}(1)^{|11|}q^{\frac{\kappa _{11}}{2}}=$$
$$=W_{1,2}q^1+W_{1,2^t}q=\frac{(2q^23q+2)\sqrt{q}}{(q1)^3},$$
(76)
and this is also precisely the coefficient which we get from crystal expansion without the factor (72) $`(1\frac{a_1}{a_2})`$ (and up to $`q`$ inversion) with two branes at $`a_1,a_2`$ on one slice of the crystal and antibrane on the other slice. |
warning/0506/math0506210.html | ar5iv | text | # Coniveau and the Grothendieck group of varieties
## 1. Filtered Hodge structures
Let $`X`$ be a smooth projective variety. Its cohomology carries a natural Hodge structure. The coniveau filtration on $`H^i(X)`$ is given by
$$N^pH^i(X)=\underset{\text{codim}Sp}{}\text{ker}[H^i(X)H^i(XS)]$$
It is a descending filtration by sub Hodge structures. The largest rational sub Hodge structure $`^pH^i(X)`$ contained in $`F^pH^i(X)`$ gives a second filtration, which we call the level filtration. We have $`N^pH^i(X)^pH^i(X)`$. We will say that $`\text{GHC}(H^i(X),p)`$ holds if $`N^pH^i(X)=^pH^i(X)`$. We will say that generalized Hodge conjecture (GHC) holds for $`X`$ if we have equality for all $`i`$ and $`p`$. Note that $`\text{GHC}(H^{2p}(X),p)`$ is just the usual Hodge conjecture.
We recall that a Hodge structure is polarizable if it admits a polarization, that is a bilinear form satisfying the Hodge-Riemann bilinear relations. We note that the Hodge structure on the cohomology of a smooth projective variety is polarizable: once an ample line bundle is chosen, a polarization is given by taking the orthogonal direct sum of the polarizations, determined in the usual way, on primitive cohomology \[W, p. 202, 207\]. Let $`๐ฎ`$ be the category of finite direct sums of pure rational polarizable Hodge structures. The category $`๐ฎ`$ is a semisimple Abelian category with tensor products \[D, 4.2.3\]. Any object $`H`$ can be decomposed into a sum
$$H=\underset{i}{}H^i$$
where $`H^i`$ is the largest sub Hodge structure of weight $`i`$. We define the category $`๐ฎ`$ whose objects are polarizable Hodge structures with finite descending filtrations by sub Hodge structures, and whose morphisms preserve the filtration. Note that this would be a filtration by polarizable sub Hodge structures, since a sub Hodge structure of a polarizable Hodge structure is again polarizable.
Given an additive category $`๐`$, we can define a Grothendieck group $`\text{K}_0^{split}(๐)`$ by generators and relations as follows. We have one generator $`[X]`$ for each isomorphism class of objects $`X๐`$, and we impose the relation $`[X_3]=[X_1]+[X_2]`$, whenever $`X_3X_1X_2`$. When the category $`๐`$ possesses exact sequences, we can define a quotient $`\text{K}_0(๐)`$ by imposing the above relation when $`X_3`$ is an extension of $`X_2`$ by $`X_1`$. Although it is not strictly necessary for our purposes, we will show in the appendix that these constructions lead to the same groups when applied to $`๐ฎ`$ and $`๐ฎ`$. Consequently, we will usually drop the label โsplitโ in the sequel.
By definition, any additive invariant on $`๐ฎ`$ or $`๐ฎ`$ factors through their
Grothendieck groups. In particular, this remark applies to the Poincarรฉ polynomial
$$P_H(t)=\underset{i}{}dimH^it^i[t,t^1]$$
and a filtered version of it
$$FP_{(H,N)}(t,u)=\underset{i,p}{}dim(N^pH^i)t^iu^p[t^{\pm 1},u^{\pm 1}]$$
We can define two functors between the categories $`๐ฎ`$ and $`๐ฎ`$:
$`\mathrm{\Gamma }`$ $`:๐ฎ๐ฎ;H(H,^{})`$
$`\mathrm{\Phi }`$ $`:๐ฎ๐ฎ;(H,N^{})H`$
where $`^{}`$ is the level filtration on $`H`$, i.e. $`^pH`$ is the largest sub Hodge structure of $`F^pH`$. These functors are clearly additive. Thus we obtain well-defined group homomorphisms $`\gamma `$ and $`\varphi `$, respectively:
$`\gamma `$ $`:\text{K}_0(๐ฎ)\text{K}_0(๐ฎ);\left[H\right]\left[(H,^{})\right]`$
$`\varphi `$ $`:\text{K}_0(๐ฎ)\text{K}_0(๐ฎ);\left[(H,N^{})\right]\left[H\right]`$
Let $`\text{K}_0(๐ฑar_{})`$ denote the Grothendieck group of the category of varieties over $``$ \[DL1\]. A more convenient description for our purposes is provided by \[B, Theorem 3.1\]:
$`\text{K}_0(๐ฑar_{})\text{K}_0^{\text{bl}}(๐ฑar_{}),`$
where $`\text{K}_0^{\text{bl}}(๐ฑar_{})`$ is the free Abelian group generated by isomorphism classes $`[X]`$ of smooth projective varieties subject to the relation $`[\text{Bl}_ZX][E]=[X][Z]`$ for every blow up $`\text{Bl}_ZX`$ of $`X`$ along a smooth closed subvariety $`ZX`$ with the exceptional divisor $`E`$.
Let $`X`$ be a smooth projective variety. Set
$`\left[X\right]_๐ฎ:={\displaystyle \underset{i}{}}(1)^i[H^i(X)]\text{K}_0(๐ฎ)`$
$`\left[X\right]_๐ฎ:={\displaystyle \underset{i}{}}(1)^i[(H^i(X),N^{})]\text{K}_0(๐ฎ)`$
where $`N^{}`$ is the coniveau filtration. In the next section, we will show that these classes depend only on $`[X]\text{K}_0^{\text{bl}}(๐ฑar_{})`$.
## 2. Coniveau of a blow up
We use the following notation throughout this section. Let $`X`$ be a smooth projective variety and let $`\sigma :\stackrel{~}{X}=\text{Bl}_ZXX`$ be the blow up of $`X`$ along a smooth closed subvariety $`Z`$ of $`X`$ of codimension $`2`$. The exceptional divisor $`E`$ can be identified with $`(N_{Z/X})`$, where $`N_{Z/X}`$ is the normal bundle. Therefore it has a tautological line bundle $`๐ช_E(1)`$. Let $`r=\text{codim}(Z,X)1=dimEdimZ`$ and let $`h=c_1(๐ช_E(1))`$.
$`\begin{array}{c}\hfill \text{}\end{array}`$
where $`i,i_1,j,j_1`$ are inclusions.
###### Lemma 2.1.
$$N^pH^i(E)=\sigma ^{}(N^pH^i(Z))(h\sigma ^{}(N^{p1}H^{i2}(Z)))\mathrm{}(h^r\sigma ^{}(N^{pr}H^{i2r}(Z)))$$
###### Proof.
Since $`hN^1H^2(E)`$, $`h^kN^kH^{2k}(E)`$ for each $`k1`$. Hence for any $`\alpha _k\sigma ^{}(N^{pk}H^{i2k}(Z))N^{pk}H^{i2k}(E)`$, we have
$$h^k\alpha _kN^kH^{2k}(E)N^{pk}H^{i2k}(E)N^pH^i(E)$$
by \[AK, Corollary 1.2\] for $`1kr`$. Hence this and $`\sigma ^{}(N^pH^i(Z))N^pH^i(E)`$ give
$$N^pH^i(E)\sigma ^{}(N^pH^i(Z))(h\sigma ^{}(N^{p1}H^{i2}(Z)))\mathrm{}(h^r\sigma ^{}(N^{pr}H^{i2r}(Z)))$$
To show the converse, first note that we have \[Le, Proposition 8.23\]
$`H^i(E)=\sigma ^{}(H^i(Z))\left({\displaystyle \underset{k=1}{\overset{r}{}}}(h^k\sigma ^{}(H^{i2k}(Z)))\right)`$
Hence for any $`\alpha N^pH^i(E)`$, we can decompose
$$\alpha =\sigma ^{}(\alpha _0)+(h\sigma ^{}(\alpha _1))+(h^2\sigma ^{}(\alpha _2))+\mathrm{}+(h^r\sigma ^{}(\alpha _r))$$
where $`\alpha _kH^{i2k}(Z)`$ for each $`k=0,\mathrm{},r`$. We will show that $`\alpha _kN^{pk}H^{i2k}(Z)`$ for each $`k=0,\mathrm{},r`$ by using descending induction on $`k`$.
First note that
$$\sigma _{}(h^i)=\{\begin{array}{cc}0\text{if }i<r\hfill & \\ 1\text{if }i=r\hfill & \end{array}$$
Therefore, by the projection formula and \[AK, Theorem 1.1.(1)\],
$$\alpha _r=\sigma _{}(h^r\sigma ^{}(\alpha _r))=\sigma _{}(\alpha )N^{pr}H^{i2r}(Z)$$
This shows the claim when $`k=r`$. Now suppose that $`\alpha _lN^{pl}H^{i2l}(Z)`$ for $`k+1lr`$. We have to show that $`\alpha _kN^{pk}H^{i2k}(Z)`$. By assumption, for $`k+1lr`$ we have
$$h^l\sigma ^{}(\alpha _l)N^lH^{2l}(E)N^{pl}H^{i2l}(E)N^pH^i(E)$$
Set
$$\beta _{k+1}=\alpha \underset{l=k+1}{\overset{r}{}}(h^l\sigma ^{}(\alpha _l))$$
Then, we have
$`\beta _{k+1}=\sigma ^{}(\alpha _0)+(h\sigma ^{}(\alpha _1))+\mathrm{}+(h^k\sigma ^{}(\alpha _k))N^pH^i(E)`$
Then by taking a cup product with $`h^{rk}N^{rk}H^{2(rk)}(E)`$, we get
$`h^{rk}\beta _{k+1}`$ $`=`$ $`(h^{rk}\sigma ^{}(\alpha _0))+(h^{rk+1}\sigma ^{}(\alpha _1))+\mathrm{}+(h^{rk+k}\sigma ^{}(\alpha _k))`$
$``$ $`N^{p+(rk)}H^{i+2(rk)}(E)`$
Then,
$$\sigma _{}(h^{rk}\beta _{k+1})=\sigma _{}(h^r\sigma ^{}(\alpha _k))=\alpha _kN^{pk}H^{i2k}(Z)$$
as we claimed. Therefore we have
$$N^pH^i(E)\sigma ^{}(N^pH^i(Z))(h\sigma ^{}(N^{p1}H^{i2}(Z)))\mathrm{}(h^r\sigma ^{}(N^{pr}H^{i2r}(Z)))$$
###### Corollary 2.2.
Given $`\alpha N^pH^i(E)`$, we can write
$$\alpha =\sigma ^{}(\alpha _0)+(h\beta )$$
where $`\alpha _0N^pH^i(Z)`$ and $`\beta N^{p1}H^{i2}(E)`$
The following is well known, but we indicate the proof for lack of a suitable reference.
###### Lemma 2.3.
There is a short exact sequence of pure Hodge structures of weight $`i`$:
(1)
###### Proof.
We have a commutative diagram with exact rows:
where some of the arrows are injections or isomorphisms as indicated. The lemma now follows from a straight forward diagram chase. โ
###### Lemma 2.4.
The sequence of the previous lemma is exact in $`๐ฎ`$, i.e.
is an exact sequence for all $`p`$.
###### Proof.
Consider the short exact sequence (1) and note that
$`(\sigma ^{}+i^{})(N^pH^i(X))N^pH^i(\stackrel{~}{X})N^pH^i(Z),`$
$`(j^{}+\sigma ^{})(N^pH^i(\stackrel{~}{X})N^pH^i(Z))N^pH^i(E)`$
since all maps preserve the coniveau.
We will check exactness in the middle. It suffices to show that
$$\text{im}(\sigma ^{}+i^{})|_{N^pH^i(X)}\mathrm{ker}(j^{}+\sigma ^{})|_{N^pH^i(\stackrel{~}{X})N^pH^i(Z)},$$
since the reverse inclusion follows from (1). Let $`(\beta ,\gamma )N^pH^i(\stackrel{~}{X})N^pH^i(Z)`$ such that $`(\beta ,\gamma )\mathrm{ker}(j^{}+\sigma ^{})`$. Then, from the exact sequence (1), there is $`\alpha H^i(X)`$ such that $`(\sigma ^{}+i^{})(\alpha )=(\beta ,\gamma )`$. In particular, we have
$$\sigma ^{}(\alpha )=\beta N^pH^i(\stackrel{~}{X})$$
Since $`\sigma :\stackrel{~}{X}X`$ is a birational map, we have
$$\alpha =(\sigma _{}\sigma ^{})(\alpha )=\sigma _{}(\beta )\sigma _{}(N^pH^i(\stackrel{~}{X}))N^pH^i(X)$$
Hence we have $`(\beta ,\gamma )=(\sigma ^{}+i^{})(\alpha )(\sigma ^{}+i^{})(N^pH^i(X))`$.
We will check surjectivity of $`(j^{}+\sigma ^{})|_{N^pH^i(\stackrel{~}{X})N^pH^i(Z)}`$. Let $`\alpha N^pH^i(E)`$. Then by Corollary 2.2 we can decompose
$$\alpha =\sigma ^{}(\alpha _0)+(h\beta )$$
where $`\alpha _0N^pH^i(Z)`$ and $`\beta N^{p1}H^{i2}(E)`$. Note the composition
is same as cupping with $`[E]|_E`$. Since $`๐ช(E)|_E=๐ช_E(1)`$ and $`h=c_1(๐ช_E(1))`$, we have
$$(j^{}j_{})(\beta )=[E]|_E\beta =h\beta $$
Thus, we have
$$\alpha =\sigma ^{}(\alpha _0)j^{}(j_{}(\beta ))$$
and $`j_{}(\beta )j_{}(N^{p1}H^{i2}(E))N^pH^i(\stackrel{~}{X})`$ and this shows that the map $`(j^{}+\sigma ^{})|_{N^pH^i(\stackrel{~}{X})N^pH^i(Z)}`$ is surjective. โ
Let $`f:XY`$ be a morphism of smooth projective varieties. Then $`f^{}`$ preserves the coniveau and induces maps $`\overline{f^{}}:Gr_N^pH^i(Y)Gr_N^pH^i(X)`$.
###### Corollary 2.5.
The following sequence is exact:
###### Proof.
This follows from lemma A.1. โ
We want to construct well-defined morphisms from $`\text{K}_0^{\text{bl}}(๐ฑar_{})`$ to $`\text{K}_0(๐ฎ)`$ and from $`\text{K}_0^{\text{bl}}(๐ฑar_{})`$ to $`\text{K}_0(๐ฎ)`$. First we need the following lemmas.
###### Lemma 2.6.
The relation
$$[\text{Bl}_ZX]_๐ฎ[E]_๐ฎ=[X]_๐ฎ[Z]_๐ฎ$$
holds. Hence, we have a well-defined group homomorphism
$`\text{K}_0^{\text{bl}}(๐ฑar_{})\text{K}_0(๐ฎ);\left[X\right][X]_๐ฎ`$
###### Proof.
By Lemma 2.3 we have a short exact sequence of pure Hodge structures:
Then by the definition of $`\text{K}_0(๐ฎ)`$, we have
$$[H^i(X)]+[H^i(E)]=[H^i(\stackrel{~}{X})H^i(Z)]=[H^i(\stackrel{~}{X})]+[H^i(Z)]$$
and the result follows immediately by taking alternating sums. โ
We define
$$\lambda :\text{K}_0^{\text{bl}}(๐ฑar_{})\text{K}_0(๐ฎ)$$
as the composition of the map constructed in the above lemma with $`\gamma `$.
###### Lemma 2.7.
The relation
$$[\text{Bl}_ZX]_๐ฎ[E]_๐ฎ=[X]_๐ฎ[Z]_๐ฎ$$
holds. Hence we have a well-defined group homomorphism
$`\nu :\text{K}_0^{\text{bl}}(๐ฑar_{})\text{K}_0(๐ฎ);\nu (\left[X\right])=[X]_๐ฎ`$
###### Proof.
We showed in Lemma 2.4 that
$$0[(H^i(X),N^{})][(H^i(\stackrel{~}{X})H^i(Z),N^{})][(H^i(E),N^{})]0$$
is exact in $`๐ฎ`$ where $`N^p(H^i(\stackrel{~}{X})H^i(Z))=N^pH^i(\stackrel{~}{X})N^pH^i(Z)`$. The rest of the argument is exactly the same as above. โ
###### Remark 2.8.
The proof of Lemma 2.4 shows that the above sequence is split exact. So in particular, as the referee has pointed to us, we can construct the homomorphism
$$\nu :\text{K}_0^{\text{bl}}(๐ฑar_{})\text{K}_0^{split}(๐ฎ)$$
directly, without appealing to the results of the appendix.
## 3. Main theorem
###### Theorem 3.1.
Let $`X`$ be a smooth projective variety. Then, the following statements are equivalent:
1. GHC holds for $`X`$.
2. $`[X]\mathrm{ker}(\nu \lambda )`$.
3. The equality $`FP_{[X]_๐ฎ}(t,u)=FP_{\gamma ([X]_๐ฎ)}(t,u)`$ of filtered Poincarรฉ polynomials holds.
###### Proof.
It is clear that GHC for $`X`$ implies $`[X]_๐ฎ=\gamma ([X]_๐ฎ)`$, or equivalently that $`[X]\mathrm{ker}(\nu \lambda )`$.
Suppose $`[X]\mathrm{ker}(\nu \lambda )`$. Then
$$[X]_๐ฎ=\gamma ([X]_๐ฎ)$$
i.e.
$`{\displaystyle \underset{i}{}}(1)^i[(H^i(X),N^{})]`$ $`=`$ $`{\displaystyle \underset{i}{}}(1)^i[(H^i(X),^{})]`$
Taking the filtered Poincarรฉ polynomial $`FP(t,u)`$ of both sides yields the third statement.
Assume the equality in 3. The coefficient of $`t^i`$ on the left is
$$(1)^i\underset{p}{}dimN^pH^i(X)u^p,$$
while on the right it is
$$(1)^i\underset{p}{}dim^pH^i(X)u^p$$
The equality of these expressions forces $`N^pH^i(X)=^pH^i(X)`$ in this case. โ
###### Remark 3.2.
The coefficient of $`t^iu^p`$ in $`FP_{[X]_๐ฎ}(t,u)`$ is $`(1)^idimN^pH^i(X)`$, and the coefficient of $`t^iu^p`$ in $`FP_{\gamma ([X]_๐ฎ)}(t,u)`$ is $`(1)^idim^pH^i(X)`$. Therefore $`\text{GHC}(H^i(X),p)`$ holds precisely when these coefficients coincide.
Let $`๐^{}=[๐ธ_{}^1]\text{K}_0(๐ฑar_{})`$ denote the Lefschetz object. Under the isomorphism
$$\text{K}_0(๐ฑar_{})\text{K}_0^{\text{bl}}(๐ฑar_{})$$
$`๐^{}`$ maps to $`๐=[^1][\mathrm{}]`$. Let
$$๐=[((1),N^{})]$$
where $`N^{}`$ is the filtration determined by $`N^1/N^2=(1)`$. Then
$$\lambda (๐)=\nu (๐)=๐$$
We have a product on the category $`๐ฎ`$ which is just the tensor product filtered by
$$N^p(H_1H_2)=\underset{r+s=p}{}N^rH_1N^sH_2$$
This gives a commutative ring structure on $`\text{K}_0(๐ฎ)`$. The group $`\text{K}_0^{\text{bl}}(๐ฑar_{})`$ also has a commutative ring structure induced by the product of varieties. It is not clear whether $`\lambda `$ or $`\nu `$ are ring homomorphisms, however we do have:
###### Lemma 3.3.
For any $`\eta \text{K}_0^{\text{bl}}(๐ฑar_{})`$,
$$\lambda (๐\eta )=๐\lambda (\eta )$$
$$\nu (๐\eta )=๐\nu (\eta )$$
###### Proof.
Under the Kรผnneth isomorphism, we have
$$N^pH^i(^1\times X)N^{p1}H^{i2}(X)(1)N^pH^i(\{\mathrm{}\}\times X),$$
and a similar statement holds for $`^p`$. The lemma is an immediate consequence. โ
The element $`๐`$ is invertible, and the above identities guarantee that $`\lambda `$ or $`\nu `$ factor through the localization $`=\text{K}_0^{\text{bl}}(๐ฑar_{})[๐^1]`$.
Recall (\[DL1\]) that there is a decreasing filtration $`F^{}`$ on $``$, where $`F^m`$ is the subgroup of $``$ generated by $`\{[X]๐^j|dimXjm\}`$. Let $`\widehat{}`$ be the completion of the ring $``$ with respect to the filtration $`F^{}`$. A similar filtration (compare \[Loo\]) can be defined on $`\text{K}_0(๐ฎ)`$ by replacing dimension by weights. More precisely, let $`L^m\text{K}_0(๐ฎ)`$ be the subgroup generated by $`\{[(H,N^{})]\mathrm{max}\{iH^i0\}m\}`$. We denote the completion of $`\text{K}_0(๐ฎ)`$ with respect to $`L^{}`$ by $`\widehat{๐ฉ}`$. The weights of the Hodge structure on cohomology of a smooth projective variety of dimension $`d`$ are bounded by $`2d`$. Therefore the induced filtrations $`\lambda (F^{})`$ and $`\nu (F^{})`$ are cofinal with a subfiltration of $`L^{}`$. It follows that we have a commutative diagram
$`\begin{array}{c}\hfill \text{}\end{array}`$
where $`\alpha `$ denotes either $`\lambda `$ or $`\nu `$, and $`f,\tau `$ are the canonical ones.
###### Lemma 3.4.
The homomorphism $`\text{K}_0(๐ฎ)[t^{\pm 1},u^{\pm 1}]`$ given by the filtered Poincarรฉ polynomial factors through the image of $`\tau `$.
###### Proof.
Let $`\eta _mL^m`$. Then for each $`m`$, $`\eta `$ can be expressed as a linear combination of classes of filtered Hodge structures of weight at most $`m`$. Thus the degree of $`FP_\eta (t,u)`$ in $`t`$ is bounded above by $`m`$, for all $`m`$. This is impossible unless $`FP_\eta (t,u)=0`$. โ
###### Corollary 3.5 (to theorem).
Let $`X,Y`$ be smooth projective varieties such that $`X,Y`$ define the same classes in $`\widehat{}`$. Then GHC holds for $`X`$ if and only if GHC holds for $`Y`$.
###### Proof.
If the images of $`[X]`$ and $`[Y]`$ coincide in $`\widehat{}`$ then they coincide in $`\widehat{๐ฉ}`$. Therefore their filtered Poincarรฉ polynomials coincide. The conclusion is now an immediate consequence of Theorem 3.1. โ
From Remark 3.2, we obtain:
###### Corollary 3.6.
Let $`X,Y`$ be smooth projective varieties such that $`X,Y`$ define the same classes in $`\widehat{}`$. Then for each $`i`$ and $`p`$, $`\text{GHC}(H^i(X),p)`$ holds if and only if $`\text{GHC}(H^i(Y),p)`$ holds.
The proof of the next corollary depends on the motivic integration theory of Kontsevich, Denef and Loeser. See \[DL1\], \[DL2\], and \[Loo\] for an introduction to these ideas.
###### Corollary 3.7.
Let $`X,Y`$ be $`K`$-equivalent smooth projective varieties, i.e. there is a smooth projective variety $`Z`$ and birational maps $`\pi _1:ZX`$ and $`\pi _2:ZY`$ such that $`\pi _1^{}K_X=\pi _2^{}K_Y`$, where $`K_X`$ (respectively $`K_Y`$) is the canonical divisor on $`X`$ (respectively $`Y`$). Then $`\text{GHC}(H^i(X),p)`$ holds if and only if $`\text{GHC}(H^i(Y),p)`$ holds.
###### Proof.
It is enough to show that $`X`$ and $`Y`$ define the same class in $`\widehat{}`$. This follows from the $`K`$-equivalence assumption by a standard application of motivic integration theory, see \[Loe\] or \[V\]. For convenience of the reader, we reproduce the argument. By the change of variables formula \[DL1, Lemma 3.3\], we have
$`f([X])={\displaystyle _{(X)}}๐\mu _X`$ $`=`$ $`{\displaystyle _{(Z)}}๐^{ord_t\pi _1^{}\omega _X}๐\mu _Z`$
$`=`$ $`{\displaystyle _{(Z)}}๐^{ord_t\pi _2^{}\omega _Y}๐\mu _Z={\displaystyle _{(Y)}}๐\mu _Y=f([Y])`$
where $`f:\text{K}_0^{\text{bl}}(๐ฑar_{})\widehat{}`$ is the canonical map, $`(X),(Y)`$ are the arc spaces, $`\omega _X,\omega _Y`$ are the canonical sheaves, and $`d\mu _X,d\mu _Y`$ are the motivic measures. Hence the corollary follows from Corollary 3.5. โ
###### Corollary 3.8.
Let $`X,Y`$ be birational Calabi-Yau varieties. Then $`\text{GHC}(H^i(X),p)`$ holds if and only if $`\text{GHC}(H^i(Y),p)`$ holds.
###### Proof.
Since $`K_X=K_Y=0`$, it follows from Corollary 3.7. โ
## Appendix A Grothendieck groups of filtered categories
We recall that an exact category consists of an additive category $`๐`$, together with a distinguished class of diagrams
$$0ABC0$$
called exact sequences satisfying appropriate conditions \[Q\]. For example, any additive category $`๐`$ can be made exact by taking the class of exact sequences to be isomorphic to the class of split sequences
$$0AACC0.$$
An Abelian category gives another example of an exact category, where exact sequences have the usual meaning. If the category is also semisimple, then this exact structure coincides with the split structure above. This remark applies to $`๐ฎ`$.
Let $`๐`$ be an Abelian category, and let $`๐`$ (respectively $`๐ข๐`$) denote the category of filtered (respectively graded) objects in $`๐`$. The category $`๐ข๐`$ is Abelian, but $`๐`$ is generally not. However $`๐`$ has a natural exact structure \[BBD, 1.1.4\] given as follows. We have a functor
$$(H,N^{})\underset{p}{}N^pH$$
from $`๐`$ to $`๐ข๐`$. We declare a sequence
$$0(H_1,N^{})(H_2,N^{})(H_3,N^{})0$$
in $`๐`$ to be exact if and only if its image in $`๐ข๐`$ is exact. For the record, we note the following alternative formulation, which is perhaps more common:
###### Lemma A.1.
The sequence
$$0(H_1,N^{})(H_2,N^{})(H_3,N^{})0$$
is exact if and only if the following sequence is exact in $`๐ข๐`$:
$$0\underset{p}{}N^pH_1/N^{p+1}H_1\underset{p}{}N^pH_2/N^{p+1}H_2\underset{p}{}N^pH_3/N^{p+1}H_30$$
###### Proof.
This is a straight forward application of the Snake lemma and induction. โ
###### Lemma A.2.
The category $`๐`$ with the above notion of exact sequence is an exact category.
Given an exact category $`๐`$, its Grothendieck group $`\text{K}_0(๐)`$ is given by generators $`[M]`$, with $`M๐`$, and relations $`[M_2]=[M_1]+[M_3]`$ for every exact sequence $`0M_1M_2M_30`$. Let us denote the Grothendieck group for $`๐`$ with its split exact structure by $`\text{K}_0^{split}(๐)`$. We see immediately that, $`\text{K}_0^{split}(๐)\text{K}_0(๐)`$ if $`๐`$ is Abelian and semisimple. This is, in particular, the case for $`๐ฎ`$. We have a homomorphism $`\text{K}_0^{split}(๐ฎ)\text{K}_0(๐ฎ)`$, which is also an isomorphism by:
###### Lemma A.3.
If $`๐`$ is semisimple Abelian, then any exact sequence in $`๐`$ is split exact.
###### Proof.
It is enough to check that given an exact sequence in $`๐`$
(4)
$$0(H_1,N^{})(H_2,N^{})\stackrel{f}{}(H_3,N^{})0,$$
there is a splitting $`s:(H_3,N^{})(H_2,N^{})`$ for $`f`$. First note that by semisimplicity of $`๐`$, we have a noncanonical decomposition
$$N^pH_iN^{p+1}H_iGr_N^pH_i$$
where $`Gr_N^pH_i=N^pH_i/N^{p+1}H_i`$, $`i=2,3`$. We define $`s_p=s|_{N^pH_3}`$ by descending induction on $`p`$ : Note that the exact sequence (4) induces exact sequences in $`๐`$
$$0N^{p+1}H_1N^{p+1}H_2\stackrel{f_{p+1}}{}N^{p+1}H_30$$
$$0Gr_N^pH_1Gr_N^pH_2\stackrel{\overline{f}_p}{}Gr_N^pH_30$$
where $`f_{p+1}=f|_{N^{p+1}H_2}`$ and $`\overline{f}_p`$ is the induced map. By induction and semisimplicity of $`๐`$, there are splittings $`s_{p+1}:(N^{p+1}H_3,N^{}N^{p+1}H_3)(N^{p+1}H_2N^{}N^{p+1}H_2)`$ and $`t_p:Gr_N^pH_3Gr_N^pH_2`$ for $`f_{p+1}`$ and $`\overline{f}_p`$, respectively. Set
$$s_p=s_{p+1}+t_p:N^pH_3N^pH_2$$
Then $`s_p`$ gives a well-defined splitting for $`f_p`$ and hence we have a splitting $`s=s_0`$ for $`f`$. This completes the proof of the Lemma. โ
###### Corollary A.4.
$`\text{K}_0^{split}(๐ฎ)\text{K}_0(๐ฎ)`$ |
warning/0506/math-ph0506028.html | ar5iv | text | # A family of hyperbolic spin Calogero-Moser systems and the spin Toda lattices
### 1. Introduction
In the theory of integrable systems, a wide range of important examples are covered by the Adler-Kostant-Symes scheme and its generalization known as classical r-matrix theory (see \[A\], \[K\], \[S\], \[RSTS1\], \[RSTS2\], \[AvM\], \[STS1\],\[STS2\], \[RSTS3\], \[FT\], \[LP1\] and the references therein). As is well-known, classical r-matrices are naturally associated with Poisson structures on Lie groups and duals of Lie algebras and the corresponding geometric objects have been used with great success in the solutions of many integrable Hamiltonian systems.
In the early 90โs, dynamical analog of the classical r-matrices was discovered in the study of Wess-Zumino-Witten (WZN) conformal field theory \[BDF\], \[F\]. Since then, these objects have cropped up in other areas as well (see, for example, \[BAB\], \[ABB\], \[Lu\], \[AM\]) and their geometric meaning was unraveled by Etingof and Varchenko in their fundamental paper \[EV\]. While classical r-matrices play a role in Poisson Lie group theory \[D\], the authors in \[EV\] showed that an appropriate geometric setting for the classical dynamical r-matrices is that of a special class of Poisson groupoids (a notion due to Weinstein \[W1\]), the so-called coboundary dynamical Poisson groupoids. If $`R`$ is an $`H`$-equivariant classical dynamical r-matrix, and $`(\mathrm{\Gamma },\{,\}_R)`$ is the associated coboundary dynamical Poisson groupoid, then it follows from Weinsteinโs coisotropic calculus \[W1\] or otherwise that the Lie algebroid dual $`A^{}\mathrm{\Gamma }`$ also has a natural Lie algebroid structure \[LP2\], \[BKS\]. We shall call $`A^{}\mathrm{\Gamma }`$ the coboundary dynamical Lie algebroid associated to $`R`$ and it is this class of Lie algebroids which we use in the study of integrable systems in \[LX2\] and in the present work.
Our purpose in this paper is twofold. First of all, we will continue to develop a general scheme (which we initiated in \[LX2\]) to study integrable systems based on realization in the dual bundles of coboundary dynamical Lie algebroids. To summarize, the class of invariant Hamiltonian systems which admits such a realization (for the genuinely dynamical case) has the following key features: (a) the systems are defined on a Hamiltonian $`H`$-space $`X`$ with equivariant momentum map $`J`$ and the Hamiltonians are the pull-back of natural invariant functions by an $`H`$-equivariant realization map, (b) the pullback of natural invariant functions do not Poisson commute everywhere on $`X`$, but they do so on a fiber $`J^1(\mu )`$ of the momentum map, (c) the reduced Hamiltonian systems on $`X_\mu =J^1(\mu )/H_\mu `$ ($`H_\mu `$ is the isotropy subgroup at $`\mu `$) admit a natural collection of Poisson commuting integrals. In this work, our main focus is on the case in which $`R`$ is a solution of the modified dynamical Yang-Baxter equation (mDYBE). As we pointed out in the announcement \[L1\], the (mDYBE) is associated with a factorization problem on the trivial Lie groupoid $`\mathrm{\Gamma }`$. By making use of the algebraic and geometric structures associated with (mDYBE) (which will be fully worked out here), we will develop an effective method to integrate the Hamiltonian flows on $`J^1(\mu )`$ (which parallels the one announced in \[L1\] for the groupoid framework) based on this factorization. Hence we can obtain the integrable flows on $`X_\mu `$ by reduction.
Our second purpose in this work is to give two important class of new examples and to illustrate our factorization theory using these examples. Our first class of examples is a family of hyperbolic spin Calogero-Moser (CM) systems and their associated integrable models, corresponding to the solutions of (mDYBE) for pairs $`(๐ค,๐ฅ)`$ of Lie algebras, as classified in \[EV\]. Here, $`๐ค`$ is simple, and $`๐ฅ๐ค`$ is a Cartan subalgebra. As such, our systems are parametrized by subsets $`\pi ^{}`$ of a simple system of roots $`\pi `$. Note that in the special case where $`\pi ^{}=\pi `$, our corresponding integrable model is isomorphic to the one in \[R\], and the $`sl(N)`$ case has also appeared in \[AB\], \[KBBT\],for example (see Remark 5.5). The second class of examples was actually discovered when an appropriate scaling limit is applied to the hyperbolic spin CM systems (the ones which are not integrable). Remarkably, the obstruction to integrability dissolves in the scaling limit, leading to a family of integrable models which we will call the spin Toda lattices (again, these are parametrized by subsets $`\pi ^{}`$ of a simple system). As it turns out, the spin Toda lattices are systems which admit realization in the dual bundle $`๐ธ`$ of the coboundary dynamical Lie algebroid $`๐ธ^{}T๐ฅ\times ๐ค`$ associated to the standard r-matrix and the reduction of these $`H`$-invariant systems lead to a family of Toda lattices parametrized by $`\pi ^{}`$. So this gives us a first nontrivial example in which a constant r-matrix is relevant.
The paper is organized as follows. In Section 2, we derive an intrinsic expression for the Lie-Poisson structure on the dual bundle of a coboundary dynamical Lie algebroid which is important for subsequent developments of our program. We also give the complete set of equations for a natural class of invariant Hamiltonian systems. In Section 3, we reprove (essentially) Theorem 3.10 in \[LX2\] using an intrinsic point of view, without having to assume the existence of an $`H`$-equivariant map $`g:XH`$ (also we do not assume $`๐ฅ=Lie(H)`$ is Abelian). We also compute how the realization map evolves under our invariant Hamiltonian systems on $`X`$ based on the development in Section 2. As the reader will see, the integrable flows on the reduced space $`X_\mu `$ are actually realized on a Poisson quotient of a coisotropic submanifold of $`A\mathrm{\Gamma }`$, which in some sense is the analog of the gauge group bundle in \[L1\]. In Section 4, we discuss the algebraic and geometric structures associated with (mDYBE), leading up to a factorization method for solving the Hamiltonian flows. In Section 5, we introduce a family of hyperbolic spin Calogero-Moser systems using Proposition 4.2 (a) and consider the associated integrable models. Then we consider scaling limits of the hyperbolic spin CM systems at the levels of the Hamiltonians, the equations of motion and the (generalized) Lax equations. At the end of the section, we work out the realization picture for the spin Toda lattices and also consider their reduction. Section 6 is concerned with the solution of the hyperbolic spin CM systems and the spin Toda lattices, utilizing the factorization method of Section 4. Here, the reader will see how the concrete factorization problems are being solved. In a remark, we will also discuss the solution of a family of hyperbolic spin Ruijenaars-Schneider models (introduced in \[L1\] and related to the affine Toda field theories \[BH\]) in the general case. We shall address the complete integrability and other aspects of the integrable models here in subsequent publications. For the solution of the systems in \[LX2\] using the method developed here, we refer the reader to the forthcoming work \[L2\] (see also Remark 5.5 (c)). Acknowledgments. The author would like to thank the referee for a helpful question which has led him to put back a missing keyword in the formulation of several results in Section 4. Special thanks are also due to Reeva Goldsmith for converting the AMS-TeX file to LaTeX.
### 2. Coboundary dynamical Lie algebroids and Lie-Poisson structures faak on their dual bundles
In this section, our main goal is to derive an intrinsic formula for the Lie-Poisson structure on the dual bundle of a coboundary dynamical Lie algebroid which is important for subsequent developments. As the reader will see, the same method of calculation is also used in Section 5 to write down the Lie-Poisson structure associated with a trivial Lie algebroid, whose vertex Lie algebra is given by a semi-direct product.
We begin by recalling the definition of a Lie algebroid.
###### Definition Definition 2.1
A Lie algebroid over a manifold $`M`$ is a smooth vector bundle $`\pi _A:AM`$ equipped with a Lie bracket $`[,]_A`$ on its space $`Sect(M,A)`$ of smooth sections and a bundle map $`a_A:ATM`$ (called the anchor map) such that (a) the bundle map $`a_A`$ induces a Lie algebra homomorphism $`Sect(M,A)Sect(TM)`$ (which we also denote by $`a_A`$), (b) for any $`X`$, $`YSect(M,A)`$ and $`fC^{\mathrm{}}(M)`$, the Leibnitz identity
$$[X,fY]_A=f[X,Y]_A+(a(X)f)Y$$
holds.
Let $`(A,[,]_A,a_A)`$ be a Lie algebroid over $`M`$, and let $`\pi _A^{}:A^{}M`$ be its dual bundle. For any $`XSect(M,A)`$, we can associate a smooth function $`l_X`$ on $`A^{}`$ by putting $`l_X(\xi )=<\xi ,X\pi _A^{}(\xi )>`$ for all $`\xi A^{}`$.
###### Theorem 2.2 \c{CDW}
There exists a unique Poisson structure on $`A^{}`$ (called the Lie-Poisson structure) which is characterized by the property
$$\{l_X,l_Y\}=l_{[X,Y]_A}$$
for all $`X`$, $`YSect(M,A)`$.
Remark 2.3 In \[C\] and \[W2\],there are two extra conditions in the characterization of the Lie-Poisson structure, namely, $`\{f\pi _A^{},g\pi _A^{}\}=0`$, and $`\{l_X,f\pi _A^{}\}=(a(X)f)\pi _A^{}`$ for all $`f`$, $`gC^{\mathrm{}}(M)`$, and $`XSect(M,A)`$. However, it is not hard to show that these are actually consequences of $`\{l_X,l_Y\}=l_{[X,Y]_A}`$ and the properties of the Lie algebroid bracket $`[,]_A`$. We shall use these two conditions in an essential way in (2.10) below.
We now recall the definition of a coboundary dynamical Lie algebroid. Let $`G`$ be a connected Lie group, and $`HG`$ a connected Lie subgroup. We shall denote by $`๐ค`$ and $`๐ฅ`$ the corresponding Lie algebras and let $`\iota :๐ฅ๐ค`$ be the Lie inclusion. Let $`U๐ฅ^{}`$ be a connected $`Ad_H^{}`$-invariant open subset, and let $`R:UL(๐ค^{},๐ค)`$ be a classical dynamical r-matrix (here and henceforth we denote by $`L(๐ค^{},๐ค)`$ the set of linear maps from $`๐ค^{}`$ to $`๐ค`$), i.e. $`R`$ is pointwise skew symmetric
$$<R(q)(A),B>=<A,R(q)B>$$
$`(2.1)`$
and satisfies the classical dynamical Yang-Baxter condition
$$\begin{array}{cc}& [R(q)A,R(q)B]+R(q)(ad_{R(q)A}^{}Bad_{R(q)B}^{}A)\hfill \\ \hfill +& dR(q)\iota ^{}A(B)dR(q)\iota ^{}B(A)+d<R(A),B>(q)=\chi (A,B),\hfill \end{array}$$
$`(2.2)`$
for all $`qU`$, and all $`A,B๐ค^{}`$, where $`\chi :๐ค^{}\times ๐ค^{}๐ค`$ is $`G`$-equivariant. The dynamical $`r`$-matrix is said to be $`H`$-equivariant if and only if
$$R(Ad_{h^1}^{}q)=Ad_hR(q)Ad_h^{}$$
$`(2.3)`$
for all $`hH,qU.`$ We shall equip $`\mathrm{\Gamma }=U\times G\times U`$ with the trivial Lie groupoid structure over $`U`$ \[M\] with target and source maps
$$\alpha (u,g,v)=u,\beta (u,g,v)=v$$
$`(2.4)`$
and multiplication map
$$m((u,g,v),(v,g^{},w))=(u,gg^{},w).$$
$`(2.5)`$
Recall that associated with an $`H`$-equivariant classical dynamical r-matrix $`R`$ is a natural Poisson structure $`\{,\}_R`$ on $`\mathrm{\Gamma }`$ \[EV\] such that the pair $`(\mathrm{\Gamma },\{,\}_R)`$ is a Poisson groupoid in the sense of Weinstein \[W1\] (see \[L1\] for the intrinsic forms of $`\{,\}_R`$). Let $`A\mathrm{\Gamma }:=_{qU}T_{ฯต(q)}\alpha ^1(q)=_{qU}\{0_q\}\times ๐ค\times ๐ฅ^{}`$ be the Lie algebroid of $`\mathrm{\Gamma }.`$ Then by Weinsteinโs coisotropic calculus \[W1\] or otherwise, the Lie algebroid dual $`A^{}\mathrm{\Gamma }=_{qU}\{0_q\}\times ๐ค^{}\times ๐ฅ`$ also has a natural Lie algebroid structure \[BKS\],\[LP2\] such that the pair $`(A\mathrm{\Gamma },A^{}\mathrm{\Gamma })`$ is a Lie bialgebroid in the sense of Mackenzie and Xu \[MX\]. We shall denote the Lie brackets on $`Sect(U,A\mathrm{\Gamma })`$ and $`Sect(U,A^{}\mathrm{\Gamma })`$ respectively by $`[,]_{A\mathrm{\Gamma }}`$ and $`[,]_{A^{}\mathrm{\Gamma }}.`$ Throughout the paper, the pair $`(A^{}\mathrm{\Gamma },[,]_A^{})`$ together with the anchor map $`a_{}:A^{}\mathrm{\Gamma }TU`$ given by
$$a_{}(0_q,A,Z)=(q,\iota ^{}Aad_Z^{}q)$$
$`(2.6)`$
will be called the coboundary dynamical Lie algebroid associated to $`R`$. Explicitly, the Lie bracket $`[,]_{A^{}\mathrm{\Gamma }}`$ on $`Sect(U,A^{}\mathrm{\Gamma })`$ is given by the following expression \[BKS\],\[LP2\]:
$$\begin{array}{cc}& [(0,A,Z),(0,A^{},Z^{})]_{A^{}\mathrm{\Gamma }}(q)\hfill \\ \hfill =& (0_q,dA^{}(q)(\iota ^{}A(q)ad_{Z(q)}^{}q)dA(q)(\iota ^{}A^{}(q)ad_{Z^{}(q)}^{}q)\hfill \\ & ad_{R(q)A(q)Z(q)}^{}A^{}(q)+ad_{R(q)A^{}(q)Z^{}(q)}^{}A(q),\hfill \\ & dZ^{}(q)(\iota ^{}A(q)ad_{Z(q)}^{}q)dZ(q)(\iota ^{}A^{}(q)ad_{Z^{}(q)}^{}q)\hfill \\ & [Z,Z^{}](q)+<dR(q)()A(q),A^{}(q)>)\hfill \end{array}$$
$`(2.7)`$
where $`A,A^{}:U๐ค^{}`$, $`Z,Z^{}:U๐ฅ`$ are smooth maps and $`<dR(q)()A(q),A^{}(q)>`$ is the element in $`๐ฅ`$ whose pairing with $`\lambda ๐ฅ^{}`$ is $`<dR(q)(\lambda )A(q),A^{}(q)>.`$
In the rest of the section, we shall make the identifications
$$A\mathrm{\Gamma }U\times ๐ฅ^{}\times ๐ค,A^{}\mathrm{\Gamma }U\times ๐ฅ\times ๐ค^{}.$$
$`(2.8)`$
Let us fix a point $`(q,\lambda ,X)A\mathrm{\Gamma }`$. In order to derive an intrinsic expression for the Lie-Poisson bracket $`\{\phi ,\psi \}_{A\mathrm{\Gamma }}(q,\lambda ,X)`$ on the dual bundle $`A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$, we need to introduce some notation. To start with, let $`Pr_i`$ be the projection map onto the $`i`$-th factor of $`U\times ๐ฅ^{}\times ๐คA\mathrm{\Gamma }`$, $`i=1,2,3.`$ If $`\phi C^{\mathrm{}}(A\mathrm{\Gamma })`$, we have $`d\phi (q,\lambda ,X)=(\delta _1\phi ,\delta _2\phi ,\delta \phi )`$, where the partial derivatives are defined by
$$\begin{array}{cc}& <\delta _1\phi ,\mu >=\frac{d}{dt}_{|_{t=0}}\phi (q+t\mu ,\lambda ,X),<\delta _2\phi ,\mu >=\frac{d}{dt}_{|_{t=0}}\phi (q,\lambda +t\mu ,X),\mu ๐ฅ^{}\hfill \\ & <\delta \phi ,Y>=\frac{d}{dt}_{|_{t=0}}\phi (q,\lambda ,X+tY),Y๐ค.\hfill \end{array}$$
We also associate with $`\phi `$ the function $`\stackrel{~}{\phi }`$ on $`U`$, defined by $`\stackrel{~}{\phi }(u)=\phi (u,\lambda ,X)`$. On the other hand, $`s(\phi ):UU\times ๐ฅ\times ๐ค^{}`$ will denote the constant section of $`U\times ๐ฅ\times ๐ค^{}`$ given by $`s(\phi )(u)=(u,\delta _2\phi ,\delta \phi )`$, where $`\delta _2\phi `$, $`\delta \phi `$ are the partial derivatives evaluated at the fixed point $`(q,\lambda ,X)`$.
Now, it is easy to check by a direct calculation that $`d(\stackrel{~}{\phi }Pr_1)(q,\lambda ,X)=(\delta _1\phi ,0,0)`$, while $`dl_{s(\phi )}(q,\lambda ,X)=(0,\delta _2\phi ,\delta \phi )`$. Thus we have
$$d\phi (q,\lambda ,X)=d(l_{s(\phi )}+\stackrel{~}{\phi }Pr_1)(q,\lambda ,X).$$
$`(2.9)`$
Therefore,
$$\begin{array}{cc}& \{\phi ,\psi \}_{A\mathrm{\Gamma }}(q,\lambda ,X)\hfill \\ \hfill =& \{l_{s(\phi )}+\stackrel{~}{\phi }Pr_1,l_{s(\psi )}+\stackrel{~}{\psi }Pr_1\}_{A\mathrm{\Gamma }}(q,\lambda ,X)\hfill \\ \hfill =& l_{[s(\phi ),s(\psi )]_{A^{}\mathrm{\Gamma }}}(q,\lambda ,X)+d\stackrel{~}{\psi }(q)a_{}(s(\phi ))(q)\hfill \\ & d\stackrel{~}{\phi }(q)a_{}(s(\psi ))(q).\hfill \end{array}$$
$`(2.10)`$
By using the expression for $`[,]_{A^{}\mathrm{\Gamma }}`$ in (2.7), we have
$`l_{[s(\phi ),s(\psi )]_{A^{}\mathrm{\Gamma }}}(q,\lambda ,X)`$
$`=`$ $`<\lambda ,[\delta _2\phi ,\delta _2\psi ]+<dR(q)()\delta \phi ,\delta \psi >>`$
$`+<X,ad_{R(q)\delta \phi \delta _2\phi }^{}\delta \psi +ad_{R(q)\delta \psi \delta _2\psi }^{}\delta \phi >.`$
Meanwhile, from the expression for the anchor map $`a_{}`$, we find
$$d\stackrel{~}{\phi }(q)a_{}(s(\phi ))(q)=<\delta _1\phi ,\iota ^{}\delta \psi ad_{\delta _2\psi }^{}q>.$$
Assembling the calculations, we have the following result.
###### Theorem 2.4
The Lie-Poisson structure on the dual bundle $`A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$ is given by
$$\begin{array}{cc}& \{\phi ,\psi \}_{A\mathrm{\Gamma }}(q,\lambda ,X)\hfill \\ \hfill =& <\lambda ,[\delta _2\phi ,\delta _2\psi ]>+<dR(q)(\lambda )\delta \phi ,\delta \psi >\hfill \\ & +<X,ad_{R(q)\delta \phi \delta _2\phi }^{}\delta \psi +ad_{R(q)\delta \psi \delta _2\psi }^{}\delta \phi >\hfill \\ & <q,[\delta _2\phi ,\delta _1\psi ]+[\delta _1\phi ,\delta _2\psi ]>+<\delta _1\psi ,\iota ^{}\delta \phi ><\delta _1\phi ,\iota ^{}\delta \psi >.\hfill \end{array}$$
Remark 2.5 In a similar fashion, we can show that the Lie-Poisson bracket on the dual bundle $`A^{}\mathrm{\Gamma }`$ of the trivial Lie algebroid $`(A\mathrm{\Gamma },[,]_{A\mathrm{\Gamma }},a)`$ is given by $`\{\phi ,\psi \}_{A^{}\mathrm{\Gamma }}(q,p,\xi )=<\delta _2\phi ,\delta _1\psi ><\delta _1\phi ,\delta _2\psi >+<\xi ,[\delta \phi ,\delta \psi ]>.`$ The reader is referred to Proposition 5.10 below for the details of a similar calculation.
If $`(P,\{,\}_P)`$ is a Poisson manifold, then for each $`fC^{\mathrm{}}(P)`$, we shall define the associated Hamiltonian vector field $`X_f`$ using the convention $`X_f.g=\{f,g\}_P`$.
###### Corollary 2.6
The Hamiltonian vector field on $`A\mathrm{\Gamma }`$ associated to $`\phi C^{\mathrm{}}(A\mathrm{\Gamma })`$ is given by
$`X_\phi (q,\lambda ,X)`$
$`=`$ $`(\iota ^{}\delta \phi ad_{\delta _2\phi }^{}q,ad_{\delta _2\phi }^{}\lambda +\iota ^{}ad_X^{}\delta \phi ad_{\delta _1\phi }^{}q,`$
$`[X,R(q)\delta \phi \delta _2\phi ]+dR(q)(\lambda )\delta \phi \delta _1\phi +R(q)(ad_X^{}\delta \phi )).`$
Now, a natural collection of invariant functions on $`A\mathrm{\Gamma }`$ is $`Pr_3^{}I(๐ค)`$, where $`I(๐ค)`$ is the ring of ad-invariant functions on $`๐ค`$. The following result is an easy consequence of Theorem 2.4 and Corollary 2.6.
###### Corollary 2.7
(a) The Hamiltonโs equation on $`A\mathrm{\Gamma }`$ generated by $`Pr_3^{}f`$, $`fI(๐ค)`$ is of the form
$`\dot{q}=\iota ^{}df(X),`$
$`\dot{\lambda }=0,`$
$`\dot{X}=[X,R(q)df(X)]+dR(q)(\lambda )(df(X)).`$
(b) For all $`f_1`$, $`f_2I(๐ค)`$, we have
$`\{Pr_3^{}f_1,Pr_3^{}f_2\}_{A\mathrm{\Gamma }}(q,\lambda ,X)`$
$`=`$ $`<dR(q)(\lambda )(df_1(X)),df_2(X)>.`$
Remark 2.8 If $`R`$ is a constant r-matrix, then it is immediate from Corollary 2.7 (b) above that functions in $`Pr_3^{}I(๐ค)`$ Poisson commute on $`A\mathrm{\Gamma }`$. In this case, the equation for $`X`$ in part (a) of the same corollary is a Lax equation in the standard r-matrix framework for Lie algebras \[STS1\]. So when $`R`$ is constant, what we have here is a slight extension of the standard framework. For an example associated with a constant r-matrix which fits into our framework, we refer the reader to Section 5.
### 3. Realization of Hamiltonian systems in coboundary dynamical Lie faak algebroids
Let $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$ be the coboundary dynamical Lie algebroid corresponding to $`R`$, and let $`\rho :XA\mathrm{\Gamma }`$ be a realization of a Poisson manifold $`(X,\{,\}_X)`$ in the dual bundle $`A\mathrm{\Gamma }`$ of the Lie algebroid $`A^{}\mathrm{\Gamma }`$, i.e., $`\rho `$ is a Poisson map. If $`Pr_i`$ is the projection map onto the $`i`$-th factor of $`U\times ๐ฅ^{}\times ๐คA\mathrm{\Gamma }`$, $`i=1,2,3,`$ we put
$$m=Pr_1\rho :XU,$$
$`(3.1)`$
$$\tau =Pr_2\rho :X๐ฅ^{},$$
$`(3.2)`$
$$L=Pr_3\rho :X๐ค.$$
$`(3.3)`$
We shall make the following assumptions:
A1. $`X`$ is a Hamiltonian $`H`$-space with an equivariant momentum map $`J:X๐ฅ^{}`$,
A2. the realization map $`\rho `$ is $`H`$-equivariant, where $`H`$ acts on $`A\mathrm{\Gamma }`$ via the formula
$$h(q,\lambda ,X)=(Ad_{h^1}^{}q,Ad_{h^1}^{}\lambda ,Ad_hX),$$
$`(3.4)`$
A3. for some regular value $`\mu ๐ฅ^{}`$ of $`J`$,
$$\rho (J^1(\mu ))U\times \{0\}\times ๐ค.$$
$`(3.5)`$
Note that the condition in (3.4) is the natural generalization of the corresponding condition in \[LX2\] since we do not assume $`๐ฅ`$ is Abelian here. On the other hand, our assumption A3 is stronger than what we had in \[LX2\]. Our purpose in this section is to exhibit the intrinsic role played by the orbit space $`(U\times \{0\}\times ๐ค)/H`$ of the action in (3.4) in the reduction to integrable flows. We also compute how the realization map evolves under our invariant Hamiltonian systems on $`X`$.
###### Proposition 3.1
With the action defined in (3.4), the dual bundle $`A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$ equipped with the Lie-Poisson structure is a Hamiltonian $`H`$-space with equivariant momentum map $`\gamma :A\mathrm{\Gamma }๐ฅ^{}`$, $`(q,\lambda ,X)\lambda .`$
###### Demonstration Proof
Denote the action by $`\mathrm{\Phi }`$. If $`\phi C^{\mathrm{}}(A\mathrm{\Gamma })`$, it follows by a direct calculation that $`\delta _i(\phi \mathrm{\Phi }_h)(q,\lambda ,X)=Ad_{h^1}\delta _i\phi (\mathrm{\Phi }_h(q,\lambda ,X))`$, $`i=1,2`$ and $`\delta (\phi \mathrm{\Phi }_h)(q,\lambda ,X)=Ad_h^{}\delta \phi (\mathrm{\Phi }_h(q,\lambda ,X)).`$ The assertion that $`\mathrm{\Phi }_h`$ is Poisson then follows upon using the formula in Theorem 2.4 and the fact that $`R`$ is $`H`$-equivariant. Now, for any $`Z๐ฅ`$, we have
$$\frac{d}{dt}_{|_{t=0}}\mathrm{\Phi }_{e^{tZ}}(q,\lambda ,X)=(ad_Z^{}q,ad_Z^{}\lambda ,ad_ZX).$$
Comparing the right hand side of the above expression with the formula in Corollary 2.6, it is clear that this is equal to $`X_{\widehat{\gamma }(Z)}(q,\lambda ,X)`$, where $`\widehat{\gamma }(Z)(q,\lambda ,X)=<\lambda ,Z>`$. Hence $`\gamma (q,\lambda ,X)=\lambda .`$ $`\mathrm{}`$
From this result, it follows that $`X=A\mathrm{\Gamma }`$, and $`\rho =id_{A\mathrm{\Gamma }}`$ satisfy assumptions A1-A3 above with $`\mu =0`$ and we have $`\gamma ^1(0)=U\times \{0\}\times ๐ค.`$
We shall denote by $`H_\mu `$ the isotropy subgroup of $`\mu `$ for the $`H`$-action on $`X`$. Then it follows by Poisson reduction \[MR\], \[OR\] (see \[OR\] for the singular case) that the variety $`X_\mu =J^1(\mu )/H_\mu `$ inherits a unique Poisson structure $`\{,\}_{X_\mu }`$ satisfying
$$\pi _\mu ^{}\{f_1,f_2\}_{X_\mu }=i_\mu ^{}\{\stackrel{~}{f_1},\stackrel{~}{f_2}\}_X.$$
$`(3.6)`$
Here, $`i_\mu :J^1(\mu )X`$ is the inclusion map, $`\pi _\mu :J^1(\mu )X_\mu `$ is the canonical projection, $`f_1`$, $`f_2C^{\mathrm{}}(X_\mu )`$, and $`\stackrel{~}{f_1}`$, $`\stackrel{~}{f_2}`$ are (locally defined) smooth extensions of $`\pi _\mu ^{}f_1`$, $`\pi _\mu ^{}f_2`$ with differentials vanishing on the tangent spaces of the $`H`$-orbits. For the case where $`X=A\mathrm{\Gamma }`$, $`\rho =id_{A\mathrm{\Gamma }}`$, it is clear that the isotropy subgroup at $`\mu =0`$ is $`H`$ itself and so we have the Poisson variety
$$(A\mathrm{\Gamma }_0=\gamma ^1(0)/H,\{,\}_{A\mathrm{\Gamma }_0}),$$
$`(3.7)`$
with the inclusion map $`i_H:\gamma ^1(0)A\mathrm{\Gamma }`$ and the canonical projection $`\pi _H:\gamma ^1(0)A\mathrm{\Gamma }_0.`$
Clearly, functions in $`i_H^{}Pr_3^{}I(๐ค)C^{\mathrm{}}(\gamma ^1(0))`$ are $`H`$-invariant, hence they descend to functions in $`C^{\mathrm{}}(A\mathrm{\Gamma }_0)`$. On the other hand, it follows from assumption A2 that the functions in $`i_\mu ^{}L^{}I(๐ค)C^{\mathrm{}}(J^1(\mu ))`$ drop down to functions in $`C^{\mathrm{}}(X_\mu ).`$ Now, by assumption A2-A3, and the fact that $`\rho `$ is Poisson, it follows from \[OR\] that $`\rho `$ induces a unique Poisson map
$$\widehat{\rho }:X_\mu A\mathrm{\Gamma }_0=(U\times \{0\}\times ๐ค)/H$$
$`(3.8)`$
characterized by $`\pi _H\rho i_\mu =\widehat{\rho }\pi _\mu .`$ Hence $`X_\mu `$ admits a realization in the Poisson variety $`A\mathrm{\Gamma }_0`$.
We shall use the following notation. For $`fI(๐ค)`$, the unique function in $`C^{\mathrm{}}(A\mathrm{\Gamma }_0)`$ determined by $`i_H^{}Pr_3^{}f`$ will be denoted by $`\overline{f}`$; while the unique function in $`C^{\mathrm{}}(X_\mu )`$ determined by $`i_\mu ^{}L^{}f`$ will be denoted by $`_\mu `$. From the definitions, we have
$$_\mu \pi _\mu =(\widehat{\rho }^{}\overline{f})\pi _\mu =i_\mu ^{}L^{}f$$
$`(3.9)`$
###### Theorem 3.2
Let $`(X,\{,\}_X)`$ be a Poisson manifold which admits a realization $`\rho :XA\mathrm{\Gamma }`$ and assume A1-A3 are satisfied. Then there exist a unique Poisson structure $`\{,\}_{X_\mu }`$ on the reduced space $`X_\mu =J^1(\mu )/H_\mu `$ and a unique Poisson map $`\widehat{\rho }`$ such that
(a) for all $`f_1`$, $`f_2I(๐ค)`$, $`xJ^1(\mu )`$, we have
$`\{\widehat{\rho }^{}\overline{f}_1,\widehat{\rho }^{}\overline{f}_2\}_{X_\mu }\pi _\mu (x)`$
$`=`$ $`<L(x),ad_{R(m(x))df_1(L(x))}^{}df_2(L(x))+ad_{R(m(x))df_2(L(x))}^{}df_1(L(x)>.`$
(b) functions $`\widehat{\rho }^{}\overline{f}`$, $`fI(๐ค)`$, Poisson commute in $`(X_\mu ,\{,\}_{X_\mu })`$, (c) if $`\psi _t`$ is the induced flow on $`\gamma ^1(0)=U\times \{0\}\times ๐ค`$ generated by the Hamiltonian $`Pr_3^{}f`$, $`fI(๐ค)`$, and $`\varphi _t`$ is the Hamiltonian flow of $`=L^{}f`$ on $`X`$, then under the flow $`\varphi _t`$, we have
$`{\displaystyle \frac{d}{dt}}m(\varphi _t)=\iota ^{}df(L(\varphi _t)),`$
$`{\displaystyle \frac{d}{dt}}\tau (\varphi _t)=0,`$
$`{\displaystyle \frac{d}{dt}}L(\varphi _t)=[L(\varphi _t),R(m(\varphi _t))df(L(\varphi _t))]+dR(m(\varphi _t))(\tau (\varphi _t))df(L(\varphi _t))`$
where the term involving $`dR`$ drops out on $`J^1(\mu )`$. Moreover, the reduction $`\varphi _t^{red}`$ of $`\varphi _ti_\mu `$ on $`X_\mu `$ defined by $`\varphi _t^{red}\pi _\mu =\pi _\mu \varphi _ti_\mu `$ is a Hamiltonian flow of $`_\mu =\widehat{\rho }^{}\overline{f}`$ and $`\widehat{\rho }\varphi _t^{red}(\pi _\mu (x))=\pi _H\psi _t(\rho (x))`$, $`xJ^1(\mu )`$.
###### Demonstration Proof
(a) Since $`\rho (J^1(\mu ))U\times \{0\}\times ๐ค`$, we have $`\tau (x)=0`$ for $`xJ^1(\mu ).`$ Therefore,
$`\{\widehat{\rho }^{}\overline{f}_1,\widehat{\rho }^{}\overline{f}_2\}_{X_\mu }\pi _\mu (x)`$
$`=`$ $`\{\overline{f}_1,\overline{f}_2\}_{A\mathrm{\Gamma }_0}\pi _H(\rho (x))`$
$`=`$ $`\{Pr_3^{}f_1,Pr_3^{}f_2\}_{A\mathrm{\Gamma }}(\rho (x))`$
$`=`$ $`<L(x),ad_{R(m(x))df_1(L(x))}^{}df_2(L(x))+ad_{R(m(x))df_2(L(x))}^{}df_1(L(x)>`$
where in the last step we have invoked the formula in Theorem 2.4 and the vanishing of $`\tau (x)`$ for $`xJ^1(\mu )`$.
(b) This is clear from part (a).
(c) Since $`\rho `$ is a Poisson map, we have $`\frac{d}{dt}\rho (\varphi _t)=X_{fPr_3}(\rho (\varphi _t))`$ from which the equations follow on invoking Corollary 2.7. Finally, the assertion on $`\varphi _t^{red}`$ is basically a corollary of Theorem 2.16 of \[OR\] and the relation $`\rho \varphi _ti_\mu =\psi _t\rho i_\mu `$. $`\mathrm{}`$
Remark 3.3 In \[LX2\], we have only written down the equation for $`L`$ under the Hamiltonian flow $`\varphi _t`$ (in the Abelian case). However, the full set of equations is important. See Section 4 and Section 6 below.
### 4. Factorization problems on Lie groupoids and exact solvability
We shall develop a factorization method to solve the (generalized) Lax equations in Corollary 2.7 (a) on the level set $`\gamma ^1(0)`$ of the momentum map $`\gamma `$. For the first part of this section, we shall use $`A\mathrm{\Gamma }=_{qU}\{0_q\}\times ๐ค\times ๐ฅ^{}`$, $`A^{}\mathrm{\Gamma }=_{qU}\{0_q\}\times ๐ค^{}\times ๐ฅ`$ , and when $`๐ค`$ has an ad-invariant non-degenerate pairing, we shall identify the Lie algebras with their duals.
As in \[L1\], we introduce the bundle map
$$:A^{}\mathrm{\Gamma }A\mathrm{\Gamma },(0_q,A,Z)(0_q,\iota Z+R(q)A,\iota ^{}Aad_Z^{}q)$$
$`(4.1)`$
and call it the r-matrix of the Lie algebroid $`A^{}\mathrm{\Gamma }`$. Also, we assume $`R`$ satisfies the modified dynamical Yang-Baxter equation (mDYBE):
$$\begin{array}{cc}& [R(q)A,R(q)B]+R(q)(ad_{R(q)A}^{}Bad_{R(q)B}^{}A)\hfill \\ \hfill +& dR(q)\iota ^{}A(B)dR(q)\iota ^{}B(A)+d<R(A),B>(q)\hfill \\ \hfill =& [K(A),K(B)]\hfill \end{array}$$
$`(4.2)`$
where $`KL(๐ค^{},๐ค)`$ is a nonzero symmetric map which satisfies $`ad_XK+Kad_X^{}=0`$ for all $`X๐ค`$,i.e, $`K`$ is $`G`$-equivariant.
The next two results were announced in \[L1\]. We give details of the proof here.
###### Lemma 4.1
If R satisfies (mDYBE), then the r-matrix $`:A^{}\mathrm{\Gamma }A\mathrm{\Gamma }`$ satisfies the equation
$$\begin{array}{cc}& [(0,A,Z),(0,A^{},Z^{})]_{A\mathrm{\Gamma }}[(0,A,Z),(0,A^{},Z^{})]_{A^{}\mathrm{\Gamma }}\hfill \\ \hfill =& (0,[K(A),K(A^{})],0)\hfill \end{array}$$
$`(4.3)`$
for all smooth maps $`A,A^{}:U๐ค^{}`$, $`Z,Z^{}:U๐ฅ`$.
###### Demonstration Proof
The calculation will be postponed to the appendix. $`\mathrm{}`$
Using $`K`$, we define
$$๐ฆ:A^{}\mathrm{\Gamma }A\mathrm{\Gamma },(0_q,A,Z)(0_q,K(A),0),$$
$`(4.4)`$
and set $`^\pm =\pm ๐ฆ,R^\pm (q)=R(q)\pm K.`$
###### Proposition 4.2
(a) $`^\pm `$ are morphisms of transitive Lie algebroids and, as morphisms of vector bundles over $`U`$, are of locally constant rank. In particular,
$$[^\pm (0,A,Z),^\pm (0,A^{},Z^{})]_{A\mathrm{\Gamma }}=^\pm [(0,A,Z),(0,A^{},Z^{})]_{A^{}\mathrm{\Gamma }}$$
$`(4.5)`$
for all smooth maps $`A,A^{}:U๐ค^{}`$, $`Z,Z^{}:U๐ฅ.`$ Moreover, $`^\pm `$ are $`H`$-equivariant, where $`H`$ acts on $`A^{}\mathrm{\Gamma }`$ via $`h(0_q,A,Z)=(0_{Ad_{h^1}^{}}q,Ad_{h^1}^{}A,Ad_hZ)`$ and the $`H`$-action on $`A\mathrm{\Gamma }`$ is given by (3.4). (b) $`Im^\pm `$ are transitive Lie subalgebroids of $`A\mathrm{\Gamma }.`$
###### Demonstration Proof
(a) If $`a`$ and $`a_{}`$ are the anchor maps of the Lie algebroids $`A\mathrm{\Gamma }`$ and $`A^{}\mathrm{\Gamma }`$ , it is easy to check that they are surjective submersions which satisfy $`a^\pm =a_{}`$. On the other hand, it follows from Lemma 4.1 that (4.5) holds if and only if
$$\begin{array}{cc}& ๐ฆ[(0,A,Z),(0,A^{},Z^{})]_{A^{}\mathrm{\Gamma }}\hfill \\ \hfill =& [(0,A,Z),๐ฆ(0,A^{},Z^{})]_{A\mathrm{\Gamma }}+[๐ฆ(0,A,Z),(0,A^{},Z^{})]_{A\mathrm{\Gamma }}.\hfill \end{array}$$
$`(4.6)`$
Now, for $`qU`$, we have
$`๐ฆ[(0,A,Z),(0,A^{},Z^{})]_{A^{}\mathrm{\Gamma }}(q)`$
$`=`$ $`(0_q,K(dA^{}(q)(\iota ^{}A(q)ad_{Z(q)}^{}q)ad_{R(q)A(q)Z(q)}^{}A^{}(q))`$
$`(AA^{},ZZ^{}),0)`$
where $`(AA^{},ZZ^{})`$ denote terms which can be obtained from the previous ones by interchanging $`A`$ and $`A^{}`$, $`Z`$ and $`Z^{}`$. On the other hand,
$`[(0,A,Z),๐ฆ(0,A^{},Z^{})]_{A\mathrm{\Gamma }}(q)`$
$`=`$ $`(0_q,K(dA^{}(q)(\iota ^{}A(q)ad_{Z(q)}^{}q))+[\iota Z(q)+R(q)A(q),K(A^{}(q))],0)`$
and similarly for $`[(0,A^{},Z^{}),๐ฆ(0,A,Z)]_{A\mathrm{\Gamma }}(q).`$ From these formulas, it follows that (4.6) holds if and only if
$`K(ad_{R(q)A(q)Z(q)}^{}A^{}(q))+[\iota Z(q)+R(q)A(q),K(A^{}(q))]`$
$`(AA^{},ZZ^{})=0.`$
But the latter follows from the $`G`$-equivariance of $`K`$ and this proves the first part of the assertion.(The fact that $`^\pm `$ are of locally constant rank follows from Theorem 1.6 on page 190 of \[M\].) To show that $`^\pm `$ are $`H`$-equivariant, note that by definition,
$`^\pm (h(0_q,A,Z))`$
$`=`$ $`(0_{Ad_{h^1}^{}q},Ad_hZ+R^\pm (Ad_{h^1}^{}q)Ad_{h^1}^{}A,\iota ^{}Ad_{h^1}^{}Aad_{Ad_hZ}^{}Ad_{h^1}^{}q).`$
But from the $`H`$-equivariance of $`R`$ and $`K`$, we have $`R^\pm (Ad_{h^1}^{}q)Ad_{h^1}^{}A=Ad_hR^\pm (q)A.`$ On the other hand, it is straightforward to check that $`ad_{Ad_hZ}^{}Ad_{h^1}^{}q=Ad_{h^1}^{}ad_Z^{}q.`$ Substituting into the above expression for $`^\pm (h(0_q,A,Z))`$, the desired conclusion follows.
(b) This is a consequence of (a). $`\mathrm{}`$
In the rest of the section, we shall assume $`๐ค`$ has an ad-invariant non-degenerate pairing $`(,)`$ such that $`(,)|_{๐ฅ\times ๐ฅ}`$ is also non-degenerate. Without loss of generality, we shall take the map $`K:๐ค^{}๐ค`$ in the above discussion to be the identification map induced by $`(,)`$. Indeed, with the identifications $`๐ค^{}๐ค`$, $`๐ฅ^{}๐ฅ`$, we shall regard $`R(q)`$ as taking values in $`End(๐ค)`$, and the left and right gradients as well as the dual maps are computed using $`(,)`$. Also, we have $`ad^{}ad`$, $`\iota ^{}\mathrm{\Pi }_๐ฅ`$, where $`\mathrm{\Pi }_๐ฅ`$ is the projection map to $`๐ฅ`$ relative to the direct sum decomposition $`๐ค=๐ฅ๐ฅ^{}`$. We shall keep, however, the notation $`A^{}\mathrm{\Gamma }`$ although as a set it can be identified with $`A\mathrm{\Gamma }.`$
We now introduce the following subbundles of the adjoint bundle $`Kera=\{(0_q,X,0)qU,X๐ค\}`$ of $`A\mathrm{\Gamma }`$:
$$^+=\{(0_q,X,0)KeraqU,^{}(0_q,X,Z)=0\text{for some}Z๐ฅ\},$$
$`(4.7a)`$
$$^{}=\{(0_q,X,0)KeraqU,^+(0_q,X,Z)=0\text{for some}Z๐ฅ\}.$$
$`(4.7b)`$
###### Proposition 4.3
$`^\pm `$ are ideals of the transitive Lie algebroids $`Im^\pm `$.
###### Demonstration Proof
We shall prove the assertion for $`^+.`$ First of all, it is easy to show that $`^+Im^+.`$ Let $`(0,\iota Z+R^+X,\mathrm{\Pi }_๐ฅX+ad_Z())Sect(U,Im^+)`$ and $`(0,X^{},0)Sect(U,^+)`$, where $`Z:U๐ฅ`$, $`X,X^{}:U๐ค`$ are smooth maps. From the expression for $`[,]_{A\mathrm{\Gamma }}`$, we have
$$\begin{array}{cc}& [(0,\iota Z+R^+X,\mathrm{\Pi }_๐ฅX+ad_Z()),(0,X^{},0)]_{A\mathrm{\Gamma }}(q)\hfill \\ \hfill =& (0_q,dX^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)+[\iota Z(q)+R^+(q)X(q),X^{}(q)],0)\hfill \end{array}$$
$`(4.8)`$
for $`qU`$. This shows
$$[(0,\iota Z+R^+X,\mathrm{\Pi }_๐ฅX+ad_Z()),(0,X^{},0)]_{A\mathrm{\Gamma }}Sect(U,Kera).$$
On the other hand, from the assumption that $`(0,X^{},0)Sect(U,^+)`$, we must have $`^{}(0_q,X^{}(q),Z^{}(q))=0`$ for some smooth map $`Z^{}:U๐ฅ.`$ Hence we obtain $`(0,X^{},0)=^+(0,X^{}/2,Z^{}/2).`$ Therefore, on using Proposition 4.2 (a), it follows that
$$\begin{array}{cc}& [(0,\iota Z+R^+X,\mathrm{\Pi }_๐ฅX+ad_Z()),(0,X^{},0)]_{A\mathrm{\Gamma }}(q)\hfill \\ \hfill =& ^+[(0,X,Z),(0,X^{}/2,Z^{}/2)]_{A^{}\mathrm{\Gamma }}(q)\hfill \\ \hfill =& ^+(0_q,X^{\prime \prime }(q),Z^{\prime \prime }(q))\hfill \\ \hfill =& (0_q,\iota Z^{\prime \prime }(q)+R^+(q)X^{\prime \prime }(q),\mathrm{\Pi }_๐ฅX^{\prime \prime }(q)+[Z^{\prime \prime }(q),q])\hfill \end{array}$$
$`(4.9)`$
where by (2.7), we find
$$\begin{array}{cc}\hfill X^{\prime \prime }(q)=& \frac{1}{2}dX^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)+\frac{1}{2}[R(q)X(q)Z(q),X^{}(q)]\hfill \\ & +\frac{1}{2}[X(q),R(q)X^{}(q)Z^{}(q)]\hfill \end{array}$$
$`(4.10)`$
and
$$\begin{array}{cc}\hfill Z^{\prime \prime }(q)=& \frac{1}{2}dZ^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)\frac{1}{2}[Z,Z^{}](q)\hfill \\ & +\frac{1}{2}(dR(q)()X(q),X^{}(q)).\hfill \end{array}$$
$`(4.11)`$
By equating the last components of the expressions in (4.8) and (4.9), we have
$$\mathrm{\Pi }_๐ฅX^{\prime \prime }(q)+[Z^{\prime \prime }(q),q]=0.$$
$`(4.12)`$
Similarly, by equating the second components of the expressions in (4.8) and (4.9), we find
$$\begin{array}{cc}& dX^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)+[\iota Z(q)+R^+(q)X(q),X^{}(q)]\hfill \\ \hfill =& \iota Z^{\prime \prime }(q)+R^+(q)X^{\prime \prime }(q)\hfill \end{array}$$
$`(4.13)`$
Now, from the relation $`^{}(0_q,X^{}(q),Z^{}(q))=0`$, it follows that (4.10) can be rewritten as
$$\begin{array}{cc}& dX^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)+[\iota Z(q)+R^+(q)X(q),X^{}(q)]\hfill \\ \hfill =& 2X^{\prime \prime }(q).\hfill \end{array}$$
$`(4.14)`$
Substitute this into (4.12), we obtain
$$\begin{array}{cc}\hfill \mathrm{\Pi }_๐ฅ& \{dX^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)+[\iota Z(q)+R^+(q)X(q),X^{}(q)]\}\hfill \\ & +[2Z^{\prime \prime }(q),q]=0.\hfill \end{array}$$
Next, (4.13) and (4.14) yield
$`R^{}(q)\{dX^{}(q)(\mathrm{\Pi }_๐ฅX(q)+ad_{Z(q)}q)+[\iota Z(q)+R^+(q)X(q),X^{}(q)]\}`$
$`=`$ $`2Z^{\prime \prime }(q).`$
From the last two relations, we can now conclude that
$$[(0,\iota Z+R^+X,\mathrm{\Pi }_๐ฅX+ad_Z(),(0,X^{},0)]_{A\mathrm{\Gamma }}Sect(U,^+),$$
as desired. $`\mathrm{}`$
Consider now the quotient vector bundles $`Im^\pm /^\pm `$ equipped with the quotient transitive Lie algebroid structures.
###### Proposition 4.4
The map $`\theta :Im^+/^+Im^{}/^{}`$ defined by
$$\theta (^+(0_q,X,Z)+_q^+)=^{}(0_q,X,Z)+_q^{}$$
is an isomorphism of transitive Lie algebroids.
###### Demonstration Proof
We first show that $`\theta `$ is well-defined. To do so, suppose $`^+(0_q,X,Z)^+(0_q,X^{},Z^{})`$ (mod $`_q^+`$). Then there exists $`Z^{\prime \prime }๐ฅ`$ such that
$$\mathrm{\Pi }_๐ฅ(XX^{})+[ZZ^{},q]=0,$$
$$R^{}(q)(\iota (ZZ^{})+R^+(q)(XX^{}))=\iota Z^{\prime \prime },$$
and
$$(ZZ^{})+\mathrm{\Pi }_๐ฅ(R^+(q)(XX^{}))=ad_qZ^{\prime \prime }.$$
From these equations, we infer that
$$\iota (ZZ^{})+\mathrm{\Pi }_๐ฅ(R^{}(q)(XX^{}))=ad_q(Z^{\prime \prime }2(ZZ^{})),$$
and
$$R^+(q)(\iota (ZZ^{})+R^{}(q)(XX^{}))=Z^{\prime \prime }2(ZZ^{}).$$
Hence we have $`^{}(0_q,X,Z)^{}(0_q,X^{},Z^{})`$ (mod $`_q^{}`$), as desired. We shall skip the argument to show that $`\theta `$ is 1:1 as it is similar to the one above. The proof of the proposition is therefore complete (it being clear that $`\theta `$ is a morphism by Proposition 4.2 (a)). $`\mathrm{}`$
To formulate our next result,introduce the Lie algebroid direct sum $`A\mathrm{\Gamma }\underset{TU}{}A\mathrm{\Gamma }.`$ Clearly, this is the Lie algebroid of the product groupoid $`P=\mathrm{\Gamma }\underset{U\times U}{\times }\mathrm{\Gamma }U`$ ($``$ the trivial Lie groupoid $`U\times (G\times G)\times U`$). For later usage, we shall denote the structure maps (target, source etc.) of $`P`$ by $`\alpha _P`$, $`\beta _P`$, and so forth.
###### Theorem 4.5
(a)The map $`(^+,^{}):A^{}\mathrm{\Gamma }A\mathrm{\Gamma }\underset{TU}{}A\mathrm{\Gamma }`$ is a monomorphism of transitive Lie algebroids. In particular, the coboundary dynamical Lie algebroid $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }})`$ is integrable. (b) $`Im(^+,^{})`$ is the Lie subalgebroid
$$\{(๐ณ_+,๐ณ_{})(Im^+\underset{TU}{}Im^{})_qqU,\theta (๐ณ_++_q^+)=๐ณ_{}+_q^{}\}$$
$`(4.15)`$
of $`Im^+\underset{TU}{}Im^{}`$.
###### Demonstration Proof
(a) See \[L1\] for the proof.
(b) Denote by $`A\mathrm{\Gamma }_R`$ the subbundle of $`Im^+\underset{TU}{}Im^{}`$ defined in (4.15). It is clear that $`Im(^+,^{})A\mathrm{\Gamma }_R.`$ Conversely, suppose $`((0_q,X_+,Z),(0_q,X_{},Z))A\mathrm{\Gamma }_R`$. Then there exist $`(0_q,X,\stackrel{~}{Z})`$, $`(0_q,X^{},\stackrel{~}{Z}^{})A\mathrm{\Gamma }`$ such that $`(0_q,X_+,Z)=^+(0_q,X,\stackrel{~}{Z})`$ and $`(0_q,X_{},Z)=^{}(0_q,X^{},\stackrel{~}{Z}^{}).`$ Moreover, from the property that $`\theta ((0_q,X_+,Z)+_q^+)=(0_q,X_{},Z)+_q^{},`$ we find $`^{}(0_q,XX^{},\stackrel{~}{Z}\stackrel{~}{Z}^{})0`$ (mod $`_q^{}`$). Let $`X^{\prime \prime }=\iota (\stackrel{~}{Z}\stackrel{~}{Z}^{})+R^{}(q)(XX^{})`$. Then it follows from the definition of $`_q^{}`$ that there exists $`Z^{\prime \prime }๐ฅ`$ such that $`^+(0_q,X^{\prime \prime },Z^{\prime \prime })=0.`$ Now, consider the element $`(0_q,X+\frac{1}{2}X^{\prime \prime },\stackrel{~}{Z}+\frac{1}{2}Z^{\prime \prime })A\mathrm{\Gamma }.`$ Clearly, $`^+(0_q,X+\frac{1}{2}X^{\prime \prime },\stackrel{~}{Z}+\frac{1}{2}Z^{\prime \prime })=(0_q,X_+,Z)`$. On the other hand,
$$\begin{array}{cc}& ^{}(0_q,X+\frac{1}{2}X^{\prime \prime },\stackrel{~}{Z}+\frac{1}{2}Z^{\prime \prime })\hfill \\ \hfill =& (0_q,X_{},Z)+(0_q,X^{\prime \prime },0)+\frac{1}{2}^{}(0_q,X^{\prime \prime },Z^{\prime \prime }).\hfill \end{array}$$
But as
$`^{}(0_q,X^{\prime \prime },Z^{\prime \prime })`$
$`=`$ $`^+(0_q,X^{\prime \prime },Z^{\prime \prime })(0_q,2X^{\prime \prime },0)`$
$`=`$ $`(0_q,2X^{\prime \prime },0),`$
it follows from the above that $`^{}(0_q,X+\frac{1}{2}X^{\prime \prime },\stackrel{~}{Z}+\frac{1}{2}Z^{\prime \prime })=(0_q,X_{},Z).`$ Thus we have shown that $`((0_q,X_+,Z),(0_q,X_{},Z))Im(^+,^{}).`$ $`\mathrm{}`$
The connection between (mDYBE) and our factorization theory is contained in the decomposition
$$\begin{array}{cc}\hfill (0_q,X,0)& =\frac{1}{2}^+(0_q,X,0)\hfill \\ & \frac{1}{2}^{}(0_q,X,0)\hfill \end{array}$$
$`(4.16)`$
where the element $`(0_q,X,0)`$ on the left hand side of (4.16) is in the adjoint bundle $`Kera`$ of $`A\mathrm{\Gamma }.`$ The reader should note that the vector bundles $`\{^\pm (0_q,X,0)qU,X๐ค\}`$ are not Lie subalgebroids of $`A\mathrm{\Gamma }`$ unless $`R`$ is a constant r-matrix. As we pointed out in \[L1\], this fact has repercussion when we try to formulate a global version of the decomposition in (4.16) (see Corollary 4.6 below).
In the rest of the section, we shall assume both $`G`$ and $`U`$ are simply-connected. Let $`\mathrm{\Gamma }^{}`$ be the unique source-simply connected Lie groupoid which integrates $`(A^{}\mathrm{\Gamma },`$ $`[,]_{A^{}\mathrm{\Gamma }})`$. Then $`(^+,^{})`$ can be lifted up to a unique monomorphism of Lie groupoids $`\mathrm{\Gamma }^{}\mathrm{\Gamma }\underset{U\times U}{\times }\mathrm{\Gamma }`$ which we shall denote by the same symbol. Now, denote by $`\mathrm{\Gamma }=\{(u,g,u)uU,gG\}`$ the gauge group bundle of $`\mathrm{\Gamma }`$. We let $`j:\mathrm{\Gamma }\underset{U\times U}{\times }\mathrm{\Gamma }\mathrm{\Gamma }`$ be the map defined by $`j(a,b)=ab^1`$ and let $`\stackrel{~}{m}=j(^+,^{})`$.
For the sake of completeness, we include the following Corollary of Theorem 4.5 (a) which (essentially) gives a global version of the decomposition in (4.16) which we mentioned above (the reader can find the proof in \[L1\]). For its formulation, note that the Lie groupoid of $`\{(0_q,0,Z)qU,Z๐ฅ\}A^{}\mathrm{\Gamma }`$ is $`H\times U`$, with target and source maps $`\alpha ^{}(h,u)=u`$, $`\beta ^{}(h,u)=Ad_hu`$ and multiplication map $`m^{}((h,u),(k,Ad_hu))=(kh,u)`$ (this is isomorphic to the Hamiltonian unit in \[LP2\]). On the other hand, the Lie groupoid of $`^\pm \{(0_q,0,Z)qU,Z๐ฅ\}`$ is given by $`E=\{(u,h,Ad_{h^1}u)uU,hH\}`$ and $`^\pm `$ embeds $`H\times U`$ in $`E`$, $`^\pm H\times U:(h,u)(u,h^1,Ad_hu).`$ Clearly, the diagonal $`\mathrm{\Delta }(E)`$ of $`E\underset{U\times U}{\times }E`$ acts on $`Im(^+,^{})`$ from the right via the simple formula
$`((u,k_+,v),(u,k_{},v)).((v,h,Ad_{h^1}v),(v,h,Ad_{h^1}v))`$
$`=`$ $`((u,k_+h,Ad_{h^1}v),(u,k_{}h,Ad_{h^1}v))`$
and the map $`jIm(^+,^{})`$ is constant on the orbits of this action.
###### Corollary 4.6
Suppose $`U`$ is simply-connected, then $`jIm(^+,^{})`$ induces a one-to-one map $`\widehat{j}:Im(^+,^{})/\mathrm{\Delta }(E)\mathrm{\Gamma }`$. Therefore, for each $`\gamma Im\stackrel{~}{m}`$, there exists unique $`[(\gamma _+,\gamma _{})]`$ in the homogeneous space $`Im(^+,^{})/\mathrm{\Delta }(E)`$ such that $`\widehat{j}([(\gamma _+,\gamma _{})])=\gamma `$.
Let $`fI(๐ค)`$ and consider the Hamiltonโs equation generated by $`F=Pr_3^{}f`$. Then according to Corollary 2.7 (a), we can express its restriction to the invariant manifold $`\gamma ^1(0)=U\times \{0\}\times ๐ค`$ in the form
$$\begin{array}{cc}& \frac{d}{dt}(q,0,X)\hfill \\ \hfill =& (\mathrm{\Pi }_๐ฅdf(X),0,[X,R(q)df(X)]).\hfill \end{array}$$
$`(4.17)`$
In the next theorem, we shall express the solution of (4.17) using the adjoint representation of $`\mathrm{\Gamma }`$ on its adjoint bundle $`Kera`$, defined by $`๐ธ๐_\gamma (q,0,X)=(q^{},0,Ad_kX)`$, for $`\gamma =(q^{},k,q)\mathrm{\Gamma }`$. We shall also make the identifications $`A\mathrm{\Gamma }`$, $`A^{}\mathrm{\Gamma }U\times ๐ฅ\times ๐ค`$ throughout. Thus the element $`(0,0,df(X_0))`$ which appears in the theorem below will denote the constant section of $`Kera`$ such that $`(0,0,df(X_0))(q)=(q,0,df(X_0))`$ for $`qU`$.
###### Theorem 4.7
Suppose that $`fI(๐ค)`$, $`F=Pr_3^{}f`$ and $`q_0U`$, where $`U`$ is simply connected. Then for some $`0<T\mathrm{}`$, there exists a unique element $`(\gamma _+(t),\gamma _{}(t))=((q_0,k_+(t),q(t)),(q_0,k_{}(t),q(t)))Im(^+,^{})`$ for $`0t<T`$ which is smooth in t, solves the factorization problem
$$exp\{2t(0,0,df(X_0))\}(q_0)=\gamma _+(t)\gamma _{}(t)^1$$
$`(4.18)`$
and satisfies
$$\begin{array}{cc}\hfill (T_{\gamma _+(t)}๐_{\gamma _+(t)^1}\dot{\gamma }_+(t),T_{\gamma _{}(t)}๐_{\gamma _{}(t)^1}\dot{\gamma }_{}(t))& (^+,^{})(\{q(t)\}\times \{0\}\times ๐ค)\hfill \end{array}$$
$`(4.19a)`$
with
$$\gamma _\pm (0)=(q_0,1,q_0).$$
$`(4.19b)`$
Moreover, the solution of (4.17) with initial data $`(q,0,X)(0)=(q_0,0,X_0)`$ (i.e. the induced flow on $`\gamma ^1(0)`$ generated by $`F`$) is given by the formula
$$(q(t),0,X(t))=๐ธ๐_{\gamma _\pm (t)^1}(q_0,0,X_0).$$
$`(4.20)`$
###### Demonstration Proof
The uniqueness of the element $`(\gamma _+(t),\gamma _{}(t))`$ is proved in the same way as in \[L1\] and makes crucial use of Corollary 4.6.
Assuming the existence of the factors for the moment, we claim that $`(q(t),0,X(t))`$ as given by (4.20) solves (4.17). First of all, we have
$`๐ธ๐_{\gamma _+(t)^1}(q_0,0,X_0)`$
$`=`$ $`(q(t),0,Ad_{k_+(t)^1}X_0)`$
$`=`$ $`(q(t),0,Ad_{k_{}(t)^1}Ad_{e^{2tdf(X_0)}}X_0)`$
$`=`$ $`(q(t),0,Ad_{k_{}(t)^1}X_0)`$
$`=`$ $`๐ธ๐_{\gamma _{}(t)^1}(q_0,0,X_0)`$
where we have used the fact that $`[df(X_0),X_0]=0`$. Take
$$(q(t),0,X(t))=๐ธ๐_{\gamma _+(t)^1}(q_0,0,X_0).$$
By differentiating the expression, we have
$`{\displaystyle \frac{d}{dt}}(q(t),0,X(t))`$
$`=`$ $`(\dot{q}(t),0,[X(t),T_{k_+(t)}l_{k_+(t)^1}\dot{k}_+(t)]).()`$
On the other hand, by rewriting (4.18) in the form
$$exp\{2t(0,df(X_0),0)\}(q_0)\gamma _{}(t)=\gamma _+(t),$$
we have, upon differentiation, that
$`T_{\gamma _+(t)}๐_{\gamma _+(t)^1}\dot{\gamma }_+(t)T_{\gamma _{}(t)}๐_{\gamma _{}(t)^1}\dot{\gamma }_{}(t)`$
$`=`$ $`2๐ธ๐_{\gamma _{}(t)^1}(q_0,0,df(X_0)).`$
But
$`๐ธ๐_{\gamma _{}(t)^1}(q_0,0,df(X_0))`$
$`=`$ $`(q(t),0,Ad_{k_{}(t)^1}df(X_0))`$
$`=`$ $`(q(t),0,df(X(t)))`$
as $`fI(๐ค)`$. Hence it follows that
$`T_{\gamma _+(t)}๐_{\gamma _+(t)^1}\dot{\gamma }_+(t)T_{\gamma _{}(t)}๐_{\gamma _{}(t)^1}\dot{\gamma }_{}(t)`$
$`=`$ $`2(q(t),0,df(X(t))).`$
From the property of $`\gamma _\pm `$ in (4.19), we can now conclude that
$$T_{\gamma _\pm (t)}๐_{\gamma _\pm (t)^1}\dot{\gamma }_\pm (t)=^\pm (q(t),0,df(X(t))).$$
But
$$T_{\gamma _+(t)}๐_{\gamma _+(t)^1}\dot{\gamma }_+(t)=(q(t),\dot{q}(t),T_{k_+(t)}l_{k_+(t)^1}\dot{k}_+(t)),$$
while
$$^+(q(t),0,df(X(t)))=(q(t),\mathrm{\Pi }_๐ฅdf(X(t)),R^+(q(t))df(X(t))).$$
By equating the two expressions, we obtain
$$\dot{q}(t)=\mathrm{\Pi }_๐ฅdf(X(t)),$$
and
$$T_{k_+(t)}l_{k_+(t)^1}\dot{k}_+(t)=R^+(q(t))df(X(t)).$$
Therefore, on substituting into (\*), we find
$`{\displaystyle \frac{d}{dt}}(q(t),0,X(t))`$
$`=`$ $`(\mathrm{\Pi }_๐ฅdf(X(t)),[X(t),R(q(t))df(X(t))]),`$
as claimed.
To prove the existence of the factors $`\gamma _\pm (t)`$, simply solve the initial value problems
$$\dot{k}_\pm (t)=T_el_{k_\pm (t)}R^\pm (q(t))df(X(t)),k_\pm (0)=1,()$$
where $`q(t)`$, $`X(t)`$ are the solutions of (4.17) with initial data $`(q,0,X)(0)=(q_0,0,X_0)`$ (which are known to exist by ODE theory). Set $`\gamma _\pm (t)=(q_0,k_\pm (t),q(t))`$. As can be easily verified, we can combine the equations for $`q(t)`$, $`k_\pm (t)`$ into one single equation for $`(\gamma _+(t),\gamma _{}(t))`$:
$`{\displaystyle \frac{d}{dt}}(\gamma _+(t),\gamma _{}(t))`$
$`=`$ $`(T_{ฯต(q(t))}๐_{\gamma _+(t)}^+(q(t),0,df(X(t))),T_{ฯต(q(t))}๐_{\gamma _{}(t)}^{}(q(t),0,df(X(t))))`$
$`=`$ $`T_{ฯต_P(\beta _P(\gamma _+(t),\gamma _{}(t)))}๐_{(\gamma _+(t),\gamma _{}(t))}^P(^+,^{})(q(t),0,df(X(t)))()`$
where $`๐_{(\gamma _+(t),\gamma _{}(t))}^P`$ represents left translation by $`(\gamma _+(t),\gamma _{}(t))`$ in the product groupoid $`P=\mathrm{\Gamma }\underset{U\times U}{\times }\mathrm{\Gamma }U`$. Clearly, what we have just written down is a well-defined equation for $`(\gamma _+(t),\gamma _{}(t))Im(^+,^{}).`$ Moreover, from the initial conditions for $`k_\pm (t)`$ and $`q(t)`$, we have $`(\gamma _+(0),\gamma _{}(0))Im(^+,^{}).`$
Now, from the equations for $`k_\pm `$ in (\**), we find
$`T_{\gamma _+(t)\gamma _{}(t)^1}๐_{(\gamma _+(t)\gamma _{}(t)^1)^1}{\displaystyle \frac{d}{dt}}\gamma _+(t)\gamma _{}(t)^1`$
$`=`$ $`(q_0,0,2df(Ad_{k_{}(t)}X(t))).`$
But from the equation for $`k_{}(t)`$ and $`X(t)`$ , we have
$`{\displaystyle \frac{d}{dt}}Ad_{k_{}(t)}X(t)`$
$`=`$ $`Ad_{k_{}(t)}\dot{X}(t)+[T_{k_{}(t)}r_{k_{}(t)^1}\dot{k}_{}(t),Ad_{k_{}(t)}X(t)]`$
$`=`$ $`[Ad_{k_{}(t)}X(t),Ad_{k_{}(t)}R(q(t))df(X(t))]`$
$`+[Ad_{k_{}(t)}R^{}(q(t))df(X(t)),Ad_{k_{}(t)}X(t)]`$
$`=`$ $`0.`$
Therefore, $`Ad_{k_{}(t)}X(t)=X_0`$ and so
$`T_{\gamma _+(t)\gamma _{}(t)^1}๐_{(\gamma _+(t)\gamma _{}(t)^1)^1}{\displaystyle \frac{d}{dt}}\gamma _+(t)\gamma _{}(t)^1`$
$`=`$ $`(q_0,0,2df(X_0)).`$
As $`\gamma _+(t)\gamma _{}(t)^1=(q_0,k_+(t)k_{}(t)^1,q_0),`$ this shows that $`k_+(t)k_{}(t)^1=e^{2tdf(X_0)}`$ and consequently,
$$exp\{2t(0,df(X_0),0)\}(q_0)=\gamma _+(t)\gamma _{}(t)^1.$$
Thus it remains to show that condition (4.19 a) is satisfied. But this is immediate from (\***). This completes the proof. $`\mathrm{}`$
###### Corollary 4.8
Let $`\psi _t`$ be the induced flow on $`\gamma ^1(0)=U\times \{0\}\times ๐ค`$ as defined in (4.20) and let $`\varphi _t`$ be the Hamiltonian flow of $`=L^{}f`$ on $`X`$, where $`L=Pr_3\rho `$ for a realization map $`\rho :XA\mathrm{\Gamma }`$ satisfying A1-A3. If we can solve for $`\varphi _t(x)`$, $`xJ^1(\mu )`$ explicitly from the relation $`\rho (\varphi _t)(x)=\psi _t(\rho (x))`$, then the formula $`\varphi _t^{red}\pi _\mu =\pi _\mu \varphi _ti_\mu `$ gives an explicit expression for the flow of the reduced Hamiltonian $`_\mu =\widehat{\rho }^{}\overline{f}`$.
Remark 4.9 (a) The reader should not feel uneasy about the use of the equation (\**) above (which involve the solutions $`q(t)`$ and $`X(t)`$) to show the existence of the factors $`k_\pm (t)`$, and which are then used in turn to construct $`q(t)`$ and $`X(t)`$. As the reader will see in Section 6 below, knowledge of the existence of the factorization facilitates its construction. (b) If we take $`K=\frac{1}{2}id_๐ค`$, which is what we will need in Section 6, then the factorization problem in (4.18) has to be replaced by $`exp\{t(0,0,df(X_0))\}(q_0)=\gamma _+(t)\gamma _{}(t)^1`$. Otherwise, the solution formula is the same as before. (c) There is a similar method for solving the Hamiltonian flows generated by natural invariant functions on the gauge group bundles of cobundary dynamical Poisson groupoids. We shall refer the reader to \[L1\] for details. (d) For the hyperbolic spin Calogero-Moser systems and the spin Toda lattices which we introduce in the next section, the assumption in Corollary 4.8 (namely, we can solve for $`\varphi _t(x)`$, $`xJ^1(\mu )`$ explicitly from the relation $`\rho (\varphi _t)(x)=\psi _t(\rho (x))`$) are actually not satisfied in general. As the reader will see, some special structure of these equations still enables us to obtain the Hamiltonian flows on $`J^1(\mu )`$ from the induced flows on $`\gamma ^1(0)`$. (e) Clearly, Theorem 4.7 also applies in the case when $`R`$ is a constant r-matrix. However, it is possible to formulate an analog of this result using the fact that the vector bundles $`\{^\pm (0_q,X,0)qU,X๐ค\}`$ are Lie subalgebroids of $`A\mathrm{\Gamma }`$ in this case, but we provide no details here.
### 5. A family of hyperbolic spin Calogero-Moser systems and the spin fak Toda lattices
In \[EV\], the authors classified solutions of (mDYBE) for pairs $`(๐ค,๐ฅ)`$ of Lie algebras, where $`๐ค`$ is simple, and $`๐ฅ๐ค`$ is a Cartan subalgebra. The purpose of this section is to introduce a natural family of hyperbolic spin Calogero-Moser systems associated with these solutions as another application of Proposition 4.2 (a). Remarkably, these models admit scaling limits, and the result is a family of Hamiltonian systems which may be regarded as a spin generalization of the Toda lattice.
Let us begin with some notation. Let $`๐ค=๐ฅ_{\alpha \mathrm{\Delta }}๐ค_\alpha `$ the root space decomposition of the simple Lie algebra $`๐ค`$ and let $`(,)`$ denote its Killing form. For each $`\alpha \mathrm{\Delta }`$, denote by $`H_\alpha `$ the element in $`๐ฅ`$ which corresponds to $`\alpha `$ under the isomorphism between $`๐ฅ`$ and $`๐ฅ^{}`$ induced by the Killing form $`(,)`$. We fix a simple system of roots $`\pi =\{\alpha _1,\mathrm{},\alpha _N\}`$ and denote by $`\mathrm{\Delta }^\pm `$ the corresponding positive/negative system. For any positive root $`\alpha \mathrm{\Delta }^+`$, we choose root vectors $`e_\alpha ๐ค_\alpha `$ and $`e_\alpha ๐ค_\alpha `$ which are dual with respect to $`(,)`$ so that $`[e_\alpha ,e_\alpha ]=H_\alpha `$. We also fix an orthonormal basis $`(x_i)_{1iN}`$ of $`๐ฅ`$. Lastly, for a subset of simple roots $`\pi ^{}\pi `$, we shall denote the root span of $`\pi ^{}`$ by $`<\pi ^{}>\mathrm{\Delta }`$ and set $`\overline{\pi }_{}^{}{}_{}{}^{\pm }=\mathrm{\Delta }^\pm <\pi ^{}>^\pm .`$
For any subset $`\pi ^{}\pi `$, we consider the following $`H`$-equivariant solution of the (mDYBE) (with $`K=\frac{1}{2}id_๐ค`$):
$$R(q)X=\underset{\alpha \mathrm{\Delta }}{}\varphi _\alpha (q)X_\alpha e_\alpha $$
$`(5.1a)`$
where
$$\begin{array}{cc}& \varphi _\alpha (q)=\frac{1}{2}\text{ for }\alpha \overline{\pi }_{}^{}{}_{}{}^{+},\varphi _\alpha (q)=\frac{1}{2}\text{ for }\alpha \overline{\pi }_{}^{}{}_{}{}^{}\hfill \\ & \varphi _\alpha (q)=\frac{1}{2}\mathrm{coth}(\frac{1}{2}(\alpha (q))\text{ for }\alpha <\pi ^{}>,\hfill \end{array}$$
$`(5.1b)`$
and $`X_\alpha =(X,e_\alpha ),\alpha \mathrm{\Delta }.`$ From now onwards, we shall assume $`G`$ and $`H`$ are simply-connected.
Consider now the coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$ which corresponds to this particular choice of $`R`$. By Proposition 4.2 (a), we know that the associated bundle maps $`^\pm `$ are morphisms of Lie algebroids, hence it follows that the dual maps $`(^\pm )^{}=^{}`$ are Poisson maps, when the domain and target are equipped with the corresponding Lie-Poisson structures. Note that the Lie-Poisson structure $`\{,\}_{A^{}\mathrm{\Gamma }}`$ on the dual bundle $`A^{}\mathrm{\Gamma }TU\times ๐ค`$ of the trivial Lie algebroid $`A\mathrm{\Gamma }`$ is a product structure, as is evident from the expression in Remark 2.5. Hence we have $`H`$-equivariant realizations of $`A^{}\mathrm{\Gamma }`$ (the dual of the trivial Lie algebroid $`A\mathrm{\Gamma }`$) in the dual vector vector bundle $`A\mathrm{\Gamma }`$ of the dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$.
To summarize, we have the following.
###### Proposition 5.1
$`(^\pm )^{}`$ are $`H`$-equivariant Poisson maps, where $`H`$ acts on $`A^{}\mathrm{\Gamma }`$, $`A\mathrm{\Gamma }TU\times ๐ค`$ by acting on the factor $`๐ค`$ by adjoint action.
To construct the spin Calogero-Moser system associated to the dynamical r-matrix $`R`$ in (5.1), introduce the quadratic function
$$Q(\xi )=\frac{1}{2}(\xi ,\xi ),\xi ๐ค.$$
$`(5.2)`$
We shall take $`\rho =(^+)^{}`$ to be our realization map (the other case with $`(^{})^{}`$ is similar) and let $`L=\mathrm{Pr}_3\rho `$, as in (3.3). Then the spin Calogero-Moser system associated to $`R`$ is the Hamiltonian system on $`A^{}\mathrm{\Gamma }TU\times ๐ค`$ generated by the Hamiltonian
$$(q,p,\xi )=L^{}Q(q,p,\xi )$$
$`(5.3)`$
Write $`p=_ip_ix_i`$, $`\xi =_i\xi _ix_i+_{\alpha \mathrm{\Delta }}\xi _\alpha e_\alpha `$, then we have
###### Proposition 5.2
The Hamiltonian of the spin Calogero-Moser system associated to the dynamical r-matrix $`R`$ in (5.1) is given by
$$\begin{array}{cc}\hfill (q,p,\xi )=& \frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{i}{}\xi _i^2+\frac{1}{2}\underset{i}{}p_i\xi _i\hfill \\ & \frac{1}{8}\underset{\alpha <\pi ^{}>}{}\frac{\xi _\alpha \xi _\alpha }{sinh^2\frac{1}{2}\alpha (q)}\hfill \end{array}$$
$`(5.4)`$
and is invariant under the Hamiltonian $`H`$-action on $`A^{}\mathrm{\Gamma }TU\times ๐ค`$:
$$h(q,p,\xi )=(q,p,Ad_h\xi )$$
$`(5.5)`$
with momentum map $`J:TU\times ๐ค๐ฅ`$ given by
$$J(q,p,\xi )=\mathrm{\Pi }_๐ฅ\xi .$$
$`(5.6)`$
Consider the level set $`J^1(0)`$ which is invariant under the flow $`\varphi _t`$ generated by $``$. Since $`J=\gamma \rho `$, where $`\gamma `$ is the momentum map in Proposition 3.1,we clearly have $`\rho (J^1(0))\gamma ^1(0)`$. Hence assumptions A1-A3 are satisfied. Therefore, the family of functions $`L^{}I(๐ค)`$ Poisson commute on $`J^1(0)`$ and hence descend to Poisson commuting functions on the reduced Poisson variety $`J^1(0)/H.`$
Remark 5.3 Note that if we consider the realization map $`\rho ^{}=(^{})^{}=^+`$ instead, then we would have the slightly different Hamiltonian
$$\begin{array}{cc}\hfill ^{}(q,p,\xi )=& \frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{i}{}\xi _i^2\frac{1}{2}\underset{i}{}p_i\xi _i\hfill \\ & \frac{1}{8}\underset{\alpha <\pi ^{}>}{}\frac{\xi _\alpha \xi _\alpha }{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}\hfill \end{array}$$
and the associated Lax operator in this case is given by $`L^{}(q,p,\xi )=pR^+(q)\xi `$.
###### Proposition 5.4
The Hamiltonian equations of motion generated by $``$ on $`A^{}\mathrm{\Gamma }`$ are given by
$$\begin{array}{cc}& \dot{q}=p+\frac{1}{2}\mathrm{\Pi }_๐ฅ\xi ,\hfill \\ & \dot{p}=\frac{1}{8}\underset{\alpha <\pi ^{}>}{}\frac{\mathrm{coth}\frac{1}{2}\alpha (q)}{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}\xi _\alpha \xi _\alpha H_\alpha ,\hfill \\ & \dot{\xi }=[\xi ,\frac{1}{4}\mathrm{\Pi }_๐ฅ\xi +\frac{1}{2}p\frac{1}{4}\underset{\alpha <\pi ^{}>}{}\frac{\xi _\alpha }{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}e_\alpha ]\hfill \\ & =[\xi ,R^+(q)L(q,p,\xi )].\hfill \end{array}$$
$`(5.7)`$
Moreover, under the Hamiltonian flow, we have
$$\begin{array}{cc}& (\mathrm{\Pi }_๐ฅ\xi )^{}=0\hfill \\ & \dot{L}(q,p,\xi )=[L(q,p,\xi ),R(q)L(q,p,\xi )]\hfill \\ & dR(q)(\mathrm{\Pi }_๐ฅ\xi )L(q,p,\xi ).\hfill \end{array}$$
$`(5.8)`$
###### Demonstration Proof
From the expression for the Poisson bracket in Remark 2.5, the equations of motion are given by $`\dot{q}=\delta _2`$, $`\dot{p}=\delta _1`$ and $`\dot{\xi }=[\xi ,\delta ].`$ Therefore, (5.7) follows by a direct computation. On the other hand, it follows from the definition of $`\rho `$ that $`m(q,p,\xi )=q`$ and $`\tau (q,p,\xi )=\mathrm{\Pi }_๐ฅ\xi `$ in the notation introduced in (3.1)-(3.2). Therefore, (5.8) is a consequence of Theorem 3.2 (c). $`\mathrm{}`$
We shall solve Eqn.(5.7) on the level set $`J^1(0)`$ (where $`\mathrm{\Pi }_๐ฅ\xi =0`$) in Section 6 below. In order to write down the equations of motion of the reduced Hamiltonian system, we have to restrict to a smooth component of $`J^1(0)/H=U\times ๐ฅ\times (๐ฅ^{}/H)`$. For this purpose, we consider the following open submanifold of $`๐ค`$:
$$๐ฐ=\{\xi ๐ค\xi _{\alpha _i}=(\xi ,e_{\alpha _i})0,i=1,\mathrm{},N\}.$$
$`(5.9)`$
Clearly, $`TU\times ๐ฐ`$ is a Poisson submanifold of $`TU\times ๐คA^{}\mathrm{\Gamma }`$ and the $`H`$-action defined by (5.5) induces a Hamiltonian action on $`TU\times ๐ฐ`$. Therefore, the corresponding momentum map is given by the restriction of the one in (5.6). To simplify notation, we shall denote this momentum map also by $`J`$ so that $`J^1(0)=TU\times (๐ฅ^{}๐ฐ)`$.
Now, recall from \[LX2\] that the formula
$$g(\xi )=\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}(C_{ji}\mathrm{log}\xi _{\alpha _j})h_{\alpha _i}\right)$$
$`(5.10)`$
defines an $`H`$-equivariant map $`g:๐ฐH`$, where $`C=(C_{ij})`$ is the inverse of the Cartan matrix and $`h_{\alpha _i}=\frac{2}{(\alpha _i,\alpha _i)}H_{\alpha _i}`$, $`i=1,\mathrm{},N`$. Using $`g`$, we can identify the reduced space $`J^1(0)/H=TU\times (๐ฅ^{}๐ฐ/H)`$ with $`TU\times ๐ค_{red}`$, where $`๐ค_{red}`$ is the affine subspace $`ฯต+_{\alpha \mathrm{\Delta }\pi }e_\alpha `$, and $`ฯต=_{j=1}^Ne_{\alpha _j}`$. Indeed, if we write $`\alpha =_{i=1}^Nm_\alpha ^i\alpha _i`$ for each $`\alpha \mathrm{\Delta }`$, then the identification map is given by
$$(q,p,[\xi ])(q,p,Ad_{g(\xi )^1}\xi ),$$
$`(5.11)`$
where explicitly,
$$Ad_{g(\xi )^1}\xi =ฯต+\underset{\alpha \mathrm{\Delta }\pi }{}\xi _\alpha \left(\underset{i=1}{\overset{N}{}}\xi _{\alpha _j}^{m_\alpha ^j}\right)e_\alpha .$$
$`(5.12)`$
Thus the natural projection $`\pi _0:J^1(0)TU\times ๐ค_{red}`$ is the map
$$(q,p,\xi )(q,p,Ad_{g(\xi )^1}\xi ).$$
$`(5.13)`$
We shall write $`s=_{\alpha \mathrm{\Delta }}s_\alpha e_\alpha `$ for $`s๐ค_{red}`$ (note that $`s_{\alpha _j}=1\text{for}j=1,\mathrm{},N`$). By Poisson reduction \[MR\], the reduced manifold $`TU\times ๐ค_{red}`$ has a unique Poisson structure which is a product structure , where the second factor $`๐ค_{red}`$ is equipped with the reduction (at 0) of the Lie- Poisson structure on $`๐ฐ`$ by the $`H`$-action. Now the symplectic leaves of $`๐ค_{red}`$ are the symplectic reduction of $`๐ช๐ฐ`$ at 0, where $`๐ช๐ค`$ is an adjoint orbit \[MR\]. In other words, any symplectic leaf of $`๐ค_{red}`$ is of the form $`(๐ช๐ฐ๐ฅ^{})/H`$, and we shall denote this by $`๐ช_{red}`$. Consequently, the symplectic leaves of $`TU\times ๐ค_{red}`$ are of the form $`TU\times ๐ช_{red}`$, which is of dimension equal to $`dim๐ช`$. Therefore, if $``$ is the Hamiltonian of the hyperbolic spin Calogero-Moser system in (5.4), then its reduction $`_0`$ on $`TU\times ๐ค_{red}`$ is given by
$$_0(q,p,s)=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{4}\underset{\alpha <\pi ^{}>^+}{}\frac{s_\alpha s_\alpha }{sinh^2\frac{1}{2}\alpha (q)},$$
$`(5.14)`$
where $`s๐ค_{red}`$.
Remark 5.5 (a) In the special case where $`\pi ^{}=\pi `$, the Hamiltonian system generated by $`_0`$ is isomorphic to the one in Reshetikhinโs paper \[R\]. (b) The family of integrable hyperbolic spin Calogero-systems constructed in this section are different from the ones in \[LX2\]. Although they look similar, however, their explicit integration requires different tools. To be more precise, the factorization problems for the systems in \[LX2\] are associated with infinite dimensional Lie groupoids whose vertex groups are loop groups. The solution of such factorization problems requires the use of algebraic geometry (compare Section 6 and \[L2\]). On the other hand, from the point of view of proving complete integrability, the two distinct families of hyperbolic systems also require totally different considerations. We shall discuss these matters in subsequent publications.
###### Proposition 5.6
The Hamiltonian equations of motion generated by $`_0`$ on the reduced Poisson manifold $`TU\times ๐ค_{red}`$ are given by
$`\dot{q}=p,`$
$`\dot{p}={\displaystyle \frac{1}{8}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{\mathrm{coth}\frac{1}{2}\alpha (q)}{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}}s_\alpha s_\alpha H_\alpha ,`$
$`\dot{s}=[s,]`$
where
$$\frac{=\frac{1}{4}\underset{\alpha <\pi ^{}>}{}\frac{s_\alpha }{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}e_\alpha +\frac{1}{4}\underset{i,j}{}C_{ji}}{\alpha <\pi ^{}>\pi ^{}\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }\frac{s_\alpha s_{\alpha _j\alpha }}{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}h_{\alpha _i}.}$$
(Here we use the notation $`[e_\alpha ,e_\beta ]=N_{\alpha ,\beta }e_{\alpha +\beta }`$ if $`\alpha +\beta \mathrm{\Delta }`$.)
###### Demonstration Proof
The first two equations are obvious from Proposition 5.4 and the definition of $`s`$. To derive the equation of $`s`$, we differentiate $`s=Ad_{g(\xi )^1}\xi `$ with respect to $`t`$, assuming $`\xi `$ satisfies the equation in Proposition 5.4 with $`\mathrm{\Pi }_๐ฅ\xi =0`$. Then we have
$`\dot{s}=`$ $`[T_{g(\xi )^1}r_{g(\xi )}{\displaystyle \frac{d}{dt}}g(\xi )^1,s]+Ad_{g(\xi )^1}\dot{\xi }`$
$`=`$ $`[s,{\displaystyle \frac{1}{2}}p{\displaystyle \frac{1}{4}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{\xi _\alpha }{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}}e^{\alpha (logg(\xi ))}e_\alpha T_{g(\xi )^1}r_{g(\xi )}{\displaystyle \frac{d}{dt}}g(\xi )^1].()`$
By a direct computation, we find $`\xi _\alpha e^{\alpha (logg(\xi ))}=\xi _\alpha (_{i=1}^N\xi _{\alpha _j}^{m_\alpha ^j})=s_\alpha `$. Meanwhile, by differentiating $`g(\xi )^1`$, we obtain
$$T_{g(\xi )^1}r_{g(\xi )}\frac{d}{dt}g(\xi )^1=\underset{i,j}{}C_{ji}\dot{\xi }_{\alpha _j}\xi _{\alpha _j}^{}{}_{}{}^{1}h_{\alpha _i}.$$
But
$`\dot{\xi }_{\alpha _j}`$ $`=(\dot{\xi },e_{\alpha _j})`$
$`=([\xi ,{\displaystyle \frac{1}{2}}p{\displaystyle \frac{1}{4}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{\xi _\alpha }{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}}e_\alpha ],e_{\alpha _j})`$
$`{\displaystyle \frac{={\displaystyle \frac{1}{2}}\alpha _j(p)\xi _{\alpha _j}+{\displaystyle \frac{1}{4}}\xi _{\alpha _j}{\displaystyle }}{\alpha <\pi ^{}>\pi ^{}}}`$
$`\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }{\displaystyle \frac{s_\alpha s_{\alpha _j\alpha }}{\mathrm{sinh}^2\frac{1}{2}\alpha (q)}}`$
whereas $`\frac{1}{2}p=_{i,j}C_{ji}\alpha _j(\frac{1}{2}p)h_{\alpha _i}`$. Therefore, on substituting the above expressions into (\*), the desired equation follows. $`\mathrm{}`$
In the rest of the section, we shall describe a scaling limit of the hyperbolic spin Calogero-Moser systems. More precisely, we consider
$$\begin{array}{cc}& q=x+2\tau w,\tau >0,\hfill \\ & \xi _i=\eta _i,1iN,\hfill \\ & \xi _\alpha =\eta _\alpha e^\tau ,\alpha \mathrm{\Delta },\hfill \end{array}$$
$`(5.15)`$
in the limit $`\tau \mathrm{}`$, where
$$w=\underset{\alpha \mathrm{\Delta }^+}{}\frac{H_\alpha }{(\alpha ,\alpha )}.$$
$`(5.16)`$
Note that this is analogous to the one in \[DP\], where the standard (spinless) elliptic Calogero-Moser system was considered. Clearly, $`\alpha (w)=(\alpha ,\delta ^{})`$, where
$$\delta ^{}=\frac{1}{2}\underset{\beta \mathrm{\Delta }^+}{}\beta ^{},\beta ^{}=\frac{2\beta }{(\beta ,\beta )}.$$
$`(5.17)`$
If for $`\alpha \mathrm{\Delta }`$, we write $`\alpha =_{i=1}^Nm_\alpha ^i\alpha _i`$, then it is not hard to show that
$$l(\alpha ):=\alpha (w)=\underset{i=1}{\overset{N}{}}m_\alpha ^i.$$
$`(5.18)`$
Therefore, $`l(\alpha )`$ is the level (or height) of $`\alpha `$. Hence $`l(\alpha )`$ is an integer, and assumes the value $`1`$ if and only if $`\alpha \pi `$.
Now, with the definition of $`x`$ and $`\eta _\alpha `$ in (5.15), it is easy to show that for $`\alpha <\pi ^{}>^+`$, we have
$$\frac{\xi _\alpha \xi _\alpha }{sinh^2\frac{1}{2}\alpha (q)}4\eta _\alpha \eta _\alpha e^{\alpha (x)2\tau (l(\alpha )1)},\tau \mathrm{}.$$
$`(5.19)`$
Therefore, if $`\alpha <\pi ^{}>\pi ^{}`$, we have
$$\underset{\tau \mathrm{}}{lim}\frac{\xi _\alpha \xi _\alpha }{sinh^2\frac{1}{2}\alpha (q)}=0.$$
$`(5.20)`$
On the other hand, if $`\alpha \pi ^{}`$, we obtain
$$\underset{\tau \mathrm{}}{lim}\frac{\xi _\alpha \xi _\alpha }{sinh^2\frac{1}{2}\alpha (q)}=4\eta _\alpha \eta _\alpha e^{\alpha (x)}.$$
$`(5.21)`$
Accordingly, the scaling limit of the Hamiltonian $``$ of the hyperbolic spin Calogero-Moser system is given by
$$\begin{array}{cc}\hfill ^s(x,p,\eta )=& \frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{i}{}\eta _i^2+\frac{1}{2}\underset{i}{}p_i\eta _i\hfill \\ & \underset{\alpha \pi ^{}}{}\eta _\alpha \eta _\alpha e^{\alpha (x)}.\hfill \end{array}$$
$`(5.22)`$
Note that in contrast to the spinless case, we do not know a priori the Poisson manifold on which $`^s`$ is defined. This issue will be settled below, but first we shall work out the scaling limits of the Hamiltonian equations of motion and the (quasi) Lax equation in Proposition 5.4 which will in fact give us some clue to this problem.
Let
$$\eta =\underset{i=1}{\overset{N}{}}\eta _ix_i+\underset{\alpha \mathrm{\Delta }}{}\eta _\alpha e_\alpha .$$
$`(5.23)`$
###### Proposition 5.7
The scaling limit of the Hamiltonian equations of motion in (5.7) is given by
$$\begin{array}{cc}& \dot{x}=p+\frac{1}{2}\mathrm{\Pi }_๐ฅ\eta ,\hfill \\ & \dot{p}=\underset{\alpha \pi ^{}}{}e^{\alpha (x)}\eta _\alpha \eta _\alpha H_\alpha ,\hfill \\ & \dot{\eta }=[\eta ,\frac{1}{4}\mathrm{\Pi }_๐ฅ\eta +\frac{1}{2}p].\hfill \end{array}$$
$`(5.24)`$
###### Demonstration Proof
The equation for $`x`$ is obvious from the equation for $`q`$. On the other hand, the equation for $`p`$ is a consequence of our previous analysis in (5.20)-(5.21) and the fact that $`\mathrm{coth}\frac{1}{2}\alpha (q)1`$ as $`\tau \mathrm{}`$ for $`\alpha \mathrm{\Delta }^+.`$ To get the last equation above, we make the substitution from (5.15) into the equation for $`\xi `$ and then divide both sides by $`e^\tau `$, this gives
$`(\mathrm{\Pi }_๐ฅ^{}\eta )^{}`$
$`=`$ $`[e^\tau \mathrm{\Pi }_๐ฅ\eta +\mathrm{\Pi }_๐ฅ^{}\eta ,{\displaystyle \frac{1}{4}}\mathrm{\Pi }_๐ฅ\eta +{\displaystyle \frac{1}{2}}p{\displaystyle \frac{1}{4}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{\eta _\alpha e^\tau }{\mathrm{sinh}^2\frac{1}{2}\alpha (x+2\tau w)}}e_\alpha ]`$
as $`(\mathrm{\Pi }_๐ฅ\eta )^{}=0`$. Therefore, upon letting $`\tau \mathrm{}`$, we find
$$(\mathrm{\Pi }_๐ฅ^{}\eta )^{}=[\mathrm{\Pi }_๐ฅ^{}\eta ,\frac{1}{4}\mathrm{\Pi }_๐ฅ\eta +\frac{1}{2}p].$$
Combining this with $`(\mathrm{\Pi }_๐ฅ\eta )^{}=0`$, the equation for $`\eta `$ follows. $`\mathrm{}`$
At this juncture, we remark that the Lax operator $`L(q,p,\xi )`$ does not actually admit a finite limit, as can be easily verified. However, we can remedy this by considering the following gauge-equivalent equation:
$$\begin{array}{cc}\hfill (Ad_{e^{\tau w}}L)^{}=& [Ad_{e^{\tau w}}L,Ad_{e^{\tau w}}R(q)L]\hfill \\ & Ad_{e^{\tau w}}dR(q)(\mathrm{\Pi }_๐ฅ\xi )L.\hfill \end{array}$$
$`(5.25)`$
Thus we introduce
$$\begin{array}{cc}& L_\tau (x,p,\eta ):=Ad_{e^{\tau w}}L(x+2\tau w,p,\mathrm{\Pi }_๐ฅ\eta +e^\tau \mathrm{\Pi }_๐ฅ^{}\eta ),\hfill \\ & M_\tau (x,p,\eta ):=Ad_{e^{\tau w}}R(x+2\tau w)L(x+2\tau w,p,\mathrm{\Pi }_๐ฅ\eta +e^\tau \mathrm{\Pi }_๐ฅ^{}\eta ).\hfill \end{array}$$
$`(5.26)`$
Using the relation $`Ad_{e^{\tau w}}e_\alpha =e^{\tau l(\alpha )}e_\alpha `$ and the $`H`$-equivariance of $`R`$, we easily find that
$$\begin{array}{cc}& L_\tau (x,p,\eta )\hfill \\ \hfill =& p+\frac{1}{2}\mathrm{\Pi }_๐ฅ\eta +\underset{\alpha <\pi ^{}>}{}\frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w)}\eta _\alpha e^{\tau (l(\alpha )1)}e_\alpha \hfill \\ & +\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}\eta _\alpha e^{\tau (l(\alpha )1)}e_\alpha ,\hfill \end{array}$$
$`(5.27)`$
whereas
$$\begin{array}{cc}& M_\tau (x,p.\eta )\hfill \\ \hfill =& \frac{1}{2}\underset{\alpha <\pi ^{}>}{}\mathrm{coth}\frac{1}{2}\alpha (x+2\tau w)\frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w)}\eta _\alpha e^{\tau (l(\alpha )1)}e_\alpha \hfill \\ & +\frac{1}{2}\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}\eta _\alpha e^{\tau (l(\alpha )1)}e_\alpha .\hfill \end{array}$$
$`(5.28)`$
Now for $`\alpha <\pi ^{}>^+`$, we have
$$\frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w)}e^{\tau (l(\alpha )1)}e^{\tau (l(\alpha )1)},\tau \mathrm{}.$$
$`(5.29)`$
Similarly, for $`\alpha <\pi ^{}>^{}`$, we obtain
$$\frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w)}e^{\tau (l(\alpha )1)}e^{\alpha (x)+\tau (l(\alpha )+1)},\tau \mathrm{}.$$
$`(5.30)`$
###### Proposition 5.8
We have
$$\begin{array}{cc}\hfill ๐(x,p,\eta ):=& \underset{\tau \mathrm{}}{lim}L_\tau (x,p,\eta )\hfill \\ \hfill =& p+\frac{1}{2}\mathrm{\Pi }_๐ฅ\eta +\underset{\alpha \pi }{}\eta _\alpha e_\alpha \underset{\alpha \pi ^{}}{}e^{\alpha (x)}\eta _\alpha e_\alpha ,\hfill \end{array}$$
$`(5.31)`$
$$\begin{array}{cc}\hfill ๐(x,p,\eta ):=& \underset{\tau \mathrm{}}{lim}M_\tau (x,p,\eta )\hfill \\ \hfill =& \frac{1}{2}\underset{\alpha \pi }{}\eta _\alpha e_\alpha \frac{1}{2}\underset{\alpha \pi ^{}}{}e^{\alpha (x)}\eta _\alpha e_\alpha .\hfill \end{array}$$
$`(5.32)`$
Moreover, the scaling limit of the (quasi) Lax equation (5.8) is given by
$$\dot{๐}=[๐,๐]=[๐,(๐)]$$
$`(5.33)`$
where $``$ is the constant $`r`$-matrix defined by
$$(\eta )=\frac{1}{2}\underset{\alpha \mathrm{\Delta }^{}}{}\eta _\alpha e_\alpha \frac{1}{2}\underset{\alpha \mathrm{\Delta }^+}{}\eta _\alpha e_\alpha .$$
$`(5.34)`$
###### Demonstration Proof
Using the asymptotics in (5.29)-(5.30), we obtain
$$\underset{\tau \mathrm{}}{lim}\frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w)}e^{\tau (l(\alpha )1)}=\{\begin{array}{cc}0,\hfill & \alpha <\pi ^{}>(\pi ^{}(\pi ^{}))\hfill \\ 1,\hfill & \alpha \pi ^{}\hfill \\ e^{\alpha (x)},\hfill & \alpha \pi ^{}\hfill \end{array}$$
from which the formulas for $`๐`$ and $`๐`$ follow. Therefore, in order to demonstrate the validity of (5.33), it remains to show that
$$\underset{\tau \mathrm{}}{lim}Ad_{e^{\tau w}}dR(x+2\tau w)(\mathrm{\Pi }_๐ฅ\eta )L(x+2\tau w,p,\mathrm{\Pi }_๐ฅ\eta +e^\tau \mathrm{\Pi }_๐ฅ^{}\eta )=\mathrm{\hspace{0.17em}0}.$$
By the $`H`$-equivariance of $`R`$ and its explicit expression,
$`Ad_{e^{\tau w}}dR(x+2\tau w)(\mathrm{\Pi }_๐ฅ\eta )L(x+2\tau w,p,\mathrm{\Pi }_๐ฅ\eta +e^\tau \mathrm{\Pi }_๐ฅ^{}\eta )`$
$`=`$ $`dR(x+2\tau w)(\mathrm{\Pi }_๐ฅ\eta )L_\tau (x,p,\eta )`$
$`=`$ $`{\displaystyle \underset{\alpha <\pi ^{}>}{}}\alpha (\mathrm{\Pi }_๐ฅ\eta )\eta _\alpha {\displaystyle \frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{(2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w))^3}}e^{\tau (l(\alpha )1)}e_\alpha .`$
But as $`\tau \mathrm{}`$, we have
$`{\displaystyle \frac{e^{\frac{1}{2}\alpha (x+2\tau w)}}{(2\mathrm{sinh}\frac{1}{2}\alpha (x+2\tau w))^3}}e^{\tau (l(\alpha )1)}`$
$``$ $`\{\begin{array}{cc}e^{\alpha (x)\tau (3l(\alpha )1)},\alpha <\pi ^{}>^+\hfill & \\ e^{2\alpha (x)+\tau (3l(\alpha )+1)},\alpha <\pi ^{}>^{}.\hfill & \end{array}`$
Hence the required property follows. $`\mathrm{}`$
Remark 5.9 (a) The constant $`r`$-matrix $``$ is the scaling limit of the dynamical $`r`$-matrix in the sense that $`(\xi )=lim_\tau \mathrm{}R(x+2\tau w)\xi .`$ (b) It is a remarkable fact that the (quasi) Lax equation (5.8) scales to the genuine Lax equation in (5.33). In other words, the obstruction to integrability dissolves in the scaling limit. (c) The reader should note that the scaling limit above is a singular limit. For this reason, the geometric structures are not preserved. As the reader will see in what follows, $`^s`$ is defined on a Poisson manifold different from that of $``$. Therefore, it is not surprising that their Hamiltonian realization would require separate consideration. We now describe a Hamiltonian formulation of the equations in Proposition 5.7. To do so, we consider the trivial Lie algebroid $`S=T๐ฅ\times ๐ค`$ over $`๐ฅ`$, where $`๐ค`$ is identified with the semi-direct product $`๐ฅ๐ฅ^{}`$ associated with the representation $`ad`$ of the Cartan subalgebra $`๐ฅ`$ in $`๐ฅ^{}`$. Thus the Lie algebroid bracket on $`S`$ is given by
$$\begin{array}{cc}& [(Z,X),(Z^{},X^{})]_S(x)\hfill \\ \hfill =& (dZ^{}(x)Z(x)dZ(x)Z^{}(x),dX^{}(x)Z(x)dX(x)Z^{}(x)\hfill \\ & +[\mathrm{\Pi }_๐ฅX(x),\mathrm{\Pi }_๐ฅ^{}X^{}(x)][\mathrm{\Pi }_๐ฅX^{}(x),\mathrm{\Pi }_๐ฅ^{}X(x)])\hfill \end{array}$$
$`(5.35)`$
where $`Z,Z^{}:๐ฅ๐ฅ`$ , $`X,X^{}:๐ฅ๐ค`$ are holomorphic maps and $`x๐ฅ`$.
###### Proposition 5.10
The Lie-Poisson structure on the dual bundle $`S^{}๐ฅ\times ๐ฅ\times ๐ค`$ of the trivial Lie algebroid $`S`$ is given by
$`\{\phi ,\psi \}_S^{}(x,p,\eta )`$
$`=`$ $`(\delta _1\psi ,\delta _2\phi )(\delta _1\phi ,\delta _2\psi )+(\eta ,[\mathrm{\Pi }_๐ฅ\delta \phi ,\mathrm{\Pi }_๐ฅ^{}\delta \psi ])+[\mathrm{\Pi }_๐ฅ^{}\delta \phi ,\mathrm{\Pi }_๐ฅ\delta \psi ])`$
and the Hamiltonian equations generated by $`\phi :S^{}`$ are:
$$\begin{array}{cc}& \dot{x}=\delta _2\phi ,\hfill \\ & \dot{p}=\delta _1\phi ,\hfill \\ & \dot{\eta }=[\eta ,\mathrm{\Pi }_๐ฅ\delta \phi ]+\mathrm{\Pi }_๐ฅ[\eta ,\delta \phi ].\hfill \end{array}$$
$`(5.36)`$
###### Demonstration Proof
Using the method of calculation in Section 2, we have
$$\begin{array}{cc}& \{\phi ,\psi \}_S^{}(x,p,\eta )\hfill \\ \hfill =& l_{[s(\phi ),s(\psi )]_S}(x,p,\eta )+(\delta _1\psi ,\delta _2\phi )(\delta _1\phi ,\delta _2\psi ).\hfill \end{array}$$
Now, from the expression for $`[,]_S`$, it is easy to check that
$`[s(\phi ),s(\psi )]_S(x)`$
$`=`$ $`(0,[\mathrm{\Pi }_๐ฅ\delta \phi ,\mathrm{\Pi }_๐ฅ^{}\delta \psi ][\mathrm{\Pi }_๐ฅ\delta \psi ,\mathrm{\Pi }_๐ฅ^{}\delta \phi ]).`$
Hence we have
$`l_{[s(\phi ),s(\psi )]_S}(x,p,\eta )`$
$`=`$ $`(\eta ,[\mathrm{\Pi }_๐ฅ\delta \phi ,\mathrm{\Pi }_๐ฅ^{}\delta \psi ]+[\mathrm{\Pi }_๐ฅ^{}\delta \phi ,\mathrm{\Pi }_๐ฅ\delta \psi ]).`$
Assembling the calculations, we obtain the formula for $`\{\phi ,\psi \}_S^{}(x,p,\eta )`$. $`\mathrm{}`$
To prepare for our next result, we need to introduce further constructs. First of all, let $`๐ธ^{}T๐ฅ\times ๐ค`$ be the coboundary dynamical Lie algebroid associated with the constant r-matrix $``$. Since $`๐ฅ`$ is Abelian, the Lie-Poisson structure on its dual bundle $`๐ธT๐ฅ\times ๐ค`$ takes the form
$$\begin{array}{cc}& \{\phi ,\psi \}_๐ธ(x,p,\eta )\hfill \\ \hfill =& (\eta ,[(\delta \phi )\delta _2\phi ,\delta \psi ]+[\delta \phi ,(\delta \psi )\delta _2\psi ])\hfill \\ & +(\delta _1\psi ,\mathrm{\Pi }_๐ฅ\delta \phi )(\delta _1\phi ,\mathrm{\Pi }_๐ฅ\delta \psi ).\hfill \end{array}$$
$`(5.37)`$
Now, define
$$\begin{array}{cc}& ๐:T๐ฅ\times ๐คS^{}๐ธT๐ฅ\times ๐ค\hfill \\ & (x,p,\eta )(x,\mathrm{\Pi }_๐ฅ\eta ,๐(x,p,\eta )).\hfill \end{array}$$
$`(5.38)`$
###### Theorem 5.11
$`๐`$ is an $`H`$-equivariant Poisson map, where the $`H`$-action on $`S^{}`$ given by
$`h(x,p,\eta )`$ $`=(x,p,Ad_h\eta )`$
$`=(x,p,\mathrm{\Pi }_๐ฅ\eta +\mathrm{\Pi }_๐ฅ^{}Ad_h\eta )`$
is Hamiltonian with equivariant momentum map $`๐:T๐ฅ\times ๐ค๐ฅ`$, $`(x,p,\eta )\mathrm{\Pi }_๐ฅ\eta `$. Moreover, the equations in Proposition 5.7 are the Hamiltonian equations generated by $`^s(x,p,\eta )=๐^{}Q(x,p,\eta )`$ in the Lie-Poisson structure $`\{,\}_S^{}`$ and admit $`๐^{}I(๐ค)`$ as a family of conserved quantities in involution.
###### Demonstration Proof
Let $`\phi `$, $`\psi C^{\mathrm{}}(๐ธ)`$. By direct calculation, we have
$`\delta (\phi \rho )(x,p,\eta )`$
$`=`$ $`\delta _2\phi (\rho (x,p,\eta ))+{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅ\delta \phi (\rho (x,p,\eta ))`$
$`+{\displaystyle \underset{\alpha \pi }{}}(\delta \phi (\rho (x,p,\eta )))_\alpha e_\alpha {\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}(\delta \phi (\rho (x,p,\eta ))_\alpha e_\alpha ,`$
$`\delta _1(\phi \rho )(x,p,\eta )=\delta _1\phi (\rho (x,p,\eta )+_{\alpha \pi ^{}}\delta \phi _\alpha \eta _\alpha e^{\alpha (x)}H_\alpha ,`$ and $`\delta _2(\phi \rho )(x,p,\eta )=\mathrm{\Pi }_๐ฅ\delta \phi (\rho (x,p,\eta )).`$ To simplify notation, let $`X=\delta \phi (\rho (x,p,\eta ))`$, $`Y=\delta _1\phi (\rho (x,p,\eta ))`$ and $`Z=\delta _2\phi (\rho (x,p,\eta ))`$ and denote the corresponding quantities associated with $`\psi `$ by $`X^{}`$, $`Y^{}`$ and $`Z^{}`$ respectively. Then it follows from the expression of $`\{,\}_S^{}`$ in Proposition 5.10 and the above calculation that
$`\{\phi \rho ,\psi \rho \}_S^{}(x,p,\eta )`$
$`=`$ $`{\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}\eta _\alpha X_\alpha ^{}\alpha \left(Z+{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅX\right)+{\displaystyle \underset{\alpha \pi }{}}\eta _\alpha X_\alpha ^{}\alpha \left(Z{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅX\right)`$
$`+(Y^{},\mathrm{\Pi }_๐ฅX)(XX^{},YY^{},ZZ^{}).`$
On the other hand,
$`\{\phi ,\psi \}_๐ธ\rho (x,p,\eta )`$
$`=`$ $`(๐(x,p,\eta ),[(X)Z,X^{}]+[X,(X^{})Z^{}])`$
$`+(Y^{},\mathrm{\Pi }_๐ฅX)(Y,\mathrm{\Pi }_๐ฅX^{}).`$
But
$`({\displaystyle \underset{\alpha \pi }{}}\eta _\alpha e_\alpha ,[(X)Z,X^{}]+[X,(X^{})Z^{}])`$
$`=`$ $`{\displaystyle \underset{\alpha \pi }{}}\eta _\alpha X_\alpha ^{}\alpha (Z{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅX)(XX^{},YY^{},ZZ^{}),`$
while
$`({\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}\eta _\alpha e_\alpha ,[(X)Z,X^{}]+[X,(X^{})Z^{}])`$
$`=`$ $`{\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}\eta _\alpha X_\alpha ^{}\alpha (Z+{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅX)(XX^{},YY^{},ZZ^{}).`$
Putting the calculations together, we conclude that $`\{\phi ,\psi \}_๐ธ\rho (x,p,\eta )`$ is identical to $`\{\phi \rho ,\psi \rho \}_S^{}(x,p,\eta )`$. Thus $`\rho `$ is a Poisson map. Alternatively, we can also establish the assertion by showing that the dual of the bundle map $`๐`$ is a morphism of Lie algebroids.
To show that the equations in Proposition 5.7 are the Hamiltonian equations generated by $`^s`$ in the Poisson structure $`\{,\}_S^{}`$, note that
$$\delta _1^s=\underset{\alpha \pi ^{}}{}\eta _\alpha \eta _\alpha e^{\alpha (x)}H_\alpha ,\delta _2^s=p+\frac{1}{2}\mathrm{\Pi }_๐ฅ\eta $$
and
$$\delta ^s=\frac{1}{2}p+\frac{1}{4}\mathrm{\Pi }_๐ฅ\eta \underset{\alpha \pi ^{}}{}e^{\alpha (x)}(\eta _\alpha e_\alpha +\eta _\alpha e_\alpha ).$$
From these formulas, it is clear that the equations for $`x`$ and $`p`$ from (5.36) with $`\phi =^s`$ are identical to the corresponding ones in Proposition 5.7. To show that the equations for $`\eta `$ are identical as well, it is enough to check that $`[\eta ,\delta ^s]๐ฅ^{}`$. Write $`\eta =\mathrm{\Pi }_๐ฅ\eta +\mathrm{\Pi }_๐ฅ^{}\eta `$ and substitute into $`[\eta ,\delta ^s]`$. As $`[๐ฅ,๐ฅ^{}]๐ฅ^{}`$, we only have to consider the term $`[\mathrm{\Pi }_๐ฅ^{}\eta ,_{\alpha \pi ^{}}e^{\alpha (x)}(\eta _\alpha e_\alpha +\eta _\alpha e_\alpha )]`$ in $`[\eta ,\delta ^s]`$. Expanding out, we have
$`[\mathrm{\Pi }_๐ฅ^{}\eta ,{\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}(\eta _\alpha e_\alpha +\eta _\alpha e_\alpha )]`$
$`=`$ $`{\displaystyle \underset{\alpha \pi ^{}}{}}{\displaystyle \underset{\beta \mathrm{\Delta }^+}{}}e^{\alpha (x)}\eta _\alpha \eta _\beta [e_\beta ,e_\alpha ]{\displaystyle \underset{\alpha \pi ^{}}{}}{\displaystyle \underset{\beta \mathrm{\Delta }^+}{}}e^{\alpha (x)}\eta _\alpha \eta _\beta [e_\beta ,e_\alpha ]`$
$`{\displaystyle \underset{\alpha \pi ^{}}{}}{\displaystyle \underset{\beta \mathrm{\Delta }^+}{}}e^{\alpha (x)}\eta _\alpha \eta _\beta [e_\beta ,e_\alpha ]{\displaystyle \underset{\alpha \pi ^{}}{}}{\displaystyle \underset{\beta \mathrm{\Delta }^+}{}}e^{\alpha (x)}\eta _\alpha \eta _\beta [e_\beta ,e_\alpha ].`$
Clearly, the first two terms of the above sum are in $`๐ฅ^{}`$. Now,
$`\mathrm{\Pi }_๐ฅ\left({\displaystyle \underset{\alpha \pi ^{}}{}}{\displaystyle \underset{\beta \mathrm{\Delta }^+}{}}e^{\alpha (x)}\eta _\alpha \eta _\beta [e_\beta ,e_\alpha ]\right)`$
$`=`$ $`{\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}\eta _\alpha \eta _\alpha H_\alpha `$
while
$`\mathrm{\Pi }_๐ฅ\left({\displaystyle \underset{\alpha \pi ^{}}{}}{\displaystyle \underset{\beta \mathrm{\Delta }^+}{}}e^{\alpha (x)}\eta _\alpha \eta _\beta [e_\beta ,e_\alpha ]\right)`$
$`=`$ $`{\displaystyle \underset{\alpha \pi ^{}}{}}e^{\alpha (x)}\eta _\alpha \eta _\alpha H_\alpha .`$
So the sum of the last two terms in the above sum is in $`๐ฅ^{}`$ as well. We shall leave the rest of the proof to the reader. $`\mathrm{}`$
We shall call the Hamiltonian systems generated by $`^s`$ in the Lie-Poisson structure $`\{,\}_S^{}`$ spin Toda lattices. To close this section, we shall consider reduction of the spin Toda lattices . As before, we consider the submanifold $`๐ฐ`$ defined in (5.9). Then clearly, the $`H`$-action defined in Theorem 5.11 induces a Hamiltonian action on $`T๐ฅ\times ๐ฐ`$. Denote the corresponding momentum map also by $`๐`$, we have $`๐^1(0)=T๐ฅ\times (๐ฅ^{}๐ฐ)`$. In this case, it is easy to verify that a generic orbit $`๐ช๐ค`$ (recall that $`๐ค๐ฅ๐ฅ^{}`$) is of dimension $`2N`$, where $`N=rank(๐ค)`$. Therefore, $`๐ช_{red}=(๐ช๐ฐ๐ฅ^{})/H`$ is a point. Indeed, we have
###### Corollary 5.13
The reduction of the Hamiltonian $`^s`$ of the spin Toda lattice on $`T๐ฅ\times ๐ช_{red}`$ is given by
$$_0^s(x,p)=\frac{1}{2}\underset{i}{}p_i^2\underset{\alpha \pi ^{}}{}c_\alpha e^{\alpha (x)}$$
$`(5.39)`$
where $`c_\alpha =s_\alpha `$ is a constant. Thus the Hamiltonian equations of motion generated by $`_0^s`$ are:
$$\begin{array}{cc}& \dot{x}=p,\hfill \\ & \dot{p}=\underset{\alpha \pi ^{}}{}c_\alpha e^{\alpha (x)}H_\alpha .\hfill \end{array}$$
$`(5.40)`$
$`\mathrm{}`$
Hence by reduction, we obtain a family of Toda lattices parametrized by subsets $`\pi ^{}`$ of $`\pi `$.
### 6. Solution of the hyperbolic spin Calogero-Moser systems and the abc spin Toda lattices
#### (a) The hyperbolic spin Calogero-Moser systems
We begin by solving the equation
$$\begin{array}{cc}& \frac{d}{dt}(q,0,L(q,p,\xi ))\hfill \\ \hfill =& (p,0,[L(q,p,\xi ),R(q)L(q,p,\xi )]).\hfill \end{array}$$
$`(6.1)`$
for the hyperbolic spin Calogero-Moser system which we obtain from Proposition 5.4 by restricting to the invariant manifold $`J^1(0)`$. As the reader will see, this will lead us eventually to the solution of the associated integrable model, whose equations are given in Proposition 5.6.
In order to set up the factorization problem properly, it is necessary to have precise knowledge of the Lie algebroids and Lie groupoids involved. Let us begin to describe these objects. Let $`๐^{}=๐ฅ+_{\alpha \mathrm{\Delta }^{}}๐ค_\alpha `$ and $`๐^+=๐ฅ+_{\alpha \mathrm{\Delta }^+}๐ค_\alpha `$ be opposing Borel subalgebras of $`๐ค`$. From the definition of $`R`$ in (5.1), we have
$$\begin{array}{cc}\hfill R^\pm (q)\xi =& \pm \frac{1}{2}\underset{i}{}\xi _ix_i\frac{1}{2}\underset{\alpha <\pi ^{}>}{}\frac{e^{\frac{1}{2}\alpha (q)}}{\mathrm{sinh}\frac{1}{2}\alpha (q)}\xi _\alpha e_\alpha \hfill \\ & \pm \underset{\alpha \overline{\pi }^{^{}}}{}\xi _\alpha e_\alpha .\hfill \end{array}$$
$`(6.2)`$
Therefore, $`R^+(q)\xi `$ is in the parabolic subalgebra
$$๐ญ_\pi ^{}^{}=๐^{}+\underset{\alpha <\pi ^{}>^+}{}๐ค_\alpha $$
$`(6.3)`$
containing $`๐^{}`$, while $`R^{}(q)\xi `$ is in the parabolic subalgebra
$$๐ญ_\pi ^{}^+=๐^++\underset{\alpha <\pi ^{}>^{}}{}๐ค_\alpha $$
$`(6.4)`$
containing $`๐^+`$. Now, recall that we have the standard decompositions \[Kn\]
$$๐ญ_\pi ^{}^\pm =๐ค_\pi ^{}+๐ซ_\pi ^{}^\pm $$
$`(6.5)`$
where
$$๐ค_\pi ^{}=๐ฅ+\underset{\alpha <\pi ^{}>}{}๐ค_\alpha $$
$`(6.6)`$
is the Levi factor of $`๐ญ_\pi ^{}^\pm `$, and
$$๐ซ_\pi ^{}^\pm =\underset{\alpha \overline{\pi }^^\pm }{}๐ค_\alpha $$
$`(6.7)`$
are the nilpotent radicals. Moreover, we have the identity
$$๐ค=๐ซ_\pi ^{}^{}+๐ค_\pi ^{}+๐ซ_\pi ^{}^+.$$
$`(6.8)`$
Let $`P_\pi ^{}^\pm `$, $`G_\pi ^{}`$ and $`N_\pi ^{}^\pm `$ be the simply-connected Lie subgroups of $`G`$ with corresponding Lie subalgebras $`๐ญ_\pi ^{}^\pm `$, $`๐ค_\pi ^{}`$ and $`๐ซ_\pi ^{}^\pm `$. Then we have
$$P_\pi ^{}^\pm =N_\pi ^{}^\pm G_\pi ^{}$$
$`(6.9)`$
and the submanifold
$$G_0=N_\pi ^{}^{}G_\pi ^{}N_\pi ^{}^+$$
$`(6.10)`$
is an open dense subset of $`G`$.
From the explicit expression for $`R^+`$ in (6.2) and the definition of $`^+`$, we find that
$$Im^+=\underset{qU}{}\{0_q\}\times ๐ญ_\pi ^{}^{}\times ๐ฅ.$$
$`(6.11)`$
Similarly, we have
$$Im^{}=\underset{qU}{}\{0_q\}\times ๐ญ_\pi ^{}^+\times ๐ฅ.$$
$`(6.12)`$
Therefore, the unique source-simply connected Lie groupoids of $`Im^\pm `$ are given by
$$\mathrm{\Gamma }_\pm =U\times P_\pi ^{}^{}\times U.$$
$`(6.13)`$
We next describe the ideals $`^\pm `$ of (4.7) for the case under consideration.
###### Lemma 6.1
$`^\pm =_{qU}\{0_q\}\times ๐ซ_\pi ^{}^{}\times \{0\}.`$
###### Demonstration Proof
We shall give the proof for $`^{}`$. Suppose $`(0_q,X,0)^{}`$, then there exists $`Z๐ฅ`$ such that $`^+(0_q,X,Z)=0`$. Equivalently, we have $`\iota Z+R^+(q)X=0`$ and $`\mathrm{\Pi }_๐ฅX=0`$. But from the explicit expression for $`R^+`$ in (6.1), we easily find that $`X_\alpha =0`$ for $`\alpha \mathrm{\Delta }^{}<\pi ^{}>^+`$. This shows that $`X๐ซ_\pi ^{}^+`$. The converse is clear by reversing the steps in the above argument. $`\mathrm{}`$
From this lemma, it follows that
$$\begin{array}{cc}\hfill Im^+/^+& =\underset{qU}{}\{0_q\}\times (๐ญ_\pi ^{}^{}/๐ซ_\pi ^{}^{})\times ๐ฅ\hfill \\ & \underset{qU}{}\{0_q\}\times ๐ค_\pi ^{}\times ๐ฅ\hfill \end{array}$$
$`(6.14)`$
where the identification map is given by
$$(0_q,X+๐ซ_\pi ^{}^{},Z)(0_q,\mathrm{\Pi }_{๐ค_\pi ^{}}^{}X,Z).$$
$`(6.15)`$
Here, $`\mathrm{\Pi }_{๐ค_\pi ^{}}^{}`$ is the projection onto $`๐ค_\pi ^{}`$ relative to the direct sum decomposition $`๐ญ_\pi ^{}^{}=๐ค_\pi ^{}+๐ซ_\pi ^{}^{}`$. Similarly,
$$\begin{array}{cc}\hfill Im^{}/^{}& =\underset{qU}{}\{0_q\}\times (๐ญ_\pi ^{}^+/๐ซ_\pi ^{}^+)\times ๐ฅ\hfill \\ & \underset{qU}{}\{0_q\}\times ๐ค_\pi ^{}\times ๐ฅ.\hfill \end{array}$$
$`(6.16)`$
This time, the identification is given by the map
$$(0_q,X+๐ซ_\pi ^{}^+,Z)(0_q,\mathrm{\Pi }_{๐ค_\pi ^{}}^+X,Z).$$
$`(6.17)`$
and $`\mathrm{\Pi }_{๐ค_\pi ^{}}^+`$ is the projection onto $`๐ค_\pi ^{}`$ relative to $`๐ญ_\pi ^{}^+=๐ค_\pi ^{}+๐ซ_\pi ^{}^+.`$
###### Proposition 6.2
The isomorphism $`\theta :Im^+/^+Im^{}/^{}`$ defined in Proposition 4.4 is given by
$$\theta (0_q,\mathrm{\Pi }_{๐ค_\pi ^{}}^{}X,Z)=(0_q,\iota Z+Ad_{e^q}\mathrm{\Pi }_{๐ค_\pi ^{}}^{}X,Z)$$
$`(6.18)`$
for all $`qU`$, $`X๐ญ_\pi ^{}^{}`$ and $`Z๐ฅ.`$
###### Demonstration Proof
From the expression for $`R^\pm (q)\xi `$, we have
$`\theta (0_q,\iota Z^{}+{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅ\xi {\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{e^{\frac{1}{2}\alpha (q)}}{\mathrm{sinh}\frac{1}{2}\alpha (q)}}\xi _\alpha e_\alpha ,\mathrm{\Pi }_๐ฅ\xi )`$
$`=`$ $`(0_q,\iota Z^{}{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅ\xi {\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{e^{\frac{1}{2}\alpha (q)}}{\mathrm{sinh}\frac{1}{2}\alpha (q)}}\xi _\alpha e_\alpha ,\mathrm{\Pi }_๐ฅ\xi ).`$
The assertion then follows from the identity $`Ad_{e^q}e_\alpha =e^{\alpha (q)}e_\alpha .`$ $`\mathrm{}`$
We shall make the natural identifications $`N_\pi ^{}^\pm \backslash P_\pi ^{}^\pm G_\pi ^{}`$ using the relation in (6.9) in what follows.
###### Corollary 6.3
The isomorphism $`\theta `$ can be lifted up to a Lie groupoid isomorphism
$$\begin{array}{cc}\hfill \mathrm{\Theta }:& U\times (N_\pi ^{}^{}\backslash P_\pi ^{}^{})\times UU\times (N_\pi ^{}^+\backslash P_\pi ^{}^+)\times U\hfill \\ & (u,๐^{}(g),v)(u,e^u๐^{}(g)e^v,v)\hfill \end{array}$$
$`(6.19)`$
where for $`gP_\pi ^{}^{}`$, $`๐^{}(g)G_\pi ^{}`$ denotes the factor in the unique factorization $`g=๐^{}(g)๐^{}(g)`$, $`๐^{}(g)N_\pi ^{}^{}.`$
###### Demonstration Proof
This is straightforward verification. $`\mathrm{}`$
We are now ready to solve Eqn.(6.1). To do so, we have to solve the factorization problem
$$exp\{t(0,0,L(q_0,p_0,\xi _0))\}(q_0)=\gamma _+(t)\gamma _{}(t)^1$$
$`(6.20)`$
for $`(\gamma _+(t),\gamma _{}(t))=((q_0,k_+(t),q(t)),(q_0,k_{}(t),q(t)))Im(^+,^{})`$ satisfying the condition in (4.19), where $`(q_0,p_0,\xi _0)J^1(0)=TU\times (๐ฐ๐ฅ^{})`$ is the initial value of $`(q,p,\xi )`$. We shall use the following notation (analogous to what we did in Corollary 6.3): for $`gP_\pi ^{}^+`$, $`๐^+(g)N_\pi ^{}^+`$, $`๐^+(g)G_\pi ^{}`$ will denote the factors in the unique factorization $`g=๐^+(g)๐^+(g)`$.
With the notation above, we have $`k_\pm (t)P_\pi ^{}^{}.`$ Therefore, the relation $`e^{tL(q_0,p_0,\xi _0)}=k_+(t)k_{}(t)^1`$ (which follows from (6.20)) can be rewritten in the form
$$e^{tL(q_0,p_0,\xi _0)}=๐^{}(k_+(t))๐^{}(k_+(t))๐^+(k_{}(t))^1๐^+(k_{}(t))^1.$$
But from Theorem 4.5 (b) and Corollary 6.3, we have $`๐^+(k_{}(t))=e^{q_0}๐^{}(k_+(t))e^{q(t)}.`$ Hence it follows that
$$\begin{array}{cc}& e^{tL(q_0,p_0,\xi _0)}\hfill \\ \hfill =& ๐^{}(k_+(t))(๐^{}(k_+(t))e^{q(t)}๐^{}(k_+(t))^1e^{q_0})๐^+(k_{}(t))^1.\hfill \end{array}$$
$`(6.21)`$
where the middle factor is in $`G_\pi ^{}`$ and $`๐^{}(k_\pm (t))N_\pi ^{}^{}.`$ We shall obtain the factors $`๐^{}(k_\pm (t))`$, $`๐^{}(k_+(t))`$ and $`q(t)`$ in several steps. First of all, from the fact that $`e^{tL(q_0,p_0,\xi _0)}P_\pi ^{}^+`$, we can find (as a consequence of (6.9)) unique $`g(t)G_\pi ^{}`$, $`n_+(t)N_\pi ^{}^+`$ satisfying $`n_+(0)=g(0)=1`$ such that
$$e^{tL(q_0,p_0,\xi _0)}=g(t)n_+(t)^1.$$
$`(6.22)`$
By comparing (6.21) and (6.22), we obtain
$$๐^{}(k_+(t))=1,๐^+(k_{}(t))=n_+(t).$$
$`(6.23)`$
Hence (6.21) reduces to the factorization problem
$$g(t)e^{q_0}=๐^{}(k_+(t))e^{q(t)}๐^{}(k_+(t))^1.$$
$`(6.24)`$
Since $`G_\pi ^{}`$ is a reductive Lie group, we can find (at least for small values of $`t`$) $`x(t)G_\pi ^{}`$ (unique up to transformations $`x(t)x(t)\delta (t)`$, where $`\delta (t)H`$) and unique $`d(t)H`$ such that
$$g(t)e^{q_0}=x(t)d(t)x(t)^1$$
$`(6.25)`$
with $`x(0)=1`$, $`d(0)=e^{q_0}`$. This uniquely determines $`q(t)`$ via the formula
$$q(t)=logd(t).$$
$`(6.26)`$
On the other hand, let us fix one such $`x(t)`$. We shall seek $`๐^{}(k_+(t))`$ in the form
$$๐^{}(k_+(t))=x(t)b(t),b(t)H.$$
$`(6.27)`$
To determine $`b(t)`$, we impose the condition in (4.19 a). After some calculations, we find that $`b(t)`$ satisfies the equation
$$\dot{b}(t)=T_el_{b(t)}\left(\frac{1}{2}\dot{q}(t)\mathrm{\Pi }_๐ฅ(T_{x(t)}l_{x(t)^1}\dot{x}(t))\right)$$
$`(6.28)`$
with $`b(0)=1`$. Solving the equation explicitly, we find that
$$๐^{}(k_+(t))=x(t)exp\{\frac{1}{2}(q(t)q_0)_0^t\mathrm{\Pi }_๐ฅ(T_{x(\tau )}l_{x(\tau )^1}\dot{x}(\tau ))๐\tau \}.$$
$`(6.29)`$
Combining (6.23) and (6.29), we finally have
$$k_+(t)=x(t)exp\{\frac{1}{2}(q(t)q_0)_0^t\mathrm{\Pi }_๐ฅ(T_{x(\tau )}l_{x(\tau )^1}\dot{x}(\tau ))๐\tau \}.$$
$`(6.30)`$
Hence we can write down the solution of Eqn.(6.1) by using (4.20). Note, however, we cannot determine $`\xi (t)`$ from the solution for $`L`$ as the expression for $`L`$ does not involve $`\xi _\alpha `$ for $`\alpha \overline{\pi }_{}^{}{}_{}{}^{}`$. The solution of Eqn.(5.7) on $`J^1(0)`$ is given in the following.
###### Theorem 6.4
Let $`(q_0,p_0,\xi _0)J^1(0)=TU\times (๐ฐ๐ฅ^{})`$. Then the Hamiltonian flow on $`J^1(0)`$ generated by
$$\begin{array}{cc}\hfill (q,p,\xi )=& \frac{1}{2}\underset{i}{}p_i^2+\frac{1}{8}\underset{i}{}\xi _i^2+\frac{1}{2}\underset{i}{}p_i\xi _i\hfill \\ & \frac{1}{8}\underset{\alpha <\pi ^{}>}{}\frac{\xi _\alpha \xi _\alpha }{sinh^2\frac{1}{2}\alpha (q)}\hfill \end{array}$$
with initial condition $`(q(0),p(0),\xi (0))=(q_0,p_0,\xi _0)`$ is given by
$$\begin{array}{cc}& q(t)=logd(t),\hfill \\ & \xi (t)=Ad_{k_+(t)^1}\xi _0,\hfill \\ & p(t)=Ad_{k_\pm (t)^1}L(q_0,p_0,\xi _0)\frac{1}{2}\underset{\alpha <\pi ^{}>}{}\frac{e^{\frac{1}{2}\alpha (q(t))}}{\mathrm{sinh}\frac{1}{2}\alpha (q(t))}\xi _\alpha (t)e_\alpha \hfill \\ & \underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}\xi _\alpha (t)e_\alpha \hfill \end{array}$$
$`(6.31)`$
where $`d(t)`$, $`k_\pm (t)`$ are constructed from the above procedure and in the expression for $`p(t)`$, the quantities $`q(t)`$ and $`\xi (t)`$ which appear on the right hand side are given by the formulas above that expression.
###### Demonstration Proof
We first show $`\xi (t)=Ad_{k_+(t)^1}\xi _0`$ solves the equation $`\dot{\xi }=[\xi ,R^+(q)L(q,p,\xi )]`$ in Proposition 5.4. To do so, we differentiate the expression for $`\xi (t)`$, this gives
$$\dot{\xi }(t)=[T_{k_+(t)^1}r_{k_+(t)}\frac{d}{dt}k_+(t)^1,\xi (t)].$$
But
$`T_{k_+(t)^1}r_{k_+(t)}{\displaystyle \frac{d}{dt}}k_+(t)^1`$
$`=`$ $`T_{k_+(t)}l_{k_+(t)^1}\dot{k}_+(t)`$
$`=`$ $`R^+(q(t))L(q(t),p(t),\xi (t))`$
from the argument in Theorem 4.7. Hence the claim. To get the formula for $`p(t)`$, we simply equate the following two expressions for $`L((q(t),p(t),\xi (t))`$, namely, $`L(q(t),p(t),\xi (t))=Ad_{k_\pm (t)^1}L(q_0,p_0,\xi _0)`$ and
$`L(q(t),p(t),\xi (t))=`$ $`p(t)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha <\pi ^{}>}{}}{\displaystyle \frac{e^{\frac{1}{2}\alpha (q(t))}}{\mathrm{sinh}\frac{1}{2}\alpha (q(t))}}\xi _\alpha (t)e_\alpha `$
$`+{\displaystyle \underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}}\xi _\alpha (t)e_\alpha .`$
This completes the proof. $`\mathrm{}`$
We next turn to the solution of the associated integrable model on $`TU\times ๐ค_{red}`$ with Hamiltonian $`_0(q,p,s)=\frac{1}{2}_ip_i^2\frac{1}{4}_{\alpha <\pi ^{}>^+}\frac{s_\alpha s_\alpha }{sinh^2\frac{1}{2}\alpha (q)}`$ and with equations of motion given in Proposition 5.6.
###### Corollary 6.5
Let $`(q_0,p_0,s_0)TU\times ๐ค_{red}`$ and suppose $`s_0=Ad_{g(\xi _0)^1}\xi _0`$ where $`\xi _0๐ฐ๐ฅ^{}.`$ Then the Hamiltonian flow generated by $`_0`$ with initial condition $`(q(0),p(0),s(0))=(q_0,p_0,s_0)`$ is given by
$$\begin{array}{cc}& q(t)=logd(t),\hfill \\ & s(t)=Ad_{\left(\stackrel{~}{k}_+(t)g\left(Ad_{\stackrel{~}{k}_+(t)^1}s_o\right)\right)^1}s_0,\hfill \\ & p(t)=Ad_{\left(\stackrel{~}{k}_+(t)g\left(Ad_{\stackrel{~}{k}_+(t)^1}s_o\right)\right)^1}(p_0R^{}(q_0)s_0)+R^{}(q(t))s(t).\hfill \end{array}$$
$`(6.32)`$
where $`\stackrel{~}{k}_+(t)=g(\xi _0)^1k_+(t)g(\xi _0)`$ and $`k_+(t)`$, $`d(t)`$ are as in Theorem 6.4.
###### Demonstration Proof
We shall obtain the Hamiltonian flow generated by $`_0`$ by reduction. Using the relation $`\varphi _t^{red}\pi _0=\pi _0\varphi _ti_0`$ from Theorem 3.2 (c), we have $`\varphi _t^{red}(q_0,p_0,s_0)=(q(t),p(t),Ad_{g(\xi (t))^1}\xi (t))`$ where $`q(t)`$, $`p(t)`$ $`\xi (t)`$ are given by the expressions in Theorem 6.4. Thus
$`s(t)=`$ $`Ad_{g(\xi (t))^1}\xi (t)`$
$`=`$ $`Ad_{g\left(Ad_{k_+(t)^1}\xi _0\right)^1}Ad_{k_+(t)^1g(\xi _0)}s_0`$
$`=`$ $`Ad_{\left(\stackrel{~}{k}_+(t)g\left(Ad_{\stackrel{~}{k}_+(t)^1}s_o\right)\right)^1}s_0`$
where we have used the $`H`$-equivariance of the map $`g`$. To express $`p(t)`$ in the desired form, introduce the gauge transformation of $`L`$: $`\stackrel{~}{L}(q,p,s)=Ad_{g(\xi )^1}L(q,p,\xi )=pR^{}(q)s`$, where as before, $`s=Ad_{g(\xi )^1}\xi `$. Then
$`\stackrel{~}{L}(q(t),p(t),s(t))=`$ $`Ad_{g(\xi (t))^1}Ad_{k_+(t)^1}L(q_0,p_0,\xi _0)`$
$`=`$ $`Ad_{\left(\stackrel{~}{k}_+(t)g\left(Ad_{\stackrel{~}{k}_+(t)^1}s_o\right)\right)^1}(p_0R^{}(q_0)s_0).`$
But $`\stackrel{~}{L}(q(t),p(t),s(t))`$ is also equal to $`p(t)R^{}(q(t))s(t)`$. By equating the two expressions, we obtain the desired expression for $`p(t)`$, as claimed. $`\mathrm{}`$
Remark 6.6 (a)It is easy to show that the element $`\stackrel{~}{k}_+(t)=g(\xi _0)^1k_+(t)g(\xi _0)`$ depends only on $`s_0`$, and not on the particular element $`\xi _0๐ฐ๐ฅ^{}`$ for which $`Ad_{g(\xi _0)^1}\xi _0=s_0`$. Indeed, from the factorization $`e^{L(q_0,p_0,\xi _0)}=k_+(t)k_{}(t)^1`$, we see that if we replace$`\xi _0`$ by $`Ad_h\xi _0`$, $`hH`$, then the factors $`k_\pm (t)`$ will be replaced by $`hk_\pm (t)h^1`$. As $`g(Ad_h\xi _0)=hg(\xi _0)`$, our assertion easily follows. Note that this is exactly the reason why we have chosen to express $`s(t)`$ and $`p(t)`$ in the form given in the above Corollary. (b) In \[L1\], we introduced a family of hyperbolic spin Ruijsenaars-Schneider models on the coboundary dynamical Poisson groupoids $`(\mathrm{\Gamma }=U\times G\times U,\{,\}_R)`$ associated to the same $`R`$โs which we use here. Recall that these are generated by nonzero multiples of $`H_i=Pr_2^{}\chi _i`$, $`i=1,\mathrm{},N,`$ where $`Pr_2`$ in this present case denotes projection onto the second factor of $`\mathrm{\Gamma }`$, and $`\chi _1,\mathrm{},\chi _N`$ are the characters of the irreducible representations corresponding to the fundamental weights $`\omega _1,\mathrm{},\omega _N`$ \[St\]. If we take the Hamiltonian $`H_i`$, say, then its Hamiltonian flow on the gauge group bundle $`\mathrm{\Gamma }`$ is defined by the equation
$`{\displaystyle \frac{d}{dt}}(q,g,q)`$
$`=`$ $`({\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅD\chi _i(g),{\displaystyle \frac{1}{2}}T_er_gR(q)(D\chi _i(g)){\displaystyle \frac{1}{2}}T_el_gR(q)(D\chi _i(g)),{\displaystyle \frac{1}{2}}\mathrm{\Pi }_๐ฅD\chi _i(g))`$
where $`D\chi _i(g)`$ is the right gradient of $`\chi _i`$. In this case, the factorization problem on $`\mathrm{\Gamma }`$ (from \[L1\] and our analysis above) gives
$`e^{tD\chi _i(g_0)}`$
$`=`$ $`๐^{}(k_+(t))(๐^{}(k_+(t))e^{q(t)}๐^{}(k_+(t))^1e^{q_0})๐^+(k_{}(t))^1`$
using the same notation as before (of course, the $`k_\pm `$ here are different from the ones above). If for small $`t`$, $`n_\pm (t)N_\pi ^{}^\pm `$, $`g(t)G_\pi ^{}`$ are the unique solution of the factorization problem
$$e^{tD\chi _i(g_0)}=n_{}(t)g(t)n_+(t)^1$$
satisfying $`n_\pm (0)=g(0)=1`$, then
$$๐^{}(k_+(t))=n_{}(t),๐^+(k_{}(t))=n_+(t).$$
Consequently, the factorization problem reduces to
$$g(t)e^{q_0}=๐^{}(k_+(t))e^{q(t)}๐^{}(k_+(t))^1$$
and therefore the solution proceeds as before. Finally, we can write down the Hamiltonian flow on $`\mathrm{\Gamma }`$ using the formula from \[L1\], namely,
$`(q(t),g(t),q(t))`$
$`=`$ $`(q_0,k_\pm (t),q(t))^1(q_0,g_0,q_0)(q_0,k_\pm (t),q(t)).`$
#### (b) The spin Toda lattices
In this final subsection, we shall discuss the solution of the spin Toda lattices. In this case, we have
$$Im^\pm =\underset{q๐ฅ}{}\{0_q\}\times ๐^{}\times ๐ฅ$$
$`(6.33)`$
where $`๐^\pm `$ are the opposing Borel subalgebra of $`๐ค`$ introduced at the beginning of the section. Let $`B^\pm `$ be the simply-connected Lie subgroups of $`G`$ which integrate $`๐^\pm `$, then the unique source-simply connected Lie groupoid of $`Im^\pm `$ are given by
$$\mathrm{\Gamma }_\pm =๐ฅ\times B^{}\times ๐ฅ.$$
$`(6.34)`$
Now, the ideals $`^\pm `$ of (4.7) in the present case are:
$$^\pm =\underset{q๐ฅ}{}\{0_q\}\times ๐ซ^{}\times \{0\}$$
$`(6.35)`$
where $`๐ซ^\pm =_{\alpha \mathrm{\Delta }^\pm }๐ค_\alpha `$. We shall denote by $`N^\pm `$ the simply-connected Lie subgroups of $`G`$ with $`Lie(N^\pm )=๐ซ^\pm `$. To cut the story short, we have the following result when we go through the analysis.
###### Theorem 6.7
Let $`(x_0,p_0,\eta _0)T๐ฅ\times ๐คS^{}`$ and let $`n_\pm (t)N^\pm `$, $`h(t)H`$ be the unique solution of the factorization problem
$$e^{t๐(x_0,p_0,\eta _0)}=n_{}(t)h(t)n_+(t)^1$$
$`(6.36)`$
(valid for $`0t<T`$ for some $`T>0`$) satisfying $`n_\pm (0)=h(0)=1`$. Then the Hamiltonian flow on $`T๐ฅ\times ๐คS^{}`$ generated by the Hamiltonian
$`^s(x,p,\eta )=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}p_i^2+{\displaystyle \frac{1}{8}}{\displaystyle \underset{i}{}}\eta _i^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}p_i\eta _i`$
$`{\displaystyle \underset{\alpha \pi ^{}}{}}\eta _\alpha \eta _\alpha e^{\alpha (x)}`$
with initial condition $`(x(0),p(0),\eta (0))=(x_0,p_0,\eta _0)`$ is given by
$$\begin{array}{cc}& x(t)=x_0+logh(t),\hfill \\ & \eta (t)=Ad_{e^{\frac{1}{2}logh(t)}}\eta _0,\hfill \\ & p(t)=Ad_{k_\pm (t)^1}๐(x_0,p_0,\eta _0)\frac{1}{2}\mathrm{\Pi }_๐ฅ\eta _0\underset{\alpha \pi }{}e^{\frac{1}{2}\alpha (logh(t))}(\eta _0)_\alpha e_\alpha \hfill \\ & +\underset{\alpha \pi ^{}}{}e^{\alpha (x_0+\frac{1}{2}logh(t))}(\eta _0)_\alpha e_\alpha \hfill \end{array}$$
$`(6.37)`$
where
$$k_\pm (t)=n_{}(t)e^{\pm \frac{1}{2}logh(t)}.$$
$`(6.38)`$
###### Demonstration Proof
The expression for $`x(t)`$ is clear if we write down the expression analogous to (6.21) and compare that with the factorization in (6.36). On the other hand, it is clear from the same expression that $`k_+(t)=n_{}(t)b_{}(t)`$ where $`b_{}(t)H`$ is to be determined from the condition given in (4.19). If we spell this out, we find that $`b_{}(t)`$ satisfies the equation
$$\dot{b}_{}(t)=\frac{1}{2}T_el_{b_{}(t)}\dot{x}(t)$$
with $`b_{}(0)=1`$. Solving the equation explicitly, we have $`b_{}(t)=e^{\frac{1}{2}(x(t)x_0)}`$. Now, in order to solve the equation for $`\eta `$ in (5.24), the crucial point to note is that we can rewrite this equation as $`\dot{\eta }=[\eta ,\frac{1}{2}\mathrm{\Pi }_๐ฅ๐(x,p,\eta )]`$. Finally, we can obtain the formula for $`p(t)`$ by equating the following two expressions for $`๐(x(t),p(t),\eta (t))`$, namely, $`๐(x(t),p(t),\eta (t))=Ad_{k_\pm (t)^1}๐(x_0,p_0,\eta _0)`$ and
$$\begin{array}{cc}\hfill ๐(x(t),p(t),\eta (t))=& p(t)+\frac{1}{2}\mathrm{\Pi }_๐ฅ\eta _0+\underset{\alpha \pi }{}\eta _\alpha (t)e_\alpha \underset{\alpha \pi ^{}}{}e^{\alpha (x(t))}\eta _\alpha (t)e_\alpha .\hfill \end{array}$$
This completes the proof. $`\mathrm{}`$
The solution of the family of Toda lattices in Corollary 5.12 is now straightforward. We shall leave the easy details to the reader.
## Appendix
#### Proof of Lemma 4.1
From the definition of $``$ and the expression for $`[,]_{A\mathrm{\Gamma }}`$, $`[,,]_{A^{}\mathrm{\Gamma }}`$, we have
$`[(0,A,Z),(0,A^{},Z^{})]_{A\mathrm{\Gamma }}(q)[(0,A,Z),(0,A^{},Z^{})]_{A^{}\mathrm{\Gamma }}(q)`$
$`=`$ $`(0_q,๐,๐ต)`$
where (after the obvious cancellations)
$`๐=`$ $`[R(q)A(q),R(q)A^{}(q)]+<dR(q)()A(q),A^{}(q)>+\{R(q)ad_{R(q)A(q)}^{}A^{}(q)`$
$`+dR(q)\iota ^{}A(q)(A^{}(q))dR(q)ad_{Z(q)}^{}q(A^{}(q))[Z(q),R(q)A^{}(q)]`$
$`R(q)ad_{Z(q)}^{}A^{}(q)(AA^{},ZZ^{})\}`$
and
$$๐ต=\iota ^{}ad_{R(q)A(q)}^{}A^{}(q)\iota ^{}ad_{R(q)A^{}(q)}^{}A(q)+ad_{<dR(q)()A(q),A^{}(q)>}^{}q.$$
Since $`R`$ is $`H`$-equivariant, we can show that $`๐ต=0`$ and
$$dR(q)ad_{Z(q)}^{}q(A^{}(q))+[Z(q),R(q)A^{}(q)]+R(q)ad_{Z(q)}^{}A^{}(q)=\mathrm{\hspace{0.17em}0}.$$
Therefore, the expression for $`๐`$ becomes
$`๐=`$ $`[R(q)A(q),R(q)A^{}(q)]+<dR(q)()A(q),A^{}(q)>`$
$`+R(q)(ad_{R(q)A(q)}^{}a^{}(q)ad_{R(q)A^{}(q)}^{}A(q))`$
$`+dR(q)\iota ^{}A(q)(A^{}(q))dR(q)\iota ^{}A^{}(q)(A(q))`$
$`=`$ $`[K(A(q)),K(A^{}(q))],`$
as desired. |
warning/0506/hep-ex0506030.html | ar5iv | text | # TOP QUARK MASS MEASUREMENTS AT THE TEVATRON
## 1 Introduction
The recent publication of the improved Run I measurement of the top mass by Dร $`^\mathrm{?}`$ was exciting for two reasons. First of all it demonstrated how much improvement in measurement precision could be achieved using a more advanced analysis technique like the Matrix Element method. Secondly, it was a reminder of how little we yet know about the properties of the top quark and that new experimental information about the top quark can have big implications for electroweak fits in the Standard Model. The current (Run I only) world average value for the top quark mass is $`178.0\pm 4.3`$ GeV$`/c^2`$. In the coming years the measurements of CDF and Dร combined should lead to a precision of about 2 GeV. Together with expected improvements in the measurement of the W boson mass this will allow to further constrain the Higgs boson mass to a relative precision of approximately 30%, as discussed elsewhere in these proceedings $`^\mathrm{?}`$.
Since the start of Run II both CDF and Dร have recorded more than 600 pb<sup>-1</sup> of data, already 5 times the Run I luminosity. The preliminary results presented here are based on fraction of the recorded data ranging from 160 to 230 pb<sup>-1</sup>.
## 2 Run II Top mass results
In $`p\overline{p}`$ collisions with $`\sqrt{s}=1.96`$ TeV at the Tevatron, top quarks can be produced via the strong interaction in $`t\overline{t}`$ pairs, or as single top quarks through the weak interaction. Single top production is predicted to have a lower cross-section and a more challenging event signature, and has not yet been observed at the time of this conference. For the Top mass measurement therefore only top pair events are used. Each top quark decays immediately to a $`W`$ boson and a $`b`$ quark, and the $`W`$ bosons decay either hadronically or leptonically, giving rise to 3 possible decay channels: di-lepton, lepton+jets and all-jets.
An overview of recent $`t\overline{t}`$ cross-section results from the CDF and Dร experiments in all three of the above final states is given elsewhere in these proceedings $`^\mathrm{?}`$. In both collaborations several top mass analyses are being developed in the di-lepton and lepton+jets decays channels, mostly based on very similar event selections. No preliminary Run II results in the all-jets channel have been presented so far.
A complete and up-to-date overview of ongoing Run II analyses can be found on the collaborationsโ public results web pages $`^{\mathrm{?},\mathrm{?}}`$. A description of all analyses is outside the scope of these proceedings. Below a few of the analyses are briefly described in order to highlight some important aspects of the top mass measurement.
### 2.1 Final states with two leptons plus jets
The striking signature due to the presence of two leptons in the final state allows for a relatively pure selection of top events, typically with a signal-to-background ratio of 4/1. The main challenge however is to fully reconstruct the kinematics of the final state, which are underconstrained due to the presence of two neutrinos. Different approaches exist to add an extra constraint to the system, and see for which value of the top mass the observed events are most likely.
In Table 1 several Run II analyses are listed with their preliminary results. Currently the most precise result was obtained by CDF with the neutrino weighting analysis using a loosened lepton identification (one lepton + one isolated track), optimizing the statistical precision by using a higher efficiency (and slightly lower purity) selection. In this method the rapidities of both neutrinos are used as extra constraints, and a weight as function of the top mass is calculated by integrating over all possible rapidity values and comparing the reconstructed missing transverse momentum with the observed momentum imbalance using a Gaussian resolution. For each event the top mass value which leads to the highest weight is plotted and fitted using Monte Carlo Templates, as shown in Figure 1.
### 2.2 Final states with one lepton plus jets
While the lepton+jets channel benefits from a higher branching ratio, it suffers from significant backgrounds from $`W`$+jets and non-$`W`$ multi-jet events.
Since only one neutrino is present the final state can be fully reconstructed. Some analyses use a constrained kinematic fit to further improve the measurement of lepton and jets beyond detector resolution. The CDF Dynamic Likelihood Method (DLM) follows a different approach, similar to the Dร Matrix Element method $`^\mathrm{?}`$; transfer functions are derived from Monte Carlo simulation describing the jet energy resolution. These functions are subsequently used in a multi-dimensional integration over phase space calculating the likelihood that the event is compatible with matrix elements describing top pair production and decay.
In order to reconstruct the invariant mass of the top decay products, a choice has to be made to assign jets and lepton to the corresponding top or anti-top quark. In a lepton+jets event 12 ways exist to do this assignment. Some analyses take only one jet assignment per event in consideration. The CDF Dynamic Likelihood Method and the Dร Ideogram analysis include all possible jet assignments in the fit.
The CDF and Dร template methods use an overall fit of Monte Carlo templates to the data in order to extract the mass. The CDF Dynamic Likelihood Method and Dร Ideogram analysis derive an event-by-event likelihood to maximize the statistical information extracted from each event. The Ideogram method also includes the hypothesis that the event could be background, weighted according to an estimated event purity.
Both experiments apply b-tagging in some of the top mass analyses. One advantage of b-tagging is to strongly reduce the backgrounds. A second advantage of b-tagging for the top mass measurement in the lepton+jets channel is the reduction of the number of possible jet assignments in the case that one or two jets are b-tagged. The CDF Template analysis combines the 0-tag, 1-tag and double tagged event samples in the fit to optimize the statistical precision. Dรโs first top mass analysis with b-tagging uses events with at least one tag, which applied to a data set of 230 pb<sup>-1</sup> leads to the most precise preliminary Run II top mass result presented so far. Figure 1 shows the fitted mass for the lowest-$`\chi ^2`$ solution for the b-tagged Dร Template analysis, compared to the Monte Carlo prediction.
An overview of the current preliminary results is shown in Table 1.
## 3 Prospects for the Top mass measurement
In all results reported here the dominant component of the systematic uncertainty is the uncertainties related to the jet energy scale. In the last year a lot of work has been done to improve the calibration of the reconstructed jet energies. CDF reports an improvement of a factor two or more in jet energy scale uncertainties compared to a year ago. Similar improvements are expected in Dร. This will have a direct effect on the systematic uncertainties quoted.
Further improvements in understanding the Jet Energy Scale can come from performing an in-situ calibration of the light-jet energy scale using the jets from the hadronic decay of the $`W`$ in the same $`t\overline{t}`$ events used to measure the top mass, and from studies in progress aimed at determining the b-jet energy scale from data.
Other systematics that are being studied are the modeling of initial state and final state gluon radiation in the $`t\overline{t}`$ Monte Carlo.
Very soon both experiments hope to present preliminary results with updated jet energy scale and an integrated luminosity of more than 300 pb<sup>-1</sup>.
All together the prospects are very good for having new top mass results this year with a precision comparable to or better than the current world average for each of the Tevatron experiments. This will open the door to an exciting new area of top physics to be further explored in the coming years at the Tevatron.
## References |
warning/0506/astro-ph0506509.html | ar5iv | text | # Capabilities of UV Coronagraphic Spectroscopy for Studying the Source Regions of Solar Energetic Particles and the Solar Wind
## 1 Introduction
In the following, we describe the capabilities of ultraviolet coronagraphic spectroscopy to address the SEP/flare/CME and solar wind problems, and we provide a brief description of a science payload that is capable of carrying out the required observations.
## 2 CMEs, Flares, and SEPs
UVCS/SOHO observations have provided new insights into the roles of shock waves, reconnection, and magnetic helicity in CME eruptions (Raymond 2002).
### 2.1 CME Shocks
The key parameters in theories of particle acceleration by shocks are the pre-shock plasma conditions (including seed particle population), the shock speed, and the angle between the magnetic field and the shock motion.
UVCS has observed CME-driven shocks through their effect on the widths of UV spectral lines (Raymond et al. 2000; Mancuso et al. 2002). These observations provide information about the compression ratio in the shock, a crucial parameter for predicting SEP spectra, and information about the thermal equilibration among electrons, protons and heavier ions. Electron heating is relatively modest, and the line widths of oxygen and silicon ions imply temperatures far higher than the proton temperatures, a potentially important consideration for models of SEP composition. In some events, the shock compression ratio can be determined from Type II burst band splitting (e.g., Vrลกnak et al. 2002), but not all radio bursts contain enough detail for this diagnostic to be useful. For a much larger fraction of events, UV spectroscopy can be used to determine the compression ratio via two independent techniques: (1) measuring pre- to post-shock temperature ratios (i.e., both $`T_p`$ and $`T_e`$) with resonant and Thomson-scattered H I Ly$`\alpha `$ line widths, then using adiabatic theory to compute the density ratio; (2) measuring ion temperatures of several species having different charges and masses and applying collisionless theories of multi-ion shock heating (e.g., Lee et al. 1987; Lee & Wu 2000) to compute the compression ratio that is most consistent with the observations (for both methods, see Mancuso et al. 2002).
UVCS routinely obtains the densities, ionization states and elemental abundances in the pre-CME corona (e.g., Raymond et al. 2003). The densities obtained by UVCS can be combined with Type II radio burst drift rates to obtain shock speeds. Upper limits on the coronal Alfvรฉn speed above active regions were inferred from the derived shock speeds by requiring that the disturbances propagate at least as fast as the local characteristic speed (Mancuso et al. 2003). For a subset of events, the resulting shock speeds are in much better agreement with LASCO CME expansion rates than shock speeds based upon average coronal density profiles, although there are some uncertainties in the measured drift rates (Mancuso & Raymond 2004). The Alfvรฉn and shock speeds can be inferred from detection of the shock arrival at different heights as determined by the timing of the increase in line widths of UV emission lines (e.g., Ciaravella et al. 2005). The angle between the shock front and the magnetic field requires the pre-shock field direction, which can be determined from streamer morphology. Severe elemental depletions are often observed in the closed field portions of streamers (Uzzo et al. 2004), providing an additional indicator of field topology. Another parameter potentially vital to the efficiency of shock acceleration is the density of suprathermal seed particles (e.g., Desai et al. 2003). While UVCS was not able to detect such particles, the improved sensitivity and instrumental profile characterization of next-generation instruments will make it possible to determine suprathermal particle densities out to $``$6 times the mean proton thermal speed (Cranmer 1998).
Although UVCS has proven the feasibility of detecting and characterizing CME shocks for several representative events, next-generation instrumentation can provide more extensive diagnostics and more complete spatial and temporal coverage. The empirical characterization of the coronal shock conditions and the ambient solar wind properties can then be used as inputs to: (1) 3D MHD models of the shock propagation through the heliosphere, and (2) multi-scale models of the SEP acceleration, transport, and energy spectrum synthesis (Zank et al. 2000; Li et al. 2003; Rice et al. 2003). Such constraints need to be applied in order to model specific events, such as the massive solar storms of OctโNov 2003. SEP acceleration and transport models can be tested for specific events by using spectroscopy and other remote-sensing diagnostics to constrain the initial parameters of the shock in the corona, then comparing with observed SEP energy spectra from, e.g., the Inner Heliospheric Sentinels. Iterative testing and refinement will ultimately result in a comprehensive validation of a predictive SEP model.
### 2.2 CME Current Sheets
Models of CMEs rely heavily on reconnection in current sheets, either trailing beneath the ejected magnetic flux rope or creating the flux rope in the first place (see Klimchuk 2001). Reconnection dumps large amounts of energy in the lower atmosphere of the Sun, creating intense heating, which accounts for the traditional flare ribbons and loops and for the current sheet containing hot plasma (Forbes et al. 1989; ล vestka & Cliver 1992; Forbes & Acton 1996; Priest & Forbes 2002).
UVCS observations made it possible for the first time to allow us to carry out diagnostics of the plasma inside the current sheet (Ciaravella et al. 2000; Ko et al. 2003). A narrow feature was seen in \[Fe XVIII\] emission in the space between the post-flare loops and the CME core, indicating electron temperatures near $`6\times 10^6`$ K. Significant progress in studying the current sheet and the process of magnetic reconnection in the current sheet was made recently when the UV spectral data of the plasma inside the current sheet during two events on January 8, 2002 and on November 18, 2003 were obtained and analyzed (Ko et al. 2003; Lin et al. 2005). Both events developed a fast CME, a growing flare loop system, and a long current sheet that connects CME and flare loops. In one of these events, the pattern of the reconnection inflow near the current sheet was well recorded in H I Ly$`\alpha `$. This allowed us to deduce the speed of the reconnection inflow directly and even to estimate the thickness of the current sheet. Also, it is possible to determine the magnetic field strength by using the observed speeds, densities, and temperatures to compute the kinetic and thermal energy densities in the reconnection region, then assume that this energy comes from the annihilation of an equal amount of magnetic energy just outside the current sheet (Ko et al. 2003). All of the parameters described above are required in order to put empirical constraints on the reconnection rate and electric field strength in the current sheet.
Such progress has very important theoretical consequences: for example, we are able to deduce the electrical resistivity (conductivity) in the current sheet, or in the reconnection region. Results obtained from UVCS and other remote-sensing instruments can provide the value of the electrical conductivity of the plasma inside the current sheet in an ongoing eruption for the first time since the impetus of applying reconnection theory to solar eruptions began six decades ago (Giovanelli 1946; and also see Priest & Forbes 2000).
With the knowledge of the dynamical process inside the current sheet, we are further able to investigate the particle acceleration taking place in the current sheet. A strong electric field is induced by magnetic reconnection in the current sheet. For a typical event, the electric field strength reaches about 5 V/cm (e.g., Wang et al. 2003; Qiu et al. 2004 for observations, and Martens & Kuin 1989; Forbes & Lin 2000; Lin 2002 for theories). An extremely high value of 50 V/cm was also reported (Xu et al. 2004). In principle, such a strong electric field is able to accelerate any charged particles. The current sheet is an assembly of waves and electric field, and accelerations can occur in various ways (see also Miller & Roberts 1995; Litvinenko 2000). We do not yet know exactly what happens in a real eruptive process, but the observations that a series of bright blobs flow successively out of the current sheet (e.g., Ko et al. 2003; Lin et al. 2005) replicate one of the main characteristics of magnetic reconnection inside the turbulent current sheet: the turbulent eddies or small magnetic islands inside the current sheet tend to merge into bigger ones before they leave the current sheet (Ambrosiano et al. 1988).
Further progress in understanding the above processes occurring in the current sheet depends on the accurate measurement of the thickness of the current sheet, plasma parameters in the current sheet (including electron and ion velocity distributions and densities), the speeds of reconnection inflow/outflow near the current sheet, as well as electric and magnetic fields in and around the current sheet. UV coronagraphic spectroscopy is uniquely suited to these requirements.
## 3 The Solar Wind
UVCS has led to fundamentally new views of the acceleration region of the solar wind. By measuring emission lines formed both by collisional excitation and by the resonant scattering of solar-disk photons, UV spectroscopy provides a multi-faceted characterization of the kinetic properties of atoms, ions, and electrons (e.g., Withbroe et al. 1982; Cranmer 2002a). The Doppler-broadened shapes of emission lines are direct probes of line-of-sight (LOS) particle velocity distributions (i.e., essentially providing $`T_{}`$ when the off-limb magnetic field is $``$radial), and red/blue Doppler shifts reveal bulk flows along the LOS. Integrated intensities of resonantly scattered lines can be used to constrain the solar wind velocity and other details about the velocity distribution in the radial direction (e.g., $`u_{}`$ and $`T_{}`$); this is the so-called โDoppler dimming/pumpingโ diagnostic (e.g., Noci et al. 1987). Intensities of collisionally dominated linesโespecially when combined into an emission measure distributionโcan constrain electron temperatures, densities, and elemental abundances in the coronal plasma. Even departures from Maxwellian and bi-Maxwellian velocity distributions are detectable with spectroscopic measurements having sufficient sensitivity and spectral resolution (e.g., Cranmer 2001).
In the high-speed solar wind, UVCS measured outflow speeds that were found to become supersonic much closer to the Sun than previously believed. In coronal holes, heavy ions (e.g., O<sup>+5</sup>) were found to flow faster, to be heated hundreds of times more strongly than protons and electrons, and to have anisotropic temperatures with $`T_{}>T_{}`$ (Kohl et al. 1997, 1998, 1999; Cranmer et al. 1999). These unexpected results have rekindled theoretical efforts to understand the heating and acceleration of the fast wind in the extended corona (e.g., Tu & Marsch 1997; Leer et al. 1998; Axford et al. 1999; Hollweg 1999; Hollweg & Isenberg 2002; Marsch 2004).
The slow solar wind was found to flow mostly along the outer edges of bright streamers, near locations with measured abundance patterns matching those of the in situ slow wind (Strachan et al. 2002; Raymond et al. 1997). The closed-field โcoreโ regions of streamers, though, exhibit heavy element abundances only 3 to 30% of those seen at 1 AU, indicating gravitational settling (e.g., Raymond 1999; Vรกsquez & Raymond 2005). UVCS observed the transition from a high-density collision-dominated plasma at low heights in streamers to a low-density collisionless plasma at large heights, the latter exhibiting high ion temperatures and anisotropies that suggest similar physics as in the fast wind (Ko et al. 2002; Frazin et al. 2003).
If the kinetic properties of additional ions were to be measured in the extended corona (i.e., a wider sampling of charge/mass combinations) we could much better constrain the specific kinds of waves that are present as well as the specific collisionless damping modes (e.g., Cranmer 2002b). Measuring the electron temperature above $``$1.5 $`R_{}`$ (never done directly before) would finally allow us to determine the bulk-plasma heating rate in different solar wind structures, thus putting the firmest ever constraints on models of why the slow \[fast\] speed wind is slow \[fast\] (e.g., Suess et al. 1999; Endeve et al. 2004; Cranmer, these proceedings). Measuring non-Maxwellian velocity distributions of electrons and positive ions would allow us to test specific models of MHD turbulence, cyclotron resonance, and velocity filtration. New capabilities such as these would be enabled by greater photon sensitivity, an expanded wavelength range, and the use of measurements that heretofore have only been utilized in a testing capacity (e.g., Thomson-scattered H I Ly$`\alpha `$ to obtain $`T_e`$; the Hanle effect to obtain constraints on the magnetic field). These would then allow the relative contributions of different physical processes to the heating and acceleration of all solar wind plasma components to be determined directly.
## 4 Implementation
A mission capable of carrying out the required observations would include two instrument units: a large-aperture ultraviolet coronagraph spectrometer (see Fig. 1) and a large-aperture visible light coronagraph. A suitable design was developed during a MIDEX Feasibility Study for a mission called ASCE. Remote external occulters supported by a deployable boom provide much larger unvignetted apertures than are possible with conventional coronagraphs. These instruments provide major improvements in sensitivity, stray light rejection, spatial resolution, minimum observable height and ultraviolet wavelength range. New spectroscopic diagnostics for the electron velocity distribution, magnetic field, and parameters for a broad range of newly observable ions are implemented by these instruments. Unprecedented cadences are possible.
This work is supported by NASA under grants NNG04GE77G, and NNG04GE84G to the Smithsonian Astrophysical Observatory. |
warning/0506/cond-mat0506458.html | ar5iv | text | # Electrodynamics of Nearly-Ferroelectric Superconductors
## I I. Introduction
This paper predicts novel effects in the electromagnetic response of a material which exhibits superconductivity (SC), and is in a โnearly-ferroelectricโ (NFE) state. We use the description โnearly-ferroelectricโ to mean a material which is a typical soft-mode phonon Cochran-Anderson system 1 , with high-static dielectric constant. Materials in this class include the sodium tungsten bronze Na<sub>x</sub>WO<sub>3</sub> and n- (or p-) doped SrTiO<sub>3</sub> systems.
Recently \[2a\] a sodium tungsten bronze Na<sub>x</sub>WO<sub>3</sub> with $`x0.05`$ was reported to be a high temperature superconductor with $`T_c90K`$. It is noteworthy that for $`0.1<x<1`$, Mattheis and collaborators long ago reported superconductivity but with $`T_c35K`$ \[2b\]. Also Mattheis and Wood \[2b\] first reported ferroelectricity in the host WO<sub>3</sub> system. Another material, n- (or p-) doped SrTiO<sub>3</sub> exhibits superconductivity, with $`T_c13K`$ 3 . The host SrTiO<sub>3</sub> is known to be a โnearly-ferroelectricโ material, where the static dielectric constant $`\epsilon (0)`$ has been measured as $`10^4`$ in the temperature range where superconductivity occurs 4 . In addition to these two materials, Weger and collaborators have recently proposed that in the high temperature cuprate superconductors, a new multicritical point occurs due to an underlying nearly-ferroelectric instability which renormalizes the electron-phonon and electron-electron interaction and enhances $`T_c`$ 5 ; 6 .
In the present work we will assume the โhostโ ionic lattice is a soft-mode NFE, with high dielectric coefficient $`\epsilon (0)`$. To clarify the scope of our work we also assume that all the carriers injected in the host by doping are fully condensed into $`s`$-wave Cooper pairs in that lattice, so there are no โfreeโ electrons. This assumption is consistent with earlier work on such systems by Cohen and Koonce 3 who presented a strong-coupling theory applied to superconductivity in SrTiO<sub>3</sub> (see also Ref. 4 ). Hence the free charge density $`\rho _e=0`$, and free current density $`๐_e=0`$.
In our present phenomenological approach, we solve the coupled equations for: a) the transverse electromagnetic field (Maxwell equations); b) soft-mode lattice vibrations (Born-Huang equations); c) the superconducting electrons in London (local) approximation. We report dispersive structure, when the electromagnetic frequency is tuned through the resonant frequency region near the basic $`\omega _{T0}`$ and $`\omega _{L0}`$ frequencies of the soft mode. The theory is applied to a semiinfinite medium, and to a film in order to calculate the reflectivity and transmittivity (impedance). We illustrate the results quantitatively, using experimental data for superconducting doped SrTiO$`_3,`$ which is one of the materials for which all the needed lattice, electronic and superconducting data are available.
We use London local electrodynamics for the superconducting sector(SC), rather than a nonlocal BCS or Pippard approach. These NFE-SC materials are all type II $`(\lambda _L>\xi `$ where $`\lambda _L`$ is the London penetration depth, and $`\xi `$ is the coherence length) and in some cases they are strongly type II $`(\lambda _L>>\xi )`$ . It is well known that London electrodynamics is valid in these cases. We also do not consider in the present work effects due to vortex creation in these systems.
## II II. Model And Dispersion
At the outset we note that all the relevant dynamical fields considered, are transverse. These are the electromagnetic fields, the lattice vibrations, and the London supercurrent. Maxwellโs equations for the macroscopic fields are:
$`๐=0,`$ $`๐=0,`$
$`\times ๐+{\displaystyle \frac{1}{c}}{\displaystyle \frac{๐}{t}}=0,`$ $`\times ๐{\displaystyle \frac{1}{c}}{\displaystyle \frac{๐}{t}}=0.`$ (1)
with the constitutive equations of the medium taken as
$$๐=๐+4\pi ๐\text{and}๐=๐+4\pi ๐=\epsilon (\omega )๐.$$
(2)
For the host ionic lattice we assume a โdiatomicโ basis, with $`๐ฐ`$ the relative displacement vector of the $`(\pm )`$ ions ($`๐ฐ=๐ฎ_+๐ฎ_{})`$. The equations of motion of the lattice, omitting damping, are 8 :
$$\frac{d^2๐ฐ}{dt^2}=b_{11}๐ฐ+b_{12}๐,๐=b_{21}๐ฐ+b_{22}๐,$$
(3)
where $`b_{ij}`$ are frequency-dependent coefficients and the electric field $`๐`$ and the polarization $`๐`$ are the same as in Eqns.(1) and (2). Now we seek a plane wave solution proportional to $`\mathrm{exp}(i๐ค๐ซi\omega t)`$ for the Eqs. (3) in the long-wave approximation. We follow Huang, and for these transverse waves, in the ionic medium with a single resonance frequency $`\omega _{T0}`$ we obtain the usual expression for the dielectric function
$$\epsilon (\omega )=\frac{\epsilon _{\mathrm{}}(\omega _{L0}^2\omega ^2)}{(\omega _{T0}^2\omega ^2)}$$
(4)
where the expressions for the coefficients $`b_{ij}`$ were given in Refs. 8 . In media such as the perovskites, bronzes, or cuprates we have to take account of the other (nonsoft) L0 and T0 modes by renormalizing $`\epsilon _{\mathrm{}}`$ to some โeffectiveโ value $`\epsilon _{\mathrm{}}^{}`$ 9 , where $`\epsilon _{\mathrm{}}^{}=\epsilon _{\mathrm{}}\mathrm{\Pi }_i^{}\omega _{Li}^2/\omega _{Ti}^2,`$ (the primed product over $`{}_{}{}^{\prime \prime }i_{}^{\prime \prime }`$ includes all oscillations except the soft mode) and we suppose the frequency $`\omega `$ is far below all frequencies $`\omega _{Li}`$ and $`\omega _{Ti}`$ except the soft mode. In what follows we use Eq. (4) with $`\epsilon _{\mathrm{}}`$ replaced by $`\epsilon _{\mathrm{}}^{}`$. The generalized Lyddane-Sachs-Teller (LST) relation is :
$$\frac{\epsilon _0}{\epsilon _{\mathrm{}}}=\frac{\mathrm{\Pi }_i\omega _{Li}^2}{\omega _{Ti}^2}.$$
(5)
Note that all oscillators, including $`i=0`$ (the soft mode) are included in Eq. (5) (no prime on the product). Finally we take the London equation 11 in the form:
$$\times ๐=\frac{1}{c}๐_s=\frac{1}{c^2\mathrm{\Lambda }}๐=\frac{n_se^2}{m^{}c^2}๐.$$
(6)
Here $`๐`$ is the vector potential for the electromagnetic field:
$$\times ๐=๐;\frac{1}{c}\frac{๐}{t}=๐.$$
(7)
The London gauge has been used, and in Eq. (6) $`n_s`$ is the superconducting electron density,and $`m^{}`$ is the electron effective mass. From Eqns. (1),(2),(4),(6) and (7) we obtain the equation for the electric field
$$\times \times ๐=\frac{1}{c^2}\epsilon (\omega )\frac{^2๐}{t^2}\frac{4\pi }{c^2\mathrm{\Lambda }}๐.$$
(8)
The field equations describe the coupled electromagnetic field, the soft mode NFE oscillator, and the London supercurrent.
We now look for a transverse (T) plane-wave solution for $`๐,`$ and recalling that $`\times \times ๐=^\mathrm{๐}๐`$ for such solutions, we obtain our new result, the dispersion equation for $`T`$ waves
$$k^2=\frac{\omega ^2}{c^2}\epsilon (\omega )\frac{4\pi }{c^2\mathrm{\Lambda }}.$$
(9)
If the London term (last on the right-hand side) is absent, then Eq. (9) is the usual equation for a transverse phonon polariton 8 ; on the other hand, if the displacement term due to the polarization is absent (first term on the right-hand side) then the field attenuates in the distance $`\lambda _L^1=\sqrt{4\pi /c^2\mathrm{\Lambda }}`$ where $`\lambda _L`$ is the London penetration depth 7 . Owing to the dispersive form (4) for $`\epsilon (\omega )`$, interesting physics emerges as $`\omega `$ is increased from zero. Depending on the sign of the right-hand side of Eq. (9), $`k`$ will either be real and electromagnetic waves will propagate, or imaginary so the wave will damp with frequency dependent penetration depth $`\lambda _L^{}`$ where
$$\lambda _L^{}=\lambda _L/\left(\left|1\epsilon (\omega )\lambda _L^2\omega ^2/c^2\right|\right)^{1/2}.$$
(10)
Figure 1 illustrates this. At $`\omega =0`$ the usual London penetration depth $`\lambda _L`$ is found. As the frequency increases $`\lambda _L^{}(\omega )`$ grows, up to a frequency $`\omega _{c1},`$ where $`k^2(\omega _{c1})=0`$ and $`\lambda _L^{}(\omega _{c1})`$ is infinite, so the field is uniform inside the medium. For $`\omega _{c1}<\omega <\omega _{T0},k^2>0`$ and the field propagates in the medium as in a dielectric. For $`\omega _{T0}<\omega <\omega _{L0},k^2<0`$ and again there is a frequency dependent $`\lambda _L^{}(\omega ),`$ i.e. Meissner effect. Exactly at $`\omega _{L0},\epsilon (\omega _{L0})=0,`$ so $`\lambda _L^{}(\omega _{LO})`$ coincides with the London penetration depth. For $`\omega _{L0}<\omega <\omega _{c2},k^2<0`$ giving increasing $`\lambda _L^{}(\omega )`$, and at $`k^2(\omega _{c2})=0,`$ a uniform field. Then, for $`\omega >\omega _{c2},k^2>0`$ and the wave propagates in the medium again, as in a dielectric.
In summary, our dispersion equation (9) as a function of frequency gives alternately regions of Meissner-like frequency-dependent penetration depth $`\lambda _L^{}(\omega )`$ that are superconducting with the magnetic field $`๐`$ excluded for $`x>\lambda _L^{},`$ then changing to regions of field propagation with $`\omega `$-dependent real wave number, as in a dielectric medium. This alternating dielectric and Meissner behavior suggests that our model has a zero-temperature phase transition between normal dielectric and superconductor phases driven by the electromagnetic frequency $`\omega ,`$ i.e. a type of โquantum phase transitionโ 12 .
Up to this point damping was ignored . To examine this we write the dielectric function (4) allowing for a phenomenological damping coefficient $`\mathrm{\Gamma }`$ (Ref. 13 ):
$$\epsilon (\omega )=\frac{\epsilon _{\mathrm{}}^{}(\omega _{L0}^2\omega ^2)}{(\omega _{T0}^2\omega ^2)2i\omega \mathrm{\Gamma }}.$$
(11)
As shown on Fig. 2, including โmoderateโ damping $`\mathrm{\Gamma }`$ of the soft-mode oscillator changes the results quantitatively but the above qualitative features remain the same.
Our first step is to determine the extent of the important frequency interval $`\omega _{c1}<\omega <\omega _{T0}.`$ The frequency $`\omega _{c1}`$ is the smallest root of the Eq. (9) for $`k=0,`$
$$\epsilon (\omega )=\frac{4\pi }{\mathrm{\Lambda }}\frac{1}{\omega ^2}.$$
(12)
The roots of this equation are,
$$\omega _{1,2}=\omega _{T0}\sqrt{\frac{\mathrm{}_1}{2}\frac{2\mathrm{\Gamma }^2}{\omega _{T0}^2(1+a^2)}\pm \frac{1}{2}\sqrt{\mathrm{}_2\frac{4\mathrm{\Gamma }^2\mathrm{}_1}{\omega _{T0}^2(1+a^2)}}}.$$
(13)
Here
$$\mathrm{}_1=1+\frac{1}{1+a^2},\mathrm{}_2=1\frac{1}{1+a^2},a^2=\frac{\omega _{L0}^2\epsilon _{\mathrm{}}^{}\mathrm{\Lambda }}{4\pi }$$
is a dimensionless parameter which can take on values of the order of unity in nearly ferroelectric materials such as n-SrTiO<sub>3</sub> (see below).
As we demonstrate below, when we calculate the actual size of the region $`\delta \omega =\omega _{TO}\omega _{c1}`$ for n-SrTiO<sub>3</sub> using measured parameters including $`\mathrm{\Gamma },\delta \omega `$ is, we believe, sufficiently wide to enable measurement of our predicted effects in a film of NFE-SC. First we consider a semi-infinite medium, then the film.
## III III. Semiinfinite medium
We consider a semiinfinite NFE-SC medium occupying the half-space $`z>0.`$ An incident electromagnetic wave propagates in the $`z`$ direction with electric and magnetic fields $`๐(๐ซ)`$ and $`๐(๐ซ)`$ polarized along $`x`$ and $`y`$ axes respectively. The Maxwell boundary conditions at the interface plane $`z=0`$ are 14
$`E_{ix}(0)+E_{rx}(0)E_{tx}(0)=0;`$
$`E_{ix}(0)E_{rx}(0){\displaystyle \frac{4\pi }{cZ}}E_{tx}(0)=0.`$ (14)
Here $`E_{ix}(0),E_{rx}(0),E_{tx}(0)`$ are the incident, reflected and transmitted electric fields, respectively and $`Z`$ is the surface impedance. First consider the case without damping $`(\mathrm{\Gamma }=0),`$ where we have
$$Z=\frac{4\pi }{c}\frac{E_{tx}(0)}{B_{ty}(0)}=\frac{8i\omega }{c^2}(\lambda _L^{})^2_0^{\mathrm{}}\frac{dk}{1\pm (\lambda _L^{}k)^2}.$$
(15)
We defined the effective penetration depth $`\lambda _L^{}`$ by Eq. (10). In the expression (15) the upper signs are used for frequency ranges where the material exhibits Meissner-like behavior $`(k^2<0),`$ while the lower signs apply to the frequency regime of dielectric-like behavior $`(k^2>0).`$ After integrating over $`k,`$ we obtain
$$Z=\frac{4\pi i\omega }{c^2}\lambda _L^{}\text{when}\omega <\omega _{c1}\text{or}\omega _{T0}<\omega <\omega _{c2};$$
and
$$Z=\frac{4\pi \omega }{c^2}\lambda _L^{}\text{when}\omega _{c1}<\omega <\omega _{T0},\text{or}\omega >\omega _{c2}.$$
(16)
Using these results we can solve the equations (14) and calculate the reflection coefficient for the Meissner-like $`(k^2<0)`$ or dielectric-like $`(k^2>0).`$ frequency ranges. We obtain the following:
$$R=\left|\frac{E_{rx}(0)}{E_{ix}(0)}\right|^2=\left(\frac{k_0\lambda _L^{}1}{k_0\lambda _L^{}+1}\right)^2;k_0=\frac{\omega }{c}$$
(17)
for the anomalous frequency ranges $`(\omega _{c1}<\omega <\omega _{T0};\omega >\omega _{c2})`$; and $`R1`$ when $`\omega <\omega _{c1}`$ or $`\omega _{T0}<\omega <\omega _{c2}.`$ Of course the latter confirms a well-known result of electrodynamics of superconductors, i.e., when we do not consider vortices, the electromagnetic field cannot penetrate into the volume of a superconducting medium. For the โanomalousโ intervals, e.g., $`\omega _{c1}<\omega <\omega _{T0}`$, the reflection coefficient tends to unity near the limiting points $`\omega _{c1}`$ and $`\omega _{T0}`$ because $`\lambda _L^{}`$ tends to infinity when $`\omega \omega _{c1}`$ and $`\lambda _L^{}`$ tends to zero when $`\omega \omega _{T0}.`$ However in the interior of this frequency range $`R`$ takes values which are considerably smaller than unity. This means that the electromagnetic field can penetrate into a semiinfinite NFE-SC medium within the โanomalousโ frequency range where the Meissner effect is suppressed. We did not include the effect of damping which will be discussed below. The damping effect is small in the practical case of n-SrTi$`O_3`$.
## IV IV. Thin film of NFE-SC
We turn to the case of a slab or film of thickness $`L`$ so that the material occupies the region $`0zL.`$ As previously, we take the electromagnetic field normally incident on the interface at $`z=0,`$ propagating in the $`z`$ direction with the electric and magnetic fields depending only on $`z`$ and parallel to $`x`$ and $`y`$ respectively. To proceed, we expand the electric field in the medium in a Fourier series, as done in Refs. 7 ,
$`E(z)`$ $`=`$ $`{\displaystyle \frac{2}{L}}{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}\left(1{\displaystyle \frac{1}{2}}\delta _{N0}\right)E_N\mathrm{cos}(k_Nz)`$ (18)
$`+`$ $`{\displaystyle \frac{2}{L}}{\displaystyle \underset{N=1}{}}\stackrel{~}{E}_N\mathrm{sin}(k_Nz).`$
Here
$`E_N`$ $`=`$ $`{\displaystyle _0^L}E(z)\mathrm{cos}(k_Nz)๐z;`$
$`\stackrel{~}{E}_N`$ $`=`$ $`{\displaystyle _0^L}E(z)\mathrm{sin}(k_Nz)๐z;k_N={\displaystyle \frac{\pi N}{L}}.`$
Using the dispersion equation (9), we get
$`E_N`$ $`=`$ $`{\displaystyle \frac{\lambda _L^2}{\lambda _L^2k_N^21}}\left[E^{}(0)(1)^NE^{}(L)\right];`$
$`\stackrel{~}{E}_N`$ $`=`$ $`{\displaystyle \frac{\lambda _L^2}{k_N(\lambda _L^2k_N^21)}}\left[E(0)(1)^NE(L)\right].`$ (19)
The upper sign in the denominators of Eq. (19) is used for those frequency ranges where $`k^2<0,`$ and the lower sign corresponds to the frequency ranges where $`k^2>0.`$
Since the new results emerge in the region $`\delta \omega =\omega _{C1}\omega _{TO}`$, when $`k^2>0`$ we consider this case in this Section \[lower sign in (19)\], and for the present we take $`\mathrm{\Gamma }=0`$. In this frequency range the material exhibits dielectriclike behavior. Substituting Eq. (19) into Eq. (18) we get
$$E_x(0)=\frac{i\omega \lambda _L^{}}{c}B_y(0)\mathrm{cot}\left(\frac{L}{\lambda _L^{}}\right)\frac{i\omega \lambda _L^{}}{c}B_y(L)\mathrm{sin}^1\left(\frac{L}{\lambda _L^{}}\right);$$
$$E_x(L)=\frac{i\omega \lambda _L^{}}{c}B_y(0)\mathrm{sin}^1\left(\frac{L}{\lambda _L^{}}\right)\frac{i\omega \lambda _L^{}}{c}B_y(L)\mathrm{cot}\left(\frac{L}{\lambda _L^{}}\right).$$
(20)
The solution of this system of equations for the magnetic induction is:
$$B_y(0)=\frac{ic}{\omega \lambda _L^{}}E_x(0)\mathrm{cot}\left(\frac{L}{\lambda _L^{}}\right)\frac{ic}{\omega \lambda _L^{}}E_x(L)\mathrm{sin}^1\left(\frac{L}{\lambda _L^{}}\right);$$
$$B_y(L)=\frac{ic}{\omega \lambda _L^{}}E_x(0)\mathrm{sin}^1\left(\frac{L}{\lambda _L^{}}\right)\frac{ic}{\omega \lambda _L^{}}E_x(L)\mathrm{cot}\left(\frac{L}{\lambda _L^{}}\right).$$
(21)
The boundary conditions on the interfaces $`z=0`$ and $`z=L`$ can be written as follows:
$`E_{ix}(0)+E_{rx}(0)=E_{tx}(0);`$
$`E_{ix}(0)E_{rx}(0)=B_y(0);`$
$`E_x(L)=B_y(L).`$ (22)
Using these boundary conditions and expressions (21) for the magnitude of the magnetic field on the interfaces $`z=0`$ and $`z=L`$ we arrive at the following expressions for the reflection and transmission coefficients:
$`R={\displaystyle \frac{1}{1+\rho ^2(\omega )\mathrm{sin}^2(L/\lambda _L^{})}};`$ (23)
$`T={\displaystyle \frac{\rho ^2(\omega )\mathrm{sin}^2(L/\lambda _L^{})}{1+\rho ^2(\omega )\mathrm{sin}^2(L/\lambda _L^{})}},`$ (24)
where the the function $`\rho (\omega )`$ has the form:
$$\rho (\omega )=\frac{2ck^{}}{\omega }\frac{1}{1c^2k^2/\omega ^2}.$$
(25)
Here $`k^{}=1/\lambda _L^{}`$ is the solution of the dispersion equation (9) in the โanomalousโ frequency range $`\omega _{C1}<\omega <\omega _{TO}`$. For frequencies very close to $`\omega _{T0}`$ we have $`k(\omega /c)\sqrt{\epsilon }`$ which corresponds to a polariton-phonon transverse wave.
We now note an interesting result following from Eqs. (23) and (24), when the effective London penetration depth $`\lambda _L^{}`$, and the Fourier component $`k=k_N`$ are in resonance. That is, for $`\lambda _L^{}=L/(\pi N)`$ or $`k^=k_N,\mathrm{sin}(L/\lambda ^{})`$ in Eqs. (23) and (24) equals zero and we get for the reflection and transmission coefficients the values $`R=0;T=1,`$ which means total transparency of the slab at $`k_N`$ (See Fig. 3).
Total transparency can also occur at a frequency for which $`k=k_0=\omega /c`$. At this frequency we have $`\epsilon (\omega )4\pi /\omega ^2\mathrm{\Lambda }=1`$ so it is located close to the left boundary frequency $`\omega _{c1}.`$ When $`k_0=k`$ the function $`\rho (\omega )`$ tends to infinity and again we get $`R=0`$ and $`T=1.`$ To explain this we can treat the quantity $`\sqrt{\stackrel{~}{\epsilon }(\omega )}=\sqrt{\epsilon (\omega )4\pi /\omega ^2\mathrm{\Lambda }}\stackrel{~}{n}(\omega )`$ as an effective reflective index of the slab material. When it equals unity the electromagnetic field from outside passes through the slab without reflection.
Our results (23) and (24) were obtained, neglecting damping. For a nonzero damping constant $`\mathrm{\Gamma }`$ we have to replace our expressions (23) and (24) near the resonances by the following:
$`R={\displaystyle \frac{1}{1+\rho ^2(\omega )+S^2(\omega )\mathrm{sin}^2(L/\lambda {}_{L}{}^{}^{})}};`$ (26)
$`T={\displaystyle \frac{\rho ^2(\omega )\mathrm{sin}^2(L/\lambda {}_{L}{}^{}^{})}{1+\rho ^2(\omega )+S^2(\omega )\mathrm{sin}^2(L/\lambda {}_{L}{}^{}^{})}},`$ (27)
where
$$S(\omega )=\rho (\omega )\mathrm{cos}\left(\frac{L}{\lambda {}_{L}{}^{}^{}}\right)+\sqrt{1+\rho ^2(\omega )}\mathrm{sinh}\left(\frac{L}{\lambda {}_{L}{}^{}^{\prime \prime }}\right)$$
(28)
and $`1/\lambda _L^{}=1/\lambda {}_{L}{}^{}^{}+i/\lambda {}_{L}{}^{}^{\prime \prime }.`$
We can use the following approximation for $`\lambda {}_{L}{}^{}^{}`$ and $`\lambda {}_{L}{}^{}^{\prime \prime }:`$
$`{\displaystyle \frac{1}{\lambda {}_{L}{}^{}^{}}}\sqrt{{\displaystyle \frac{\omega ^2}{c^2}}\epsilon ^{}(\omega ){\displaystyle \frac{1}{\lambda _L^2}}};`$ (29)
$`{\displaystyle \frac{1}{\lambda {}_{L}{}^{}^{\prime \prime }}}\left({\displaystyle \frac{2\omega \mathrm{\Gamma }}{\omega _{T0}^2\omega ^2}}\right)^{1/2}{\displaystyle \frac{1}{\lambda _L}}.`$ (30)
It was shown before that damping in the lattice system reduces the anomalous frequency range $`\delta \omega .`$ The frequency interval $`\omega _{c1}<\omega <\omega _{T0}`$ has to be replaced by the narrower interval $`\omega _1<\omega <\omega _2.`$ The imaginary part of $`1/\lambda _L^{}`$ depends on frequency and increases when the latter increases. As we will illustrate below, moderate damping, as in n-SrTiO<sub>3</sub> will only slightly modify the sharp resonances which occur at $`k=k_N`$ and $`k=k_0.`$
To the best of our knowledge the general behavior we predict for an NFE-SC system has not been observed. In order to examine the practicality of making such observations we have calculated the effects quantitatively for n-doped SrTiO<sub>3</sub>. This material was much studied, as a nearly-ferroelectric, and as a superconductor, prior to the era of high T<sub>c</sub>, so that the needed physical parameters are available.
## V V. Application to n-SrTiO<sub>3</sub>, a Prototype NFE-SC
In order to evaluate the magnitudes and locations of the predicted effects, we require lattice, dielectric, and electronic data. So far as we can determine such measured data is only available for n-SrTiO<sub>3</sub>, and we will use it below. Nearly ferroelectric SrTiO<sub>3</sub> has been studied for some time. The temperature dependent static dielectric coefficient $`\epsilon (0)`$ has been measured in a wide temperature range, extending down to about 1K . The frequency dependent $`\epsilon _{\mathrm{}}`$ is known at various temperatures . Lattice dynamic studies were carried out based on the shell model , and neutron scattering , and infra-red and Raman (including electric field dependence) optical properties have been measured . As a result the lattice dynamical parameters are known, and the generalized LST relation (5) has been verified. From the line-widths in infra-red and Raman properties , the damping parameter $`\mathrm{\Gamma }`$ for the soft mode lattice frequency $`\omega _{TO}`$ is known. Hence all the physical parameters of the dielectric host are actually known. The temperature independent quantities are
$$\epsilon _{\mathrm{}}=5.5,\omega _{LO}=5.2\times 10^{12}s^1,\frac{\mathrm{\Pi }_i^{}\omega _{Li}^2}{\omega _{Ti}^2}=4.1.$$
Using the LST relation we have at $``$1K
$$\omega _{TO}=1.6\times 10^{11}s^1$$
and at 1K
$$\epsilon (0)2.25\times 10^4.$$
The damping constant for the soft TO mode as in eqn(13) is measured from optical line width 9 :
$$\mathrm{\Gamma }^2/\omega _{T0}^20.2$$
at the temperature of interest. (Actually the measurement was at a higher temperature than the region of T<sub>c</sub> so the relevant $`\mathrm{\Gamma }`$ for our work will be smaller, but we use the above value.)
For the needed electronic parameters we use the effective mass $`m^{}10m_e`$ where $`m_e`$ is the free electron mass. A range of $`n`$ doping from $`10^{17}`$ to $`10^{21}`$ cm<sup>-3</sup> was used for SC n-SrTiO<sub>3</sub>. We choose $`n_s=9\times 10^{17}`$ cm<sup>-3</sup> for our calculation. We recall that n-doped SrTiO<sub>3</sub> exhibits concentration dependent T<sub>c</sub>, with T<sub>c</sub> = 0.3K in the optimum doped system 3 . In this temperature range $`\epsilon (0)`$ has the high value given above, as was noted also by Cohen 3 .
We now determine the anomalous frequency interval just below $`\omega _{TO}`$. We first need the parameter $`a^2=\omega _{LO}^2\epsilon _{\mathrm{}}^{}\mathrm{\Lambda }/4\pi `$ \[see Eq. (13)\]. For SrTiO<sub>3</sub> with $`n_s=9\times 10^{17}`$ cm<sup>-3</sup>, and $`m^{}10m_e`$, we get $`a^22.1`$. Taking, first, $`\mathrm{\Gamma }=0`$, and solving Eq. (12) we find $`\omega _{c1}=0.9\times 10^{11}c^1`$. This gives $`\delta \omega =\omega _{TO}\omega _{c1}=0.43\omega _{TO}.`$ Using the measured $`\mathrm{\Gamma }`$ for SrTiO<sub>3</sub>, we obtain that $`\omega _2`$ shifts negligibly while $`\omega _{c1}`$ shifts slightly to $`\omega _{c1}^{}=0.96\times 10^{11}c^1`$, and the anomalous frequency region becomes slightly smaller
$$\delta \omega ^{}=\omega _{TO}\omega _{c1}^{}0.4\omega _{TO}.$$
We expect that this interval $`\delta \omega ^{}`$ is still sufficiently wide to examine optical response experimentally by reflection/transmission studies in n-SrTiO<sub>3</sub>. The small effect of damping makes it likely also that $`\mathrm{\Gamma }=0`$ is a good approximation for the calculations that follow.
We now turn to calculating the reflection/transmission coefficients for a $`n`$-SrTiO<sub>3</sub> film of NFE-SC of thickness $`L`$, first for $`\mathrm{\Gamma }=0`$, using the formulas (23)-(25). In this case the results are shown on Fig. 3. We predict a โcomb-likeโ structure of narrow perfect transmission $`T_\rho =1`$ spikes alternating with regions where $`T_\rho =0.`$ The spikes occur at the resonance condition where $`\lambda _L^{}=L/\pi N`$. Next we estimate the effect including damping. For weak damping $`\mathrm{\Gamma }^2/\omega _{TO}^2=0.2`$ the value of the factor $`\left[2\omega \mathrm{\Gamma }/(\omega _{T0}^2\omega ^2)\right]^{1/2}`$ changes from 0.9 to 1.5 when $`\omega `$ runs over the interval between $`\omega _1`$ and $`\omega _2.`$ Thus, when the thickness of our film $`L`$ is small or of the same order as the $`\omega =0`$ London penetration depth $`\lambda _L`$ (here $`\lambda _L1.8\times 10^3`$ cm) moderately weak damping as in SrTiO<sub>3</sub> cannot significantly change the height of the peak in the transmission coefficient. Peaks arise due to the propagation of the transverse mode corresponding to the solution of the dispersion equation (9). We believe that these peaks can be observed in any NFE-SC material in the appropriate frequency ranges, and at temperatures below $`T_c.`$ Such an observation would give support for theory proposed here.
It is not possible to make realistic estimates of the interval $`\delta \omega `$ (or $`\delta \omega ^{}`$ with damping) for Na<sub>x</sub>WO<sub>3</sub>, or for cuprates pending the availability of the necessary material constants as we had in $`n`$-SrTiO<sub>3</sub>. The theory presented will apply to any such NFE-SC material, so the experimental search for the resonance comb-like structure in T($`\omega `$) just below $`\omega =\omega _{TO}`$ will be highly valuable.
## VI VI. Summary
We solved the coupled โsoft-modeโ lattice dynamics (long wavelength) equations of motion, together with Maxwell equations for the associated electromagnetic field and the London equations for the supercurrent, in a NFE-SC. The analysis can apply to several material systems of current interest: sodium tungsten bronzes Na<sub>x</sub>WO<sub>3</sub>, doped $`n`$-SrTiO<sub>3</sub>, and possibly high temperature cuprates. We note that several authors have discussed microscopic theories for the superconductivity in $`n`$-SrTiO<sub>3</sub> (Ref. 15 ) in the strong-coupling framework. However, our macroscopic approach has not previously been reported, to our knowledge. We assume all the free carriers are condensed in Cooper pairs, as in some earlier work on the SrTiO<sub>3</sub> (See Regs. 3 ; 4 ). The resulting coupled modes can be considered as phonon-polaritons dressed by the supercurrent or electromagnetic waves in a London superconductor dressed by the transverse optic waves (TO phonons).
A result, which we illustrated, is the alternation of the system response between Meissner-like superconductor and a dielectric-like medium, as the incident electromagnetic frequency is continuously varied. Particularly of interest is the โcomb-likeโ series of transmission resonance peaks we predict, which should be measurable. Such measurements would test our theory. One NFE-SC material of choice would be $`n`$-SrTiO<sub>3</sub> for which we made quantitative estimations.
## VII Acknowledgments
We thank G.M. Zimbovsky for his help with the manuscript. An inservice grant from the PSC-CUNY is also acknowledged for partial support.
## VIII Appendix: Partition of Energy Between Radiation, Lattice Polarization, and London Field
It is useful to obtain the frequency-dependent partition of energy between the three propagating fields. This can be seen as extending Huangโs early treatment of the radiation-lattice coupling (phonon-polariton 8 ) to include the London supercurrent, or conversely extending London treatment 16 of the energy momentum theorem, to include the dielectric polarization.
Within our treatment all fields are transverse and we consider the infinite medium. We will obtain the continuity equation for energy density $`U`$ and Poynting energy flux vector $`๐ฌ.`$ From Maxwellโs equations (1) we have:
$$๐ฌ=\frac{1}{4\pi }๐\frac{๐}{dt}\frac{1}{4\pi }๐\frac{๐}{dt},$$
$`(A.1)`$
where
$$๐ฌ=\frac{๐}{\mathrm{๐}\pi }[๐\times ๐].$$
$`(๐\mathbf{.2})`$
In order to simplify the right hand side of (A1) we use the lattice equation of motion (3) to obtain:
$$\frac{1}{c}๐\frac{๐}{dt}=\frac{1}{2c}\epsilon _{\mathrm{}}\frac{}{t}(๐^2)+\frac{2\pi }{c}\frac{}{t}\left(\frac{}{t}๐ฐ^2+\omega _{T0}^2๐ฐ^2\right).$$
$`(A.3)`$
Then, using the first of equations (2), and equations (6) and (7) we have
$$[\times ๐]=\frac{1}{c}๐=\frac{1}{c^2\mathrm{\Lambda }}๐$$
$`(๐\mathbf{.4})`$
and
$$[\times [\times ๐]]=^2๐=\frac{1}{c^2\mathrm{\Lambda }}๐.$$
$`(๐\mathbf{.5})`$
Now we will assume plane monochromatic electromagnetic wave propagation $`(\mathrm{exp}(i๐ค๐ซi\omega t)).`$ So
$$^2k^2;\frac{}{t}i\omega .$$
Then
$$k^2๐=\frac{1}{c^2\mathrm{\Lambda }}๐,$$
$`(A.6)`$
so
$$๐=๐\frac{1}{\lambda _L^2k^2}๐,$$
$`(๐\mathbf{.7})`$
or with $`๐=\mu ๐`$ we get:
$$\frac{1}{\mu }=1+\frac{1}{\lambda _L^2k^2}$$
$`(A.8)`$
and
$$\frac{1}{c}๐\frac{๐}{t}=\frac{1}{c}\frac{1}{\mu }๐\frac{๐}{t}=\frac{1}{2c}\left(1+\frac{1}{\lambda _L^2k^2}\right)\frac{}{t}๐^2.$$
$`(A.9)`$
Since all waves are transverse take $`๐=(E_x,0,0);๐=(0,B_y,0);๐ค=(0,0,k)`$ giving $`B_y=(ck/\omega )E_x.`$ We eliminate $`๐`$ from (A.1), and combining all terms, we get
$$\frac{U}{dt}=\frac{1}{4\pi }๐\frac{๐}{t}+\frac{1}{4\pi }๐\frac{๐}{t}=\frac{1}{8\pi }\left(\epsilon _{\mathrm{}}\frac{}{t}๐^2+\frac{๐^2}{t}\right)$$
$$\frac{1}{2}\frac{}{t}\dot{๐ฐ}^2+\frac{\omega _{T0}^2}{2}\frac{}{t}๐ฐ^2+\frac{1}{8\pi }\frac{c^2}{\omega ^2\lambda _L^2}\frac{}{t}๐^2.$$
$`(A.10)`$
It is simple to identify the various terms in the energy density $`U`$ as belonging to radiation field, lattice polarization, and London supercurrent. Thus
$$U_{RAD}=\frac{1}{8\pi }\epsilon _{\mathrm{}}๐^2+\frac{1}{8\pi }๐^2,$$
$`(A.11)`$
$$U_{LATTICE}=\frac{1}{2}\dot{๐ฐ}^2+\frac{1}{2}\omega _{TO}^2๐ฐ^2,$$
$`(A.12)`$
$$U_{SC}=\frac{1}{8\pi }\frac{1}{\lambda _L^2k_0^2}๐^2,$$
$`(A.13)`$
where $`k_0=\omega /c`$. All these contributions to the energy density can be expressed in terms of $`๐^2`$ in our plane-wave case, giving
$$U_{RAD}=\frac{1}{4\pi }\left(\epsilon _{\mathrm{}}+\frac{k^2}{k_0^2}\right)๐^2,$$
$`(A.14)`$
$$U_{LATTICE}=\frac{(\epsilon _0\epsilon _{\mathrm{}})\omega _{TO}^2(\omega ^2+\omega _{TO}^2)}{8\pi (\omega _{TO}^2\omega ^2)^2}๐^2$$
$`(A.15)`$
and, as before:
$$U_{SC}=\frac{1}{8\pi }\frac{1}{\lambda _L^2k_0^2}๐^2.$$
$`(A.16)`$
Before proceeding we note that the expression for the lattice term in energy density is unchanged from that found by Huang 8 for the phonon-polariton. Formally that is also true for the radiative term, however hidden in $`U_{RAD}`$ is the effect of the London term via the changed dispersion equation (9), which includes $`\mathrm{\Lambda }`$. Note that substituting the dispersion equation (9) into $`U_{RAD}`$ we get
$$U_{RAD}=\frac{1}{8\pi }\left(\epsilon _{\mathrm{}}+\epsilon (\omega )\frac{1}{\lambda _L^2k_0^2}\right)๐^2$$
$`(A.17)`$
or
$$U_{RAD}=\frac{1}{8\pi }\left(2\epsilon _{\mathrm{}}+\frac{\epsilon _{\mathrm{}}(\omega _{LO}^2\omega _{TO}^2)}{\omega _{TO}^2\omega ^2}\frac{1}{\lambda _L^2k_0^2}\right)๐^2,$$
$`(A.18)`$
and
$$U_{SC}=\frac{1}{8\pi }\frac{1}{\lambda _L^2k_0^2}๐^2.$$
$`(A.19)`$
We may see that the energy-density contribution $`U_{SC}`$ will exactly cancel the London contribution to the radiation energy density $`U_{RAD}`$. Thus, if we calculate the lattice (polarization) fraction of the energy we obtain exactly the result of Huang 8 despite the coupling of the London supercurrent. We have
$`\rho _{LATTICE}`$ $`=`$ $`{\displaystyle \frac{U_{LATTICE}}{U_{TOTAL}}}={\displaystyle \frac{U_{LATTICE}}{U_{LATTICE}+U_{RAD}+U_{SC}}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(\omega _{TO}^2+\omega ^2)(\omega _{LO}^2\omega _{TO}^2)}{(\omega _{TO}^2\omega ^2)^2+\omega _{TO}^2(\omega _{LO}^2\omega ^2)}}.(A.20)`$
This result shows that the mode described by the dispersion equation (9) is a mixed photon-phonon mode within all frequency range $`\omega _{c1}<\omega <\omega _{TO}`$. We emphasize this result is only meaningful in the intervals where $`k^2>0`$. The lattice contribution to the energy enhances when the frequency increases and gets its maximum when $`\omega `$ goes to $`\omega _{TO}`$. At the frequencies close to $`\omega _{TO}`$ $`\rho 1`$, so all the energy density is concentrated in the ion system. Taking into account that the frequency $`\omega _{TO}`$ is a boundary frequency for the โanomalousโ frequency range, we can expect that the ion oscillations at the frequency $`\omega _{TO}`$ being excited by the electromagnetic field which can penetrate to a NFE-SC material within the โanomalousโ frequency interval, will exist at frequencies larger than $`\omega _{TO}`$ where the material behaves as a superconductor. |
warning/0506/quant-ph0506129.html | ar5iv | text | # A GENERAL RELATIVISTIC GENERALIZATION OF BELL INEQUALITY
## 1 References
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warning/0506/astro-ph0506109.html | ar5iv | text | # The nature of Composite Seyfert/Star-forming galaxies revealed by X-ray observations
## 1. Introduction
Evidence for a link between intense star formation and nuclear activity has grown steadily in recent years (e.g. Veilleux, 2001, Gonzales Delgado et al. 2001, and references therein). On one hand, the presence of circumnuclear starbursts in many local AGNs (e.g., Levenson et al. 2004, Levenson et al. 2001) suggests a connection not yet fully understood. The presence of a starburst has been invoked for instance to produce absorption in low luminosity AGNs (Fabian et al. 1998, Ohsuga & Umemura, 2001). On the other hand, the ubiquity of supermassive black holes in the nuclei of normal galaxies (Kormendy et al. 2002) and the proportionality between the black hole and the spheroidal masses (Ferrarese & Merrit 2000) evidence a direct link between the formation of ellipticals and spiral bulges and the growth of central black holes. Therefore, the interplay between accretion on supermassive black holes and galaxy formation and evolution has become a fundamental ingredient for theoretical models in this field (e.g., Springel et al. 2005, Sazonov et al. 2005, Nipoti et al. 2003). This implies that our view on the star formation history of the universe, as deduced by galaxy luminosity functions (see Springel et al. 2005), as well as on the chemical enrichment and feedback processes in the early universe, might change in a way that we cannot foresee at this moment. Understanding the connection between starburst and AGN in the local universe is therefore of crucial importance.
A spectroscopic optical survey of bright IRAS and X-ray selected sources from the ROSAT All Sky Survey revealed a small enigmatic class of low redshift galaxies with optical spectra dominated by the features of H II galaxies but presenting very subtle Seyfert signatures as well. These objects had X-ray luminosities typical of broad line AGNs, ranging from 1.5 $`\times `$ 10<sup>42</sup> erg/s to 5 $`\times `$ 10<sup>43</sup> erg/s in the ROSAT band (Moran et al. 1996, hereafter M96). They were named Composite after their Composite optical spectra. The diagnostic emission line ratio diagrams (Veilleux & Osterbrock 1987) classify them either on the boundary between Seyfert and HII regions or as pure star-forming galaxies. Most of these galaxies show quite narrow emission lines (FWHM$`<`$300 km/s) as HII galaxies normally do. Yet, some of them present \[O III\]$`\lambda `$$`\lambda `$4959,5007 lines significantly broader than all other narrow lines in the spectrum and weak and elusive broad H$`\alpha `$ wings. A possible scenario to explain the optical and X-ray mismatch proposed by M96 invokes absorption as responsible for the obscuration of the optical features. M96 also suggested that it is not likely that the starburst could overpower a Seyfert optical nuclear spectrum, given that the starburst component in these objects is not particularly strong. On the contrary here we propose that the AGN and the starburst components in these objects are present at the same level of activity making them ideal laboratories in which the two phenomena can be studied together.
Until recently little or no information in the X-ray spectral domain were available for the M96 sample of Composite galaxies except for the ROSAT data. In the last couple of years, Chandra and XMM-Newton observations of IRAS 00317$``$2142 and IRAS 20051$``$1117, two of the brightest objects of the sample, have revealed the presence of an active nucleus with a spectral shape typical of type 1 AGNโs (a photon index $`\mathrm{\Gamma }`$ $``$ 1.8-1.9) and no absorption. Long term variability has been observed in IRAS 00317$``$2142 (Georgantopoulos, Zezas & Ward 2003), while no flux variation has been measured in IRAS 20051$``$1117 (Georgantopoulos et al. 2004).
In order to acquire the missing information in X-rays for the rest of the sample, and to derive the overall properties for the whole sample, we first proposed to observe these objects with BeppoSAX (Boella et al. 1997) and we obtained data for 3 sources of the original list of 6 of Moran et al. (1996). Then we observed four objects of the sample with Chandra (Weisskopf et al. 2000). Here we present the results obtained from the Chandra and BeppoSAX observations and discuss possible explanations for the enigmatic behaviour of this class of objects.
The paper is organized as follows: the sample is briefly outlined in Section 2, new X-ray observations are described in Sections 3 through 6 and discussed together with a multi-frequency analysis in Section 7. Our conclusions are drawn in Section 8. Throughout the paper we assume H$`{}_{0}{}^{}=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup>.
## 2. The Composite sample
We list the 6 objects defined as Composite by M96 in Table 1. We note that the original M96โs list included 7 objects, but IRAS 10113+1736 is no longer a valid Composite, since it has become evident that infrared and X-ray emissions originate from different sources (Condon et al. 1998a).
Table 1 lists all the sample properties, i.e. J2000 position, redshift, Hubble type (from LEDA<sup>1</sup><sup>1</sup>1http://leda.univ-lyon1.fr/ , Paturel et al. 1997), Galactic N<sub>H</sub> from Dickey & Lockman (1990), IRAS far-infrared luminosities; we also report the logarithmic luminosities of the H$`\alpha `$ (narrow+broad and only broad) and \[O III\]$`\lambda `$5007 emission lines obtained by M96 in a $`2^{\prime \prime }\times 4^{\prime \prime }`$ aperture.
We note that IRAS 00317$``$2142 belong to the compact group HCG 4. However, no extended emission is detected in the group at a limit of 10<sup>41</sup> erg s<sup>-1</sup> in the ROSAT PSPC (Mulchaey et al. 2003), so that no group emission is expected to contaminate our results.
The galaxies are local (z $`<`$ 0.04). They are all detected by the NVSS (Condon et al. 1998b) in the range 6-44 mJy and by the 2MASS (Cutri et al. 2003) with K<sub>s</sub> mag between 10.3 and 11.6. Their optical magnitudes range between 14 and 16 (from NED).
For completeness we include the results of the ASCA, XMM-Newton and Chandra observations of IRAS 00317$``$2142 and IRAS 20051$``$1117 (Georgantopoulos 2000, Georgantopoulos, Zezas & Ward 2003, Georgantopoulos et al. 2004) in the X-ray table results (see Sect. 5) and in the final discussion of the whole sample.
## 3. X-ray data
### 3.1. BeppoSAX Observations
We observed three Composite objects with the Narrow Field Instruments (NFI) of the BeppoSAX satellite. The observation dates and the total effective exposure time in ks for the LECS (0.1-4 keV), MECS (1.3-10 keV) and PDS (13-200 keV) instruments are listed in the bottom part of Table LABEL:infobs. All objects have been detected by the MECS instrument, none by the LECS detector. The only source detected by the PDS is IRAS 20051$``$1117 (S/N $``$$`>`$ 3) with a total count rate of (13.73$`\pm `$4.48)$`\times `$10<sup>-2</sup> cts/s in the 20-200 keV band; however the PDS data fall above the simple power-law extrapolation from the MECS data suggesting a possible contamination from a different object as discussed in Sec 4.
The extraction of the source and background spectra was done with the XSELECT package. The MECS spectra have been extracted from a 4 arcmin radii circular regions; the background for each source has been extracted from source free off-axis circular regions and subtracted. The redistribution matrices and ancillary response files released in September 1997 have been applied. MECS net count rates and errors are listed in Table LABEL:infobs.
### 3.2. Chandra Observations
We observed four objects with the Advanced CCD Imaging Spectrometer ACIS-I on board Chandra between March and September 2003 using exposures of $``$ 25 ks. Details of the observations can be found in the top part of Table LABEL:infobs. The data reduction was performed using the CIAO version 3.1 software. Level 1 event files have been created applying the time-dependent gain correction. The standard procedure provided by the CIAO โthreadsโ has been followed in order to obtain level 2 event files. Background light curves have been examined and in two cases (IRAS 01072+4954 and IRAS 01319$``$1604) background flares have been found and removed. In all cases, โgoodโ exposure times (listed in column 3 of Table LABEL:infobs) are always $``$ 90% of the total observation times. Based on the RASS fluxes, we had requested observations in 1/4 sub-array mode in order to minimize pile-up effects. However, we found that the observed count rates are below the expected values, therefore no pile-up affects the observations.
Nuclear spectra have been extracted from a circular region of a 3<sup>โฒโฒ</sup> radius; the background has been evaluated in nearby source free regions and, as shown by the errors on the ACIS-I net count rates (Table LABEL:infobs), it contributes less than few percent to the total count rate.
## 4. Chandra Imaging Analysis
We have run a detection algorithm (wavdetect in CIAO) in fields of view of $``$5$`\times `$2 of the ACIS-I3. Point-like emissions coincident with the optical position of the nucleus have been detected in all objects, confirming the X-ray/optical association.
We list the off-nuclear sources detected within the optical extent (D25) of each galaxy in Table LABEL:off. We give the positions in RA & DEC (J2000), the distance from the nucleus in arcsec and kpc, the number of net counts and relative errors, the source significance and the 2-10 keV fluxes and luminosities computed assuming a $`\mathrm{\Gamma }`$=1.8 power-law spectrum modified by Galactic column densities, and the redshift of the corresponding galaxy. We do not detect any off-nuclear source within the optical extent of IRAS 01072+4954 (at the given exposure the flux limit corresponds to a luminosity of 2.5$`\times `$10<sup>39</sup> erg/s). All the off-nuclear sources are significantly detected at $`>`$ 3$`\sigma `$. If they are actually associated with the host galaxies, the observed luminosities would be well in excess of the Eddington luminosity for a solar-mass black hole or neutron star X-ray binary, therefore they should be considered as Ultra-Luminous X-ray objects (ULXs, Swartz et al. 2004). Note that despite the presence of these luminous sources, the galaxy X-ray emission is completely dominated by the nucleus in each case.
Given the large RASS point spread function of about 45<sup>โฒโฒ</sup> even for bright sources and the positional accuracy of 25<sup>โฒโฒ</sup> (90% confidence level, Voges et al. 1999), other sources could in principle contaminate the RASS detection. To verify whether any other X-ray source was present within the ROSAT aperture we examined the $`2^{}`$ region around each source. We have found a source with $`50`$ cts at 121<sup>โฒโฒ</sup> from the nucleus of IRAS 01072+4954, a source with $`15`$ cts at 55<sup>โฒโฒ</sup> from IRAS 01319$``$1604 and a source with $``$ 45 cts at 79.5<sup>โฒโฒ</sup> from IRAS 04392$``$0123. In the present data, these sources contribute at most 15% of the nuclear flux, so that it is unlikely that they could have been dominant in the ROSAT observations.
In order to look for extended emission, the radial profile of each source has been derived in annular regions centered on the nuclear position and reaching out to 5<sup>โฒโฒ</sup>, where the signal fades into the background. We also obtained a radial profile for the Chandra PSF, by using the CIAO task chart, which simulates, by ray-tracing, a source with the same spectrum as the target and an exposure time appropriate to derive the statistics needed to define the PSF shape with virtually no errors (we choose 100 ks). After normalizing for the different exposure times, we compared the source and the Chandra simulated PSF, as described in the CIAO Threads, using the SHERPA software. In all cases but one, the source is consistent with a point source. In Figure 1 we show the radial profile of IRAS 01072+4954; this source is characterized by the presence of faint extended emission between 1$`{}_{}{}^{\prime \prime }3^{\prime \prime }`$ from the center, which correspond to a region of 0.5-1.3 kpc (at the redshift of the source). This effect is more pronounced in the soft (0.3-2 keV) than in the hard (2-10 keV) energy band.
## 5. Spectral Analysis
Both Chandra and BeppoSAX spectral data have been grouped to at least 20 counts per bin and spectra in the range of 0.3-10 keV have been analyzed using the XSPEC software package. We first fitted Chandra and BeppoSAX spectra separately. Each spectrum has been initially fitted with a single power-law plus Galactic absorption. In all objects this simple parametrization describes well the observed spectra and the spectral shapes are within the typical ranges observed for AGNs, i.e. photon indeces from 1.7 and 2, except for IRAS 04392$``$0123 which shows a flatter spectrum ($`\mathrm{\Gamma }1.3`$) both in Chandra and BeppoSAX observations. Intrinsic absorption values, when measured, are of the order of 10<sup>21</sup> cm<sup>-2</sup> or slightly in excess of the Galactic value (Dickey & Lockman, 1990). In the case of the flat spectrum of IRAS 04392$``$0123, the upper limit derived on the column density (N<sub>H</sub> $`<`$ 6$`\times `$10<sup>20</sup> cm<sup>-2</sup>) remains consistent with the Galactic absorption even when fixing the spectral slope to the AGN canonical value ($`\mathrm{\Gamma }=1.9`$).
BeppoSAX spectral results are listed in Table LABEL:bep. Although at lower significance, the BeppoSAX spectral parameters are in agreement within the errors with those derived from the Chandra data. We have therefore fitted all the available data for the same source together. The final spectra are plotted in Figure 2 while the spectral results are presented in Table LABEL:spec. We also report the results from Georgantopoulos, Zezas & Ward (2003), Georgantopoulos et al. (2004) in the same table for ease of reference. The quoted errors on the spectral parameters correspond to the 90% confidence level for one interesting parameter.
We discuss here each source in turn. IRAS 20051$``$1117 is the brightest source of the sample observed by BeppoSAX. Its spectrum is well described by a simple power-law, with a photon index ($`\mathrm{\Gamma }`$$``$1.9) in agreement with the one reported by Georgantopoulos et al. (2004), using Chandra and XMM-Newton data. The addition of a Gaussian component consistent with an FeK$`\alpha `$ line (at 6.29$`{}_{0.48}{}^{}{}_{}{}^{+0.28}`$ keV, rest frame) is significant at only 2$`\sigma `$ (EW = 282$`{}_{214}{}^{}{}_{}{}^{+215}`$ eV), again in agreement with Georgantopoulos et al. (2004) Chandra results. Residuals are also present around 6.9 keV (see Figure 2) but at a lower statistical significance. At high energies (13-200 keV), the PDS data fall above the simple power-law extrapolation from the MECS (1.3-10 keV), by a factor of $`5`$. We explored the $`1^{}`$ region around the target to look for possible contaminant sources. We found 1RXS J200433.5-111345, detected by RASS at $`50^{}`$ from IRAS 20051$``$1117 and brighter by a factor 1.3, corresponding to a galaxy detected by IRAS (no other relevant information like redshift, morphological classification, etc. is available). The source could be the contaminant if it hosts an absorbed nucleus. Given the discrepancy in the relative fluxes, the partial contamination from 1RXS J200433.5-111345 or other sources is likely, and therefore we make no further use of the PDS spectrum. On the other hand a Compton thick AGN in IRAS 20051$``$1117 is ruled out by the lack of a FeK line with high EW (see below).
IRAS 01072+4954 and IRAS 04392$``$0123 have been observed by both Chandra and BeppoSAX. In Figure 2 we show Chandra and BeppoSAX spectra fitted simultaneously leaving the relative normalization free to vary. Their values are in the range 1-1.3, indicating that the Chandra and BeppoSAX normalizations are within 30%, that would either indicate a minimal flux variation between the two observations, or the uncertainty in the relative calibration between instruments.
In the Chandra spectrum of IRAS 01072+4954 a thermal component (MEKAL) has been added to the model to account for residuals visible below 2 keV above the simple power law. Such component has a 0.3-2 keV flux of 4.6$`\times `$10<sup>-14</sup> erg cm<sup>-2</sup> s<sup>-1</sup> (L<sub>0.3-2keV</sub> $``$ 5.8$`\times `$10<sup>40</sup> erg s<sup>-1</sup>) and it is significant at 99.99% (via an F-test), consistently with the presence of the faint extended soft component revealed by its radial profile.
IRAS 20069+5929 has the highest fitted value for intrinsic N<sub>H</sub>, which is however still consistent with absorption from the host galaxy.
The presence of an Fe line at 6.4 keV is not statistically significant in any of the sources observed by Chandra. In the case of IRAS 00317$``$2142 and IRAS 01072+4954 the upper limit to the equivalent widths at 90% (as shown in Table LABEL:spec) is not stringent, due to the poor statistics above 6 keV. In the other cases, the upper limits are always consistent with expectation from the low intrinsic absorption measured, ruling out the Compton thick hypothesis (Bassani et al. 1999).
## 6. Timing Analysis
### 6.1. Short-term variability
The ACIS-I light curves have been examined in order to look for short term variability. In Figure 3 we show the 0.3-10 keV Chandra light curves of all sources. They have been extracted from circular regions of 3<sup>โฒโฒ</sup> radius, binned at 1 ks and fitted with a constant. The resulting constant values and $`\chi `$$`{}_{}{}^{2}/dof`$ are given in Table 6 for the 0.3-10 keV, 0.3-2 keV and 2-10 keV energy ranges. In the cases of IRAS 01072+4954, IRAS 01319$``$1604 and IRAS 20069+5929 the hypothesis of a constant flux is rejected at the 99% confidence level. In particular, IRAS 01072+4954 and IRAS 01319$``$1604 show larger flux variations in the soft X-ray band. To illustrate this point, in Figure 4 we plot the soft and hard light curves for IRAS 01319$``$1604, the most variable source of the sample. The non-detection of variability in IRAS 04392$``$0123 is possibly due to the low statistics available. We examine also the BeppoSAX light curve of IRAS 20051$``$1117 (see Fig. 5), in bin sizes of 3 ks. Fitting a constant to the observed count rate we reject the constant flux hypothesis at more than 95% level (see Table 6). However, Chandra or XMM-Newton short-term light curves show no variability either in this source (Georgantopoulos et al. 2004), nor is detected in Chandra for IRAS 00317$``$2142 (Georgantopoulos, Zezas & Ward 2003). The variability of a factor 2-3 for IRAS20051-1117 over a time scale of $``$4 ks implies (using light crossing arguments) that the dimensions of the emitting region are typical of a nuclear source ($`1.2\times 10^{14}`$ cm, i.e. $`4\times 10^5`$ pc).
When the statistics allows it, we have attempted to extract the spectrum at different intervals to check for possible spectral variations, suggested by the higher variability measured in the soft band. However, the spectral parameters measured in the different states (high and low flux) are always consistent within errors.
### 6.2. Long-term variability
In order to examine the presence of long-term variability we compared the observed fluxes from all available X-ray measurements, both from literature and from our BeppoSAX and Chandra data presented here. Figure 6 shows the long term light curves for the six sources in the sample. Errorbars have been plotted for all the fluxes except those derived from the literature for which errors on the count rates were not available. For a better comparison, all ROSAT, BeppoSAX and Chandra 0.3-2 keV fluxes have been recomputed, assuming the Chandra best fit spectral shape as reported in Table LABEL:spec. The resulting 0.3-2 keV unabsorbed fluxes are given in Table 7. All our sources have been observed by the ROSAT All-Sky Survey (RASS) with the Position Sensitive Proportional Counter (PSPC) in the 0.1-2.4 energy band. The RASS count rates and errors have been taken from Boller et al. (1992). IRAS 00317$``$2142 and IRAS 01319$``$1604 have also been observed in a PSPC pointed observation. In the 2 yr period between the RASS and the PSPC pointed observation, IRAS 00317$``$2142 does not show significant variability while IRAS 01319$``$1604 has nearly doubled its flux (measured by the same instrument). ASCA and Chandra count rates of IRAS 00317$``$2142 have been taken from Georgantopoulos (2000) and Georgantopoulos, Zezas & Ward (2003), respectively. In the case of IRAS 20051$``$1117, Chandra and XMM-Newton count rates, taken from Georgantopoulos et al. (2004), were both obtained on 2002 April 1 and the measured soft fluxes are comparable within a few percent; therefore we only consider the Chandra data for the long-term analysis.
Unfortunately, the statistics in the ROSAT observations is not sufficient to derive a reliable spectral measurement, therefore a possible spectral variation could have occurred between the ROSAT and the BeppoSAX/Chandra epochs, even if we do not observe spectral variability between the BeppoSAX and Chandra observations. While we do not observe flux variations on time scales of $``$ 2 yr between the BeppoSAX and the Chandra observations, IRAS 04392$``$0123 has experienced a variation by a factor of $``$ 24 in $``$ 10 yr between the ROSAT and the Chandra observations, similarly to the case of IRAS 00317$``$2142 which varied by a factor of $``$ 20 (see also Georgantopoulos, Zezas & Ward, 2003). Also IRAS 01072+4954 and IRAS 01319$``$1604 varied, even though only by a factor of $``$ 3. The strong long term variability in the X-rays seems very common in this class.
## 7. Discussion
The X-ray analysis of Composite galaxies has revealed their AGN dominance in this spectral domain.
In what follows we make use of various diagnostic diagrams and multi-frequency data to investigate the reasons for the contradictory optical/X-ray classification of these objects, and suggest a global explanation of the phenomenon.
### 7.1. Diagnostic Diagrams
The optical diagnostic diagrams obtained using standard emission line ratios (Veilleux & Osterbrock, 1987) show that Composites are very close to the boundary region between starbursts and AGNs confirming the M96 classification. This is shown in Figure 7 where the \[O III\]/H$`\beta `$ ratio is plotted vs. the \[NII\]/H$`\alpha `$ (the other two standard diagnostics which make use of \[SII\]/H$`\alpha `$ and \[OI\]/H$`\alpha `$ ratios give similar results). The lines represent the theoretical starburst limits, a standard one which have been taken from Kewley et al. (2001) (together with the dotted lines which indicate the error range) and an updated estimate for the starburst boundary derived from the SDSS observations (from Kauffmann et al. 2003). The location of the โCompositeโ region should be between these two lines (Hornschemeier et al. 2005).
However, when the flux of the optical emission line \[O III\]$`\lambda `$5007 is combined with the infrared and X-ray fluxes, Composite galaxies are classified as AGNs rather than starburst. In Figure 8 the diagram with the combination of F<sub>X</sub>/F$`_{[OIII]}`$ vs. F$`_{[OIII]}`$/F<sub>IR</sub> ratios is shown. These flux ratios have been used to separate the AGN and the starbursts contributions by Panessa & Bassani (2002) and Braito et al. (2004), based on the fact that the \[O III\]$`\lambda `$5007 flux is associated with the AGN and the far-infrared emission is associated mainly with the star-forming activity. At the same time they are a powerful tool in the detection of heavy obscuration not seen in X-rays below 10 keV (as in Compton thick objects, Bassani et al. 1999). Clearly Composite objects all fall in the AGN region, in good agreement with the relative broadness of the \[O III\]$`\lambda `$5007 line found by M96 which points to an AGN origin. The diagram further shows that they should all be classified as Compton thin AGN (note that the X-ray fluxes used in the plot are those observed in the โlow-stateโ epoch, and therefore the classification would be also valid in the โhigh-stateโ). This indicates very little amount of absorption, in agreement with the results obtained from the X-ray analysis, i.e. the absence of a strong Fe-line and X-ray obscuration. Therefore the relative weakness of the AGN optical emission lines cannot be explained as due to the presence of absorbing material in the line of sight, as suggested by M96.
### 7.2. X-ray vs. infrared luminosity diagram
Having assessed the importance of the AGN component from the X-ray, infrared and \[O III\] emissions combined, we want to estimate here the amount of X-ray emission expected to be produced by the presence of a starburst. Therefore we derived the Star Formation Rate (SFR) from the far-infrared luminosity (Kennicutt 1998), assuming that the latter (given in Table 1) is all due to the starburst. The SFRs obtained are in the range of 5-34 M/yr (note that these values are upper limits since we have not subtracted the possible AGN contribution from the FIR luminosity). Subsequently, the L<sub>X</sub> in the 2-10 keV band associated to the derived SFRs have been estimated as in Nandra et al. (2002), Grimm et al. (2003), Ranalli et al. (2003), Persic et al. (2004). The derived correlations have been plotted in Figure 9, together with the observed hard L<sub>X</sub> vs. L<sub>FIR</sub> for each Composite object. For all sources, regardless of the correlation considered, the X-ray observed emission is well above the value expected to be produced by a starburst. The SFR derived here for the Composites correspond to a starburst that is not bright enough to produce the observed X-ray emission, which is then mostly given by the AGN, as suggested by our X-ray analysis. Observed X-ray emission in excess to that predicted from SFR has also been observed in high-redshift submillimeter sources (Alexander et al. 2003), that might point to the presence of an active nucleus also for this class of objects. Only in the case of IRAS 00317$``$2142 the observed X-ray luminosity is close to the expected value for a starburst; this is in agreement with the results of the multi-wavelength analysis for this object presented in the next Section.
### 7.3. Spectral Energy Distributions
Spectral Energy Distributions (SEDs) have been assembled from radio ($`\nu 10^8`$ Hz) to hard X-rays ($`\nu 10^{18}`$ Hz) for our sample sources. Radio, far to near infrared, optical, X-ray data have been taken from NED, and complemented with the Chandra data from this work. All X-ray data points have been plotted with the spectral shape measured in the Chandra observations.
In Figure 10 we compare the observed SEDs with the templates of Medium Energy Distribution for radio quiet quasars (Elvis et al. 1994), starburst galaxies (Schmitt et al. 1997) and normal spiral galaxies (Elvis et al. 1994). The templates we show are normalized to match, and not exceed data points. In particular, the starburst template is normalized to the radio-IR portion of the spectrum, while the X-ray one is normalized to the different X-ray states observed. We stress that we are only interested in deriving an overall consistent picture, without attempting a quantitative computation of the contributions of the different components, that would require at least simultaneous data to account for the observed variability in the X-ray band. Moreover, it must be taken into account that all data points have been taken using different apertures, in particular for the optical band where the emission is heavily contaminated by the host galaxy. The SEDs show clearly the Composite nature of these objects: the AGN dominates at X-ray wavelengths, while the starburst is the most important contributor to the mid$`/`$far IR emission; the host galaxy template accounts for the optical appearance.
It is evident that, in the optical band, the AGN contribution is always less than that of the starburst, and, except for a couple of sources during their โhighโ state, even by a factor of 10. Note instead that in IRAS 00317$``$2142, when it is in the low state, the contributions to the X-ray emission from the AGN and the starburst become comparable (in agreement with Figure 9).
The AGN contributes to the total bolometric luminosity from less than 10% in the โlow flux stateโ, to 15-30% in the โhigh flux stateโ thus making the SB contribution dominant in the bolometric output, except possibly during bright AGN flares.
Mid-infrared and L-band spectroscopy could provide a powerful way to disentangle the starburst from the AGN component (Genzel et al. 1998, Imanishi 2002, Risaliti et al. 2003, Lutz et al. 2003). Moreover broad band data from the near to the far-infrared frequencies could provide a more detailed characterization of the SEDs, in particular by exploiting the Spitzer Space Telescope unprecedented capabilities (see e.g. Franceschini et al. 2005).
### 7.4. A weak and low M<sub>BH</sub> AGN?
IRAS 00317$``$2142, IRAS 20051$``$1117 and IRAS 20069$`+`$5929 are the only three sources for which a weak broad component of the H$`\alpha `$ emission line has been detected in their optical spectra, while in all objects the narrow H$`\alpha `$ component is probably due to the starburst. A correlation between the 2-10 keV X-ray luminosity versus the H$`\alpha `$ luminosity has been widely observed in high and low luminosity AGNs (Ward et al. 1988, Ho et al. 2001). IRAS 00317$``$2142, IRAS 20051$``$1117 and IRAS 20069$`+`$5929, for which we consider only the broad H$`\alpha `$ components, follow the Ho et al. (2001) correlation (which applies to both high and low luminosity AGNs). The remaining three Composites for which the broad H$`\alpha `$ was not measured are generally fainter both in X-rays and in the optical band; therefore it is likely that an optical emission line luminosity of a factor of 2-3 lower could be difficult to measure. Using the Kennicutt (1998) relation, we derived the SFR from the H$`\alpha `$ luminosity produced by the narrow component and we estimated the expected 2-10 keV luminosity to be produced by that SFR (Ranalli et al. 2003). The expected 2-10 keV luminosities are in the range of 10<sup>39-40</sup> erg s<sup>-1</sup> , i.e. a factor of 50-100 lower than the observed ones. This result is consistent with our previous interpretation that in all objects the narrow H$`\alpha `$ component is probably due to the starburst, and the X-rays to the AGN.
For the three objects for which we have a measure of the broad H$`\alpha `$ component, we attempt an estimate of the Black Hole Mass. We make use of the formula $`M_{BH}=v^2R_{BLR}/G`$ (McLure & Dunlop 2001) and of a few additional assumptions. The quantities $`v=1.5FWHMH_\beta `$ and $`R_{BLR}=32.9(\lambda L_{5100}/10^{44}`$ erg s<sup>-1</sup>)<sup>0.7</sup> in light days (Kaspi et al. 2000) require a measure of the broad $`H_\beta `$ emission, which is not detected because swamped by the narrow line produced by the starburst and of the AGN continuum, which we do not measure directly. We substitute the FWHM of H$`\beta `$ with the FWHM of the broad H$`\alpha `$ that we attribute to the AGN (from M96). If anything, the H$`\alpha `$ is usually broader than H$`\beta `$, therefore we will overestimate the resulting mass. We infer the AGN continuum from the template fitted to the X-ray luminosity. This is a rough approximation of the real illuminating continuum, but it should be correct within an order of magnitude at least. Plugging these numbers in the above formula we derive numbers for the mass of the order of a few 10<sup>5-6</sup> M$``$. An order of magnitude uncertainty on the L<sub>5100ร
</sub> changes the mass by 10<sup>0.7</sup>. We therefore can conclude that most probably the masses of these black holes are small and that they could possibly undergo strong changes in accretion rate when they brighten. We have found a possible analogy between our sources and the Seyfert 1.5 NGC 4395. It is a low luminosity AGN with a small black hole mass, that shows X-ray variability, although more violent than what observed in the Composites (Vaughan et al. 2005).
The Black Hole Mass range is consistent with all these objects being late type galaxies (as indicated by their morphological types in Table 1) with small stellar bulges. In fact, from their total B-band magnitude and an average bulge/total flux ratio typical of their morphological type (de Jong 1996), we derived the black hole mass values which are consistent with those previously found, when assuming that the $`L_{B,bulge}M_{BH}`$ correlation (e.g., Yu & Tremaine 2002) still holds for late type spirals.
Detailed follow up observations (e.g. high spatial resolution optical spectroscopy) are needed in order to provide a more reliable estimate of the black hole mass for these objects.
## 8. Conclusions
New Chandra and BeppoSAX observations of 4 and 3 Composites, respectively, have deepened our knowledge of this class of sources. Based on the X-ray analysis presented here, Composite galaxies behave like typical type 1 AGNs in the X-ray band: their emission is dominated by a bright nuclear source, whose spectral properties are typical of this class. A power-law model with spectral index ($`\mathrm{\Gamma }`$ = 1.7-2.1) and little intrinsic absorption (N<sub>H</sub> $`<`$ 4$`\times `$ 10<sup>21</sup> cm<sup>-2</sup>) well describe their spectra. Iron lines are not significantly detected (upper limits on their equivalent widths are below $``$ 400 eV). Large flux variability is found on many different temporal scales.
A broad analysis of the whole Composite sample, carried out by adding to our X-ray data the results from Georgantopoulos (2000), Georgantopoulos, Zezas & Ward (2003), Georgantopoulos et al. (2004) and multi-wavelength data from the literature, has revealed that the study of this class is relevant both for the investigation of the AGN-starburst connection and for the X-ray properties in medium/low luminosity AGNs.
Interestingly enough, AGN and starburst activity seem to be present with almost the same intensity in this class of objects. Spectral Energy Distributions have clearly shown that the different emission components contribute to different spectral energy bands: the infrared emission is probably dominated by the starburst and the X-ray one derives from the AGN. The optical continuum is mainly contaminated by the host galaxy light while the emission line spectrum shows narrow emission lines produced by starburst. The H$`\alpha `$ emission line is probably associated to both: the narrow component to the starburst and the broad component, when bright enough to be detected, correlates well with the hard X-ray luminosity, and therefore is to be ascribed to photoionization by the AGN.
It is unlikely that heavy obscuration that could explain the weakness of the AGN in the optical band is present, as shown by the pronounced X-ray variability and low column densities and by the flux diagnostic diagrams shown in section 7.1. A dusty clumpy ionized absorber, able to selectively obscure the optical emission, that leaves the X-ray emission almost unabsorbed (Georgantopoulos 2000; Maiolino et al. 2001) is clearly possible. However the detection of broad H$`\alpha `$ with FWHM typical of type 1 objects makes this hypothesis remote. In conclusion, the lack of clear indications of the presence of an AGN in the optical spectra of the Composites is probably due to a combination of the faintness of the AGN itself and of the masking effect of the starburst.
As mentioned above, the most striking characteristic of Composite galaxies is the strong, long- and short-term, X-ray variability observed. ROSAT observation taken between 1991 and 1992 reveal bright soft nuclei at the level of 10<sup>42-43</sup> erg/s. Chandra and BeppoSAX observations taken nearly 10 years later reveal that, in most cases, these objects are fainter. Two objects show a variation in flux by a factor of 20, while in two other objects the fluxes decreased by a factor of 2-3. It is likely that the inclusion of these sources in the RASS-IRAS correlation of M96 is partly due to their bright state at that time. What is the cause for their brightening? Are they now fading away or just changing state every once in a while?
Concerning the flux variation of a factor of $``$20 over $``$10 yrs, we cannot say much: it is not an extreme value, but certainly not the most common (see eg. Ulrich et al 1997). Furthermore, the data sampling is very scanty and does not allow us to determine the time dependence of the variability. On the contrary, the flux variability on short time scales, with variations by a factor of 2-3 on time scales of $``$ 4 ks, coupled to an estimate of low values for the central black hole mass (even though based more on assumptions that on measurements), could fit with the observed trend of higher variability for smaller black hole masses, and could also be linked to variations in the accretion rate (OโNeill et al. 2005, Mushotzky et al. 1993).
Even if small, this sample has a statistically sound definition that makes results applicable to other candidates. We cannot estimate their space density, however we notice that several examples of objects with similar characteristics have already appeared in the literature (Griffiths et al. 1996, Della Ceca et al. 2001, Guainazzi et al. 2000), which suggests that they might constitute a non-negligible fraction of the AGN family.
We therefore suggest the use of a combination of diagnostic ratios, such as those based on the \[OIII\]/FIR ratio and the presence of X-ray emission to pinpoint members of this class of Composites. New Spitzer Space Telescope and Chandra data would be of crucial importance to deepen our knowledge on the Composite sample and to enlarge the number of candidates belonging to the same class.
FP warmly thanks for hospitality the Center for Astrophysics where much of this work has been realized. We thank the anonymous referee for providing useful comments. We acknowledge the contribution of Massimo Cappi and Mauro Dadina for an earlier involvement in the paper, Raffaella Landi for her help with the PDS data, Jonathan McDowell for the SED-making program TIGER, Martin Elvis, and Andreas Zezas for useful comments. This work has been supported by the NASA grants GO3-4131X and NAS8-39073 to the Chandra X-Ray Center. This work has received partial financial support by ASI and Cofin Miur. RDC acknowledge partial financial support from the MIUR (Cofin-03-02-23). This research makes use of the NASA/IPAC Extra-galactic Database (NED) and of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. |
warning/0506/math0506547.html | ar5iv | text | # Coarse dimensions and partitions of unity
## 1. Introduction
There are three concepts of dimension associated with variants of the coarse category of proper metric spaces. The original one, the asymptotic dimension of Gromov , and dimensions $`asdim^{}(X)`$ and $`\text{dim}^c(X)`$ introduced by Dranishnikov . All three dimensions are defined in seemingly different ways:
1. The asymptotic dimension of Gromov (see or Definitions 1-2 in on p.1103) is the smallest integer $`n`$ such that for every $`M>0`$ there is a uniformly bounded family $`๐ฐ`$ of Lebesque number at least $`M`$ and multiplicity (or order) at most $`n+1`$.
2. The asymptotic dimension $`asdim^{}(X)`$ of Dranishnikov (see Definition 3 in on p.1104) is the smallest integer $`n`$ such that for every proper function $`f:XR_+`$ there is a contracting map $`\varphi :XK`$ to an $`n`$-dimensional asymptotic polyhedron such that for each $`M>0`$ there is a compact subset $`C`$ of $`X`$ with the property that $`\varphi ^1(B(\varphi (x),M))B(x,f(x))`$ for all $`xXC`$.
3. The coarse dimension $`\text{dim}^c(X)`$ of Dranishnikov (see Definition 4 in on p.1105) is the smallest integer $`n`$ such that $`R^{n+1}`$ is an absolute extensor of $`X`$ in the category of proper asymptotically Lipschitz functions. That dimension coincides with the dimension of the Higson corona $`\nu (X)`$ of $`X`$ (see Theorem 6.6 in on p.1111).
One of the main motivations behind the research in asymptotic dimension is the result of Yu (see and ) that the Novikov Conjecture holds for groups of finite asymptotic dimension.
In this paper we work in the coarse category of all metric spaces and we devise a unified way of defining five dimensions: coarse dimension $`\text{dim}_{rse}^{coa}(X)`$, major coarse dimension $`\text{dim}_{RSE}^{COA}(X)`$, asymptotic dimension $`asdim(X)`$, minor asymptotic dimension $`ad(X)`$, and large scale dimension $`\text{dim}_{scale}^{large}(X)`$.
In case of proper metric spaces, three of them coincide with the above dimensions. Namely, $`\text{dim}_{RSE}^{COA}(X)=dim^{}(X)`$, $`\text{dim}_{rse}^{coa}(X)=\text{dim}^c(X)`$, and $`asdim(X)`$ coincides with Gromovโs asymptotic dimension. The fourth one, the minor asymptotic dimension, is a variant of Gromovโs dimension. The large scale dimension is always equal to the coarse dimension and the reason we are introducing it is to simplify proofs of the relations between the three basic dimensions which we do in a much simpler way than as described in Dranishnikovโs paper . The main relations between dimensions are as follows:
1. There are two strands of inequalities: $`asdim(X)\text{dim}_{RSE}^{COA}(X)\text{dim}_{rse}^{coa}(X)`$ and $`asdim(X)ad(X)\text{dim}_{rse}^{coa}(X)`$,
2. In each strand (for unbounded spaces $`X`$), finiteness of a larger dimension implies its equality with all smaller dimensions in the strand.
We do not know of any unbounded space $`X`$ such that a larger dimension in a strand is infinite and a smaller dimension is finite.
Our fundamental concept is that of a coarse family and we follow the well-established route of defining the covering dimension by refining covers with covers of a prescribed multiplicity. In classical dimension theory one deals with two cases: finite covers and infinite covers. There, for paracompact spaces, the two concepts coincide. In the case of coarse covers we get two concepts of coarse dimension whose equality remains unresolved.
A finite family $`๐ฐ`$ of subsets of $`X`$ is coarse if and only if there is a slowly oscillating partition of unity $`f`$ on $`XB`$ for some bounded subset $`B`$ of $`X`$ whose carriers $`Carr(f)`$ refine $`๐ฐ`$. That explains why, in the case of a proper metric space $`X`$, its coarse dimension equals the covering dimension of the Higson corona of $`X`$.
Our basic strategy is to associate natural functions with objects and declare those objects to be coarse, asymptotic, or large scale if the function is coarsely proper. A function $`f`$ is coarsely proper if $`f(E_n)\mathrm{}`$ whenever $`E_n\mathrm{}`$. Elements $`E_n`$ related to objects could be points in a metric space, bounded subsets in a metric space, or covers of a metric space (in which case divergence to infinity is measured by the size of the Lebesque number). In (p.1089) coarsely proper functions were defined as those $`f:XY`$ such that $`f^1(A)`$ is bounded whenever $`A`$ is bounded in $`Y`$. Notice that our definition generalizes the one from .
The authors are grateful to Jose Higes for helpful comments.
## 2. Preliminaries
Given a subset $`A\mathrm{}`$ of a metric space $`X`$ the most basic function is the distance function $`d_A:XR_+`$: $`d_A(x)=dist(x,A)`$.
###### Definition 2.1.
Given a subset $`A`$ of a metric space $`(X,d_X)`$ the ball $`B(A,M)`$ is defined to be the set $`\{xXdist(x,A)<M\}`$ if $`M>0`$, it is defined to be the set $`\{xXdist(x,XA)>M\}`$ if $`M<0`$, and it is simply $`A`$ if $`M=0`$.
The distance function leads to the first concept of divergence to infinity: $`x_n\mathrm{}`$ if $`d_X(x_n,x_0)\mathrm{}`$ for some (and hence for all) $`x_0X`$. However, $`dist(x,A)`$ is a function of two arguments and we can use the second one to define divergence to infinity for bounded subsets of $`X`$. Here is a more general concept.
###### Definition 2.2.
A family $`๐ฐ`$ of bounded subsets of $`X`$ is called coarsely proper if the function $`Ud_U(x_0)`$ is coarsely proper for some (and hence for all) $`x_0X`$. Here $`๐ฐ`$ is considered as a subspace of all bounded subsets of $`X`$ with the Hausdorff metric.
Notice that a sequence $`\{A_n\}`$ of bounded subsets of $`X`$ containing points $`x_nA_n`$ so that $`x_n\mathrm{}`$ is coarsely proper if and only if every bounded subset of $`X`$ intersects at most finitely many elements of the sequence. In that case we write $`A_n\mathrm{}`$ and that form of divergence to infinity is of most interest to us.
###### Lemma 2.3.
If $`๐ฐ`$ is a coarsely proper cover of $`X`$, then every selection function $`\varphi :X๐ฐ`$ (that means $`x\varphi (x)`$) is coarsely proper.
###### Proof.
Suppose $`x_n\mathrm{}`$ and $`x_nU_n๐ฐ`$. Clearly, $`U_n\mathrm{}`$ in the Hausdorff metric. Pick $`x_0X`$. Since $`d_{U_n}(x_0)\mathrm{}`$, every bounded subset of $`X`$ intersects at most finitely many elements of the sequence $`\{U_n\}`$ and any selection function $`\varphi `$ is coarsely proper. โ
###### Definition 2.4.
Given a family $`๐ฐ`$ in $`X`$, the local Lebesque number $`L_๐ฐ(x)R_+\mathrm{}`$ is defined as the supremum of $`dist(x,XU)`$, $`U๐ฐ`$. If $`U=X`$ for some $`U๐ฐ`$ it is defined to be infinity.
Notice that either $`L_๐ฐ\mathrm{}`$ at all points or it is a natural Lipschitz function associated with $`๐ฐ`$. More precisely $`|L_๐ฐ(x)L_๐ฐ(y)|d_X(x,y)`$.
###### Definition 2.5.
The Lebesque number $`L(๐ฐ,A)`$ is $`inf\{L_๐ฐ(x)xA\}`$.
###### Definition 2.6.
A family of subsets $`๐ฐ`$ of a metric space $`X`$ is called coarse if $`L_๐ฐ`$ is coarsely proper (as a function from $`X`$ to $`R\mathrm{}`$).
An alternative way to define coarse families is to require $`L(๐ฐ,A)\mathrm{}`$ as $`A\mathrm{}`$. Yet another way is to state that $`L(๐ฐ,XB(x_0,t))\mathrm{}`$ as $`t\mathrm{}`$.
###### Proposition 2.7.
* A family $`๐ฐ=\{A\}`$ consisting of one subset $`A`$ of $`X`$ is coarse if and only if $`XA`$ is bounded.
* A family $`๐ฐ=\{X_1,X_2\}`$ consisting of two subsets of $`X`$ is coarse if and only if $`d_X`$ restricted to $`(XX_1)\times (XX_2)`$ is coarsely proper.
* A family $`๐ฐ=\{X_1,X_2,\mathrm{},X_n\}`$ consisting of finitely many subsets of $`X`$ is coarse if and only if the function $`d_๐ฐ(x):=\underset{i=1}{\overset{n}{}}dist(x,XX_i)`$ is coarsely proper.
###### Proof.
1. If $`XA`$ is bounded, then $`L_๐ฐ(x)dist(x,XA)`$ and $`L_๐ฐ`$ is coarsely proper. If $`XA`$ is unbounded, then $`L_๐ฐ(x)=0`$ at all $`xXA`$ and $`L_๐ฐ`$ is not coarsely proper.
2. Suppose $`๐ฐ=\{X_1,X_2\}`$ is coarse and $`x_n\mathrm{}`$, $`y_n\mathrm{}`$, for some $`x_nXX_1`$, $`y_nXX_2`$. Notice $`L_๐ฐ(x_n)d_X(x_n,y_n)`$, so $`d_X(x_n,y_n)\mathrm{}`$.
If $`๐ฐ=\{X_1,X_2\}`$ is not coarse, then there is a sequence $`z_n\mathrm{}`$ with $`L_๐ฐ(z_n)`$ bounded by $`M`$. We can produce $`x_nXX_1`$ and $`y_nXX_2`$ so that $`d_X(z_n,x_n)<M+1`$ and $`d_X(z_n,y_n)<M+1`$ for all $`n`$. Now, $`d_X(x_n,y_n)<2M+2`$, a contradiction.
3. Notice $`d_๐ฐ(x)L_๐ฐ(x)`$ and $`mL_๐ฐ(x)d_๐ฐ(x)`$. โ
###### Definition 2.8.
Given a function $`f:XY`$ of metric spaces, its Lebesque number transfer $`L^f:R_+R_+\mathrm{}`$ is the supremum of all functions $`\alpha :R_+R_+\mathrm{}`$ such that $`L(๐ฐ,Y)t`$ implies $`L(f^1(๐ฐ),X)\alpha (t)`$ for all families $`๐ฐ`$ of subsets of $`Y`$.
###### Definition 2.9.
A function $`f:XY`$ of metric spaces is coarse if the Lebesque number transfer $`L^f`$ is coarsely proper.
An alternative definition of coarse functions is to require the function $`๐ฐL(f^1(๐ฐ),X)`$ to be coarsely proper on the set of covers of $`Y`$.
Let us show that our definition of coarse functions coincides with that of Roe .
###### Proposition 2.10.
A function $`f:XY`$ is coarse if and only if for every $`R>0`$ there is $`M>0`$ such that $`d_X(x,y)R`$ implies $`d_Y(f(x),f(y))M`$ for all $`x,yX`$.
###### Proof.
Notice that if $`M>0`$ and $`N>0`$ are numbers such that $`d_X(x,y)<M`$ implies $`d_Y(f(x),f(y))<N`$, then $`L^f(N)M`$. Therefore $`f`$ being coarse in the sense of Roe implies $`L^f`$ being coarsely proper.
Conversely, if $`L^f(N)M`$, then consider the cover $`๐ฐ=\{B(z,N)\}_{zY}`$ whose Lebesque number is clearly at least $`N`$. If $`d_X(x,y)<M`$, then there is $`z`$ so that $`x,yf^1(B(z,N))`$. Hence $`d_Y(f(x),f(y))<2N`$ and $`f`$ is coarse. โ
Dranishnikov (p.1088) defined asymptotically Lipschitz functions $`f:XY`$ as those for which there are constants $`M>0`$ and $`A`$ such that $`d_Y(f(x),f(y))Md_X(x,y)+A`$ for all $`x,yX`$. Let us relate this concept to the Lebesque number transfer.
###### Proposition 2.11.
A function $`f:XY`$ is asymptotically Lipschitz if and only if there is a linear function $`tmt+b`$ so that $`m>0`$ and $`L^f(t)mt+b`$ for all $`t`$.
###### Proof.
Suppose there are constants $`M>0`$ and $`A`$ such that $`d_Y(f(x),f(y))Md_X(x,y)+A`$ for all $`x,yX`$. Given a cover $`๐ฐ`$ of $`Y`$ with $`L(๐ฐ,Y)t`$ and given $`xX`$, the ball $`B(x,(tA\delta )/M)`$ is mapped by $`f`$ into the ball $`B(f(x),t\delta )`$ which is contained in an element of $`๐ฐ`$ for all $`\delta >0`$. That shows the Lebesque number of $`f^1(๐ฐ)`$ to be at least $`(tA)/M)`$. Conversely, if $`L^f(t)mt+b`$ for all $`t`$ and $`m>0`$, then we claim $`d_Y(f(x),f(y))<2d_X(x,y)/m+2(1b)/m`$. Indeed, put $`d_X(x,y)=s`$ and consider the cover $`๐ฐ=\{B(z,(s+1b)/m)\}_{zY}`$ whose Lebesque number is clearly at least $`(s+1b)/m`$. There is $`z`$ so that $`x,yf^1(B(z,(s+1b)/m))`$. Hence $`d_Y(f(x),f(y))<2(s+1b)/m`$ and $`f`$ is asymptotically Lipschitz. โ
###### Proposition 2.12.
Given a function $`f:XY`$ of metric spaces the following conditions are equivalent:
* $`f`$ sends bounded subsets of $`X`$ to bounded subsets of $`Y`$ and $`f^1(๐ฐ)`$ is coarse for every coarse family $`๐ฐ`$ in $`Y`$.
* $`f`$ is coarse and coarsely proper.
###### Proof.
1$``$2. Given a bounded subset $`A`$ of $`Y`$ the family $`\{YA\}`$ is coarse (see 2.7). Since $`\{f^1(YA)\}`$ is coarse and $`f^1(YA)=Xf^1(A)`$, $`f^1(A)`$ must be bounded and $`f`$ is coarsely proper.
If $`f`$ is not coarse, we find sequences $`x_n,y_nX`$ so that $`d_Y(f(x_n),f(y_n))>n`$ for each $`n`$ but $`d_X(x_n,y_n)<M`$ for all $`n`$. Since $`f`$ sends bounded subsets of $`X`$ to bounded subsets of $`Y`$, we may assume $`x_n\mathrm{}`$, hence $`y_n\mathrm{}`$. Put $`A=\{x_n\}`$ and $`B=\{y_n\}`$. Using 2.7 we see that $`๐ฐ=\{Yf(A),Yf(B)\}`$ is a coarse family in $`Y`$. Since $`f^1(๐ฐ)`$ is coarse, the family $`๐ฑ=\{XA,XB\}`$, to which $`๐ฐ`$ is a shrinking, is coarse as well. That however contradicts 2.7.
2$``$1. Obviously, coarse functions $`f:XY`$ send bounded subsets of $`X`$ to bounded subsets of $`Y`$. Put $`๐ฑ=f^1(๐ฐ)`$ for some coarse family $`๐ฐ`$ in $`Y`$. To find points $`xX`$ such that $`L_๐ฑ(x)>t`$ we find $`s>0`$ so that $`L^f(s)>t`$ and we find $`u>0`$ such that $`L_๐ฐ(y)>s`$ for $`yYB(y_0,u)`$. Put $`๐ฒ=๐ฐ\{B(y_0,u+s)\}`$. Note $`L(๐ฒ,Y)>s`$. Since $`L(f^1(๐ฒ),X)>t`$, points $`x`$ lying outside of the bounded set $`f^1(B(y_0,u+s))`$ satisfy $`L_๐ฑ(x)>t`$. โ
In the end of this section let us demonstrate the usefulness of the concept of a coarse family by rewording notions from .
In section 5.2 of the concept of asymptotic neighborhood $`W`$ of a subset $`A`$ of $`X`$ is introduced by requiring $`\underset{r\mathrm{}}{lim}dist(AB(x_0,r),XW)=\mathrm{}`$ for some (and hence for all) $`x_0X`$.
###### Proposition 2.13.
$`W`$ is an asymptotic neighborhood of $`A`$ if and only if the pair $`\{XA,W\}`$ is coarse.
###### Proof.
According to part 2 of 2.7 the pair $`\{XA,W\}`$ is coarse if and only if $`d_X`$ restricted to $`A\times (XW)`$ is coarsely proper. That can be easily seen as equivalent to $`\underset{r\mathrm{}}{lim}dist(AB(x_0,r),XW)=\mathrm{}`$ for some (and hence for all) $`x_0X`$. โ
In section 5.2 of (see also ) the concept of asymptotically disjoint subsets $`A`$ and $`B`$ of $`X`$ is introduced by requiring $`\underset{r\mathrm{}}{lim}dist(AB(x_0,r),BB(x_0,r))=\mathrm{}`$ for some (and hence for all) $`x_0X`$.
###### Proposition 2.14.
$`A`$ and $`B`$ are asymptotically disjoint if and only if the pair $`\{XA,XB\}`$ is coarse.
###### Proof.
Apply part 2 of 2.7. โ
Also notice that the concept of an asymptotic separator of (see section 5.2) can be introduced without referring to the Higson corona.
###### Definition 2.15.
A subset $`C`$ of $`X`$ is an asymptotic separator between asymptotically disjoint subsets $`A`$ and $`B`$ if there are asymptotic neighborhoods $`W_A`$ of $`A`$ and $`W_B`$ of $`B`$ such that $`C=X(W_AW_B)`$ and $`W_AW_B=\mathrm{}`$.
## 3. Multiplicity and higher Lebesque numbers
###### Definition 3.1.
Given a family $`๐ฐ`$ of subsets of $`X`$ we define the multiplicity function $`m_๐ฐ:XZ_+\mathrm{}`$ by setting $`m_๐ฐ(x)`$ to be equal to the number of elements of $`๐ฐ`$ containing $`x`$. The global multiplicity $`m(๐ฐ,A)`$ is the supremum of $`m_๐ฐ(x)`$, $`xA`$.
By a coarse refinement $`๐ฑ`$ of a coarse family $`๐ฐ`$ we mean a coarse family such that every element $`V`$ of $`๐ฑ`$ is contained in an element $`U`$ of $`๐ฐ`$. $`๐ฑ`$ is called a shrinking of $`๐ฐ`$ if they are indexed by the same set $`S`$ and $`V_sU_s`$ for all $`sS`$. If $`๐ฑ`$ is a coarse refinement of $`๐ฐ`$ indexed by a different set $`T`$, then one can create a shrinking $`๐ฑ^{}`$ of $`๐ฐ`$ as follows: find a function $`\varphi :TS`$ satisfying $`V_tU_{\varphi (t)}`$ for all $`tT`$. Define $`V_s^{}`$ as $`\{V_ts=\varphi (t)\}`$. Notice that $`๐ฑ^{}`$ has multiplicity at most that of $`๐ฑ`$ and is a coarse shrinking of $`๐ฐ`$.
Given a family $`\varphi =\{\varphi _s:XR_+\}_{sS}`$ of functions its carrier family $`Carr(\varphi )`$ is the family $`\{\varphi _s^1(0,\mathrm{})\}_{sS}`$. The multiplicity $`m(\varphi )`$ of $`\varphi `$ is defined as the multiplicity of its carrier family and its Lebesque number $`L(\varphi )`$ is defined as the Lebesque number of its carrier family.
###### Lemma 3.2.
If $`๐ฐ=\{U_s\}_{sS}`$ is family in $`X`$ such that $`L_๐ฐ(x_0)=\mathrm{}`$ for some $`x_0X`$, then it has a coarse refinement $`๐ฑ`$ of multiplicity at most $`2`$.
###### Proof.
Put $`V_n=\{xX(n1)^2d(x,x_0)<(n+1)^2\}`$ for $`n1`$. โ
###### Lemma 3.3.
If $`๐ฐ=\{U_s\}_{sS}`$ is a family in $`X`$ of multiplicity at most $`n+1`$, then it can be refined by $`๐ฑ=\underset{i=1}{\overset{n+1}{}}๐ฑ^i`$ such that $`L_๐ฑ(x)L_๐ฐ(x)/(2n+2)`$ for each $`xX`$ and each $`๐ฑ^i`$ consists of disjoint sets.
###### Proof.
Define $`f_s(x)=dist(x,XV_s)`$. For each finite set $`T`$ of $`S`$ define $`W_T=\{xX\mathrm{min}\{f_t(x)tT\}>sup\{f_s(x)sST\}\}`$. Notice $`W_T=\mathrm{}`$ if $`T`$ contains at least $`n+2`$ elements. Also, notice that $`W_TW_F=\mathrm{}`$ if both $`T`$ and $`F`$ are different but contain the same number of elements. Let us estimate the Lebesque number of $`๐ฒ=\{W_T\}_{TS}`$. Given $`xX`$ arrange all non-zero values $`f_s(x)`$ from the largest to the smallest. Add $`0`$ at the end and look at gaps between those values. The largest number is at least $`L_๐ฐ(x)`$, there are at most $`n+1`$ gaps, so one of them is at least $`L_๐ฐ(x)/(n+1)`$. That implies the ball $`B(x,L_๐ฐ(x)/(2n+2))`$ is contained in one $`W_T`$ ($`T`$ consists of all $`t`$ to the left of the gap). Define $`๐ฑ_i`$ as $`\{W_T\}`$, all $`T`$ containing exactly $`i`$ elements. โ
###### Lemma 3.4.
If $`๐ฐ=\{U_s\}_{sS}`$ is a coarse family in $`X`$, then it has a coarse refinement $`๐ฑ`$ that is coarsely proper. Moreover, if $`๐ฐ`$ is of finite multiplicity, then we may require $`๐ฑ`$ to be of finite multiplicity as well.
###### Proof.
Let $`๐ฑ=\{V_{s,m}\}_{(s,m)S\times N}`$, where $`V_{s,m}=\{xU_s2^m<d(x,x_0)2^{m+2}\}`$. Notice $`๐ฑ`$ is coarse of multiplicity at most $`2m(๐ฐ)`$. Also, it consists of bounded sets so that for any sequence $`x_k\mathrm{}`$ the conditions $`x_kV_{s(k),m(k)}`$ imply $`V_{s(k),m(k)}\mathrm{}`$. โ
###### Proposition 3.5.
If $`๐ฐ=\{U_s\}_{sS}`$ is a coarse family in $`X`$, then it has a coarse shrinking $`๐ฑ=\{V_s\}_{sS}`$ such that for any $`M>0`$ there is a bounded subset $`A_M`$ of $`X`$ with the property that $`B(x,M)V_s\mathrm{}`$ implies $`B(x,M)U_s`$ provided $`xXA_M`$.
###### Proof.
Pick $`x_0X`$ and define $`f(x)=\mathrm{min}(d(x,x_0)/2,L_๐ฐ(x)/2)`$. Notice $`f`$ is a coarsely proper function of Lipschitz constant $`1/2`$. For each $`xX`$ pick $`s(x)S`$ so that $`B(x,f(x))U_{s(x)}`$. Define $`V_s`$ as the union of those balls $`B(x,f(x)/2)`$ so that $`s=s(x)`$. It suffices to observe that $`B(x,M)V_s\mathrm{}`$ and $`M<f(x)/3`$ implies $`B(x,M)U_s`$. Indeed, $`B(y,f(y))U_s`$ for some $`yB(x,M)`$. Since $`f(x)f(y)d(x,y)/2<M/2`$, one has $`f(y)>f(x)M/2>3MM/2>2M`$ and $`B(x,M)B(y,f(y))U_s`$. โ
###### Lemma 3.6.
If $`๐ฐ`$ is a coarse family in $`X`$ that is coarsely proper, then there is a coarsely proper function $`f:๐ฐR_+`$ such that the family $`\{B(U,f(U))\}_{U๐ฐ}`$ is coarse.
###### Proof.
Define $`f(U)=inf\{L_๐ฐ(x)/4xU\}`$. Notice $`f`$ is a coarsely proper function. Pick $`s(x)S`$ so that $`B(x,L_๐ฐ(x)/2)U_{s(x)}`$. $`f(U_{s(x)})L_๐ฐ(x)/4`$ which implies $`B(x,L_๐ฐ(x)/4)B(U_{s(x)},f(U_{s(x)}))`$. Thus $`\{B(U,f(U))\}_{U๐ฐ}`$ is coarse. โ
In the large scale geometry one should think of bounded subsets of $`X`$ as points. Here is a generalization of the Lebesque number.
###### Definition 3.7.
Let $`n0`$. Suppose $`๐ฐ`$ is a family in $`X`$ and $`A`$ is a bounded subset of $`X`$. The $`n`$-th Lebesque number $`L^n(๐ฐ,A)`$ is the supremum of $`t[0,\mathrm{}]`$ such that $`๐ฐ|_A`$ has a refinement of multiplicity at most $`n+1`$ and Lebesque number at least $`t`$.
Notice such supremum exists as the cover of $`A`$ consisting of points is of Lebesque number $`0`$ and multiplicity $`1`$.
Observe that $`L^n(๐ฐ,A)`$, $`n0`$, form an increasing sequence of numbers bounded by $`L(๐ฐ,A)`$. If $`๐ฐ|_A`$ is of finite order, then they eventually stabilize and are equal to $`L(๐ฐ|_A,A)`$.
Let us point out that Spernerโs Lemma can be used to estimate higher Lebesque numbers as follows: Consider a 2-simplex $`\mathrm{\Delta }`$ with vertices labeled $`0`$, $`1`$, and $`2`$. Let $`๐ฐ`$ be the cover of $`\mathrm{\Delta }`$ by stars $`U_i`$, $`i=0,1,2`$, of its vertices. Consider a subdivision $`L`$ of $`\mathrm{\Delta }`$ with mesh $`M`$ (in this case it coincides with the longest edge in the subdivision). Let $`X=A`$ be the set of vertices of $`L`$. Suppose $`๐ฑ=\{V_0,V_1,V_2\}`$ is a shrinking of $`๐ฐ|_A`$. Obviously, there is a shrinking of multiplicity $`1`$. However, if we request $`๐ฑ`$ to be of large Lebesque number, we run into problems. Namely, $`L^1(๐ฐ,A)M`$. Indeed, if $`L(๐ฑ)>M`$, we assign to each vertex $`v`$ of $`L`$ number $`i`$ such that $`vV_i`$. We are in the situation of the classical Spernerโs Lemma: vertices on the edges of $`\mathrm{\Delta }`$ must be labeled with a number of one of the vertices of that edge. Therefore one has a simplex in $`L`$ whose vertices were assigned all three numbers $`0,1,2`$. Since $`L(๐ฑ)>M`$, the three vertices belong to $`V_0V_1V_2`$ and multiplicity of $`๐ฑ`$ is $`3`$. Thus $`L^1(๐ฑ,A)M`$.
We will use the observation above in the case of $`M`$-scale connected spaces.
###### Definition 3.8.
Suppose $`M`$ is a positive number. A metric space $`X`$ is called $`M`$-scale connected if for every two points $`x,yX`$ there is a chain of points $`x=x_1,x_2,\mathrm{},x_k=y`$ such that $`d_X(x_i,x_{i+1})<M`$ for all $`i<k`$.
Here is an application of Spernerโs Lemma for $`1`$-simplices.
###### Lemma 3.9.
Let $`M`$ be a positive number and $`X`$ be an $`M`$-scale connected metric space. If $`L^0(๐ฐ,X)>M`$ for some cover $`๐ฐ`$ of $`X`$, then $`๐ฐ`$ contains $`X`$ as an element.
###### Proof.
Suppose $`๐ฑ`$ is a refinement of $`๐ฐ`$ of multiplicity at most $`1`$ and Lebesque number bigger than $`M`$. If $`X`$ is not an element of $`๐ฑ`$, then there are disjoint non-empty elements $`V_1,V_2๐ฑ`$. Pick a chain of points $`x=x_1,x_2,\mathrm{},x_k=y`$ such that $`d_X(x_i,x_{i+1})<M`$ for all $`i<k`$ and $`xV_1`$, $`yV_2`$. There is an index $`j<k`$ such that $`x_jV_1`$ and $`x_{j+1}V_1`$. The ball $`B(x_{j+1},M)`$ is contained in an element $`W`$ of $`๐ฑ`$ and intersects $`V_1`$. Therefore $`W=V_1`$, a contradiction. โ
## 4. The coarse category
Let us introduce the coarse category in a way that explains why two coarse functions are considered equivalent if their distance is bounded.
###### Definition 4.1.
Given a metric space $`(X,d_X)`$ and its two subsets $`X_1`$ and $`X_2`$ the notation $`X_1X_2`$ means there is a positive number $`R`$ such that $`X_1`$ is contained in the ball $`B(X_2,R)=\{xXdist(x,X_2)<R\}`$.
###### Proposition 4.2.
A function $`f:XY`$ of metric spaces is coarse if and only if it preserves the relation $``$ of sets. Thus, $`X_1X_2`$ implies $`f(X_1)f(X_2)`$.
###### Proof.
Suppose $`f:XY`$ preserves the relation $``$ of sets but not in the sense of Roe. Therefore, for some $`M>0`$ there is a sequence of points $`x_n,y_n`$ so that $`d_X(x_n,y_n)<M`$ for each $`n`$ but $`d_Y(f(x_n),f(y_n))\mathrm{}`$ as $`n\mathrm{}`$. If $`f(A)`$ is bounded for some subsequence $`A`$ of $`x_n`$, then $`f(B)`$ is bounded for the corresponding subsequence of $`y_n`$ (in view of $`f(B)f(A)`$) contradicting $`d_Y(f(x_n),f(y_n))\mathrm{}`$ as $`n\mathrm{}`$. Thus $`f(x_n)\mathrm{}`$ and $`f(y_n)\mathrm{}`$ as $`n\mathrm{}`$. By induction define a subsequence $`a_n`$ of $`\{x_n\}_{n1}`$ and the corresponding subsequence $`b_n`$ of $`\{y_n\}_{n1}`$ with the property that $`d_Y(f(a_k),f(b_i))>k`$ and $`d_Y(f(b_k),f(a_i))>k`$ for all $`ki`$. Since $`A=\{a_n\}_{n1}B=\{b_n\}_{n1}`$ one has $`f(A)f(B)`$, a contradiction.
Suppose $`f:XY`$ is coarse in the sense of Roe and $`X_1X_2`$ in $`X`$. Pick $`R>0`$ so that $`X_1B(X_2,R)`$ and choose $`M>0`$ satisfying $`d_Y(f(x),f(y))<M`$ if $`d_X(x,y)<R`$ for all $`x,yX`$. Given $`xX_1`$ pick $`yX_2`$ so that $`d_X(x,y)<R`$ since $`d_Y(f(x),f(y))<M`$ one gets $`f(X_1)B(f(X_2),M)`$. Thus $`f(X_1)f(X_2)`$. โ
Notice that $`X_1X_2`$ for every bounded subset $`X_1`$ of $`X`$ provided $`X_2\mathrm{}`$. Also, $`X_1X_2`$ implies $`X_1`$ is bounded provided $`X_2`$ is bounded. Therefore $`f(A)`$ is bounded for every bounded subset $`A`$ of $`X`$ and every coarse function $`f:XY`$.
Given a function $`f:XY`$ of metric spaces one can identify it with its graph $`\mathrm{\Gamma }(f)X\times Y`$. Therefore it makes sense to ponder the meaning of $`\mathrm{\Gamma }(f)\mathrm{\Gamma }(g)`$ for $`f,g:XY`$.
###### Proposition 4.3.
Suppose $`f,g:XY`$ are functions of metric spaces.
* If $`g`$ is coarse, then $`\mathrm{\Gamma }(f)\mathrm{\Gamma }(g)`$ implies that the distance $`dist(f,g)`$ between $`f`$ and $`g`$ is finite. In particular, $`f`$ is coarse.
* If $`dist(f,g)`$ is finite, then $`\mathrm{\Gamma }(f)\mathrm{\Gamma }(g)`$.
###### Proof.
1. Suppose the distance $`dist(f,g)`$ is not finite, so there are points $`x_nX`$ with $`d_Y(f(x_n),g(x_n)>n`$ for all $`n1`$. Let $`R>0`$ be a number such that $`B(\mathrm{\Gamma }(g),R)`$ contains $`\mathrm{\Gamma }(f)`$. For each $`n`$ pick $`y_nX`$ satisfying $`d_X(x_n,y_n)+d_Y(f(x_n),g(y_n))<R`$. There is $`M>0`$ so that $`d_Y(g(x_n),g(y_n))<M`$ for all $`n1`$ as $`g`$ is coarse. Now, $`d_Y(f(x_n),g(x_n)d_Y(f(x_n),g(y_n))+d_Y(g(y_n),g(x_n))<R+M`$ for all $`n1`$, a contradiction.
2. Notice $`\mathrm{\Gamma }(f)B(\mathrm{\Gamma }(g),dist(f,g))`$. โ
###### Definition 4.4.
Given a function $`f:XY`$ of metric spaces we define the forward distance transfer function $`d_f:R_+R_+\mathrm{}`$ as the infimum of all functions $`\alpha :R_+R_+\mathrm{}`$ with the property that $`d_X(x,y)t`$ implies $`\alpha (t)d_Y(f(x),f(y))`$ for all $`x,yX`$.
The reverse distance transfer function $`d^f:R_+R_+\mathrm{}`$ as the infimum of all functions $`\alpha :R_+R_+\mathrm{}`$ with the property that $`d_Y(f(x),f(y)))t`$ implies $`d_X(x,y)\alpha (t)`$ for all $`x,yX`$.
Notice that $`f`$ is coarse if and only if $`d_f`$ maps $`R_+`$ to $`R_+`$, i.e. the values of $`d_f`$ are finite. Also, $`f`$ is asymptotically Lipschitz if and only if $`d_f`$ is bounded by a linear function.
###### Proposition 4.5.
If $`f,g:XY`$ are two coarsely proper coarse functions, then the following conditions are equivalent:
* $`dist(f,g)`$ is finite.
* For every coarse family $`๐ฐ=\{U_s\}_{sS}`$ in $`Y`$ the family $`\{f^1(U_s)g^1(U_s)\}_{sS}`$ is coarse.
###### Proof.
1$``$2. Let $`dist(f,g)<M`$. Consider $`๐ฑ=\{B(U_s,M)\}_{sS}`$. It is a coarse family, so $`f^1(๐ฑ)`$ is coarse by 2.12. Notice $`f^1(B(U_s,M))f^1(U_s)g^1(U_s)`$ for all $`sS`$ which is sufficient to establish coarseness of $`\{f^1(U_s)g^1(U_s)\}_{sS}`$.
2$``$1. If $`dist(f,g)`$ is not finite, there is a sequence $`x_n\mathrm{}`$ such that $`d_Y(f(x_n),g(x_n))>n`$ for all $`n`$. Put $`A=\{x_n\}_{n1}`$. By 2.7, the family $`๐ฐ=\{Yf(A),Yg(A)\}`$ is coarse. However, $`\{f^1(U_s)g^1(U_s)\}_{sS}`$ is not coarse as it refines $`\{XA\}`$ which is not coarse. โ
Our category is that of metric spaces and equivalence classes of coarse functions. $`fg`$ if $`d_Y(f(x),g(x))`$ is a bounded function of $`x`$.
Generalizing the concept of $`AB`$ for subsets of a given metric space $`X`$, we say $`Y`$ coarsely dominates $`X`$ (notation: $`X_{rse}^{coa}Y`$) if there are coarse functions $`f:XY`$ and $`g:YX`$ such that $`gf`$ is at a finite distance from $`id_X`$.
###### Proposition 4.6.
Suppose $`f:XY`$ and $`g:YX`$ are coarse functions. If $`gf`$ is at a finite distance from $`id_X`$, then both $`f:Xf(X)`$ and $`g:f(X)X`$ are coarsely proper and $`fg`$ is at finite distance from $`id_{f(X)}`$.
###### Proof.
Suppose $`x_n\mathrm{}`$. None of the subsequences of $`\{f(x_n)\}`$ can be bounded as $`g`$ would send it to a bounded subset of $`X`$. Thus $`f(x_n)\mathrm{}`$. If $`f(x_n)\mathrm{}`$, then none of subsequences of $`\{x_n\}`$ is bounded. Therefore none of the subsequences of $`\{g(f(x_n))\}`$ is bounded and $`g:f(X)X`$ is coarsely proper. If $`d_X(g(f(x)),x)<M`$ for all $`xX`$, then $`d_Y(f(g(f(x))),f(x))d_f(M)`$ and $`fg`$ is at finite distance from $`id_{f(X)}`$. โ
###### Proposition 4.7.
A surjective coarse function $`f:XY`$ of metric spaces is a coarse isomorphism if and only if the reverse distance transfer function $`d^f`$ is finite.
###### Proof.
If there is a coarse function $`g:YX`$ such that $`gf`$ is at finite distance $`M`$ to $`id_X`$, then $`d^f(a)d_g(a)+2M`$ is finite.
Assume $`d^f`$ is finite and pick a right inverse $`g:YX`$. Notice $`d_X(g(x),g(y))d^f(d_Y(x,y))`$, so $`g`$ is coarse. โ
## 5. Coarse dimensions
###### Definition 5.1.
The coarse dimension $`\text{dim}_{rse}^{coa}(X)`$ (respectively, the major coarse dimension $`\text{dim}_{RSE}^{COA}(X)`$) is the smallest integer $`n`$ such that any finite coarse family in $`X`$ (respectively, any coarse family in $`X`$) has a coarse refinement with multiplicity at most $`n+1`$.
###### Remark 5.2.
Using Proposition 4.4 on p.1104 in (notice that the words โuniformly boundedโ are erroneously inserted there) one can show that, for proper metric spaces $`X`$, the major coarse dimension of $`X`$ coincides with the asymptotic dimension of Dranishnikov. In view of 8.2, our coarse dimension and Dranishnikov coarse dimension are identical.
Given a coarse family $`๐ฐ=\{U_s\}_{sS}`$ in a subset $`A`$ of $`X`$ one can extend it to a coarse family $`๐ฐ^{}=\{U_s(XA)\}_{sS}`$ in $`X`$. Notice that $`๐ฑA`$ is a coarse refinement of $`๐ฐ`$ for any coarse refinement $`๐ฑ`$ of $`๐ฐ^{}`$. Therefore the following holds.
###### Corollary 5.3.
If $`A`$ is a subset of a metric space $`X`$, then $`\text{dim}_{rse}^{coa}(A)\text{dim}_{rse}^{coa}(X)`$ and $`\text{dim}_{RSE}^{COA}(A)\text{dim}_{RSE}^{COA}(X)`$.
###### Proposition 5.4.
If $`Y`$ coarsely dominates $`X`$, then $`\text{dim}_{rse}^{coa}(X)\text{dim}_{rse}^{coa}(Y)`$ and $`\text{dim}_{RSE}^{COA}(X)\text{dim}_{RSE}^{COA}(Y)`$.
###### Proof.
The proof is almost the same for both dimensions. Suppose $`๐ฐ`$ is a coarse family in $`X`$ and $`f:XY`$, $`g:YX`$ are coarse functions such that there is $`M>0`$ with $`d_X(x,g(f(x)))<M`$ for all $`xX`$. Replacing $`Y`$ by $`f(X)`$ we may assume $`f`$ is onto and both $`f`$ and $`g`$ are coarsely proper (see 4.6). The idea of the proof is to refine $`g^1(๐ฐ)`$ by $`๐ฑ`$ and then refine $`f^1(๐ฑ)`$ to obtain a desired refinement $`๐ฒ`$ of $`๐ฐ`$ of multiplicity at most $`n+1`$, where $`n`$ is the dimension of $`Y`$. Consider $`๐ฐ^{}=\{B(U_s,M)\}_{sS}`$. It is a coarse family in $`X`$, so $`\{g^1(B(U_s,M))\}_{sS}`$ is coarse and it has a coarse shrinking $`๐ฑ=\{V_s\}_{sS}`$ of multiplicity at most $`n+1`$. Suppose $`xf^1(V_s)U_s`$. Since $`d_X(x,g(f(x)))<M`$, $`g(f(x))B(U_s,M)`$. However, $`f(x)V_sg^1(B(U_s,M))`$, a contradiction. โ
###### Definition 5.5.
The minor asymptotic dimension $`ad(X)`$ (respectively, the asymptotic dimension $`asdim(X)`$) is the smallest integer $`n`$ such that the function $`๐ฐL^n(๐ฐ,X)`$ is coarsely proper on the space of finite covers (respectively, arbitrary covers) $`๐ฐ`$ of $`X`$.
Let us show that our definition of asymptotic dimension is equivalent to that of Gromov.
###### Proposition 5.6.
$`asdim(X)n`$ if and only if for each $`M>0`$ there is a uniformly bounded family $`๐ฐ`$ in $`X`$ of Lebesque number at least $`M`$ and multiplicity at most $`n+1`$.
###### Proof.
If $`asdim(X)n`$ as in 5.5 and $`M>0`$, then there is $`N>0`$ such that every cover $`๐ฑ`$ of $`X`$ satisfying $`L(๐ฑ,X)N`$ has a refinement $`๐ฐ`$ of multiplicity at most $`n+1`$ and Lebesque number at least $`M`$. Pick $`๐ฑ`$ to be the cover of $`X`$ by balls of radius $`N`$. The resulting $`๐ฐ`$ is uniformly bounded.
Suppose for each $`M>0`$ there is a uniformly bounded family $`๐ฐ^M`$ of multiplicity at most $`n+1`$ and Lebesque number at least $`M`$. Let $`\alpha (M)`$ be the supremum of diameters of elements of $`๐ฐ^M`$. Given any family $`๐ฑ`$ of Lebesque number at least $`\alpha (M)+1`$, $`๐ฐ^M`$ is a refinement of of $`๐ฑ`$ which proves that the function $`๐ฑL^n(๐ฑ,X)`$ is coarsely proper on the space of all covers $`๐ฑ`$ of $`X`$. โ
Quite often it is useful to have even stronger conditions imposed on covers appearing in 5.6.
###### Proposition 5.7 (Gromov).
If Gromov asymptotic dimension $`asdim(X)`$ does not exceed $`n`$, then for any $`M,N>0`$ there exist uniformly bounded families $`๐ฐ^i`$, $`1in+1`$, such that each $`๐ฐ^i`$ is $`N`$-disjoint and $`๐ฐ=\underset{i=1}{\overset{n+1}{}}๐ฐ^i`$ is of Lebesque number at least $`M`$.
###### Proof.
Consider a uniformly bounded family $`๐ฑ=\{V_s\}_{ss}`$ of multiplicity at most $`n+1`$ and Lebesque number at least $`2(n+1)(M+N)`$. Lemma 3.3 says it can be refined by $`๐ฑ^{}=\underset{i=1}{\overset{n+1}{}}๐ฑ^i`$ such that $`L_๐ฑ(x)L_๐ฐ(x)/(2n+2)M+N`$ for each $`xX`$ and each $`๐ฑ^i`$ consists of disjoint sets. Define $`๐ฐ_i`$ as $`\{B(W,N)\}`$, $`W๐ฑ^i`$. โ
Let us characterize spaces of asymptotic dimension $`0`$.
###### Proposition 5.8.
$`asdim(X)>0`$ if and only if there exist a number $`M>0`$ and a coarsely proper sequence $`\{(x_n,y_n)\}_{n=1}^{\mathrm{}}`$ of pairs of points in $`X`$ such that $`dist(x_n,y_n)\mathrm{}`$ and the points $`x_n`$ and $`y_n`$ can be $`M`$-scale connected in $`XB(x_0,n)`$.
###### Proof.
If $`asdim(X)=0`$, then for any $`M>0`$ there exists an $`M`$-disjoint cover of $`X`$ by uniformly bounded sets. Therefore, the distance between two points $`x`$ and $`y`$ which can be $`M`$-scale connected in $`X`$ is uniformly bounded.
Suppose $`asdim(X)>0`$. Let $`n`$ be a positive integer and $`x_0`$ be the base point in $`X`$. There is $`L>0`$ such that $`X`$ does not have a uniformly bounded cover of Lebesque number bigger than $`L`$ and multiplicity 1. Define an equivalence relation on $`XB(x_0,n)`$ by saying $`xy`$ iff $`x`$ and $`y`$ can be $`2L`$-scale connected in $`XB(x_0,n)`$. The cover of $`X`$ by the equivalence classes has Lebesque number at least $`2L`$, therefore these classes are not uniformly bounded by the choice of $`L`$. Thus, there exist points $`x_n`$ and $`y_n`$ which can be $`2L`$-scale connected in $`XB(x_0,n)`$ such that $`dist(x_n,y_n)`$ is arbitrarily large. โ
###### Proposition 5.9.
If $`Y`$ coarsely dominates $`X`$, then $`asdim(X)asdim(Y)`$ and $`ad(X)ad(Y)`$.
###### Proof.
The proof is almost the same for both dimensions. Suppose $`๐ฐ`$ is a coarse family in $`X`$ and $`f:XY`$, $`g:YX`$ are coarse functions such that there is $`M>0`$ with $`d_X(x,g(f(x)))<M`$ for all $`xX`$. By replacing $`Y`$ with $`f(X)`$ we may assume $`f`$ is onto and both $`f`$ and $`g`$ are coarsely proper (see 4.6). The idea of the proof is to refine $`g^1(๐ฐ)`$ by $`๐ฑ`$ and then refine $`f^1(๐ฑ)`$ to obtain a desired refinement $`๐ฒ`$ of $`๐ฐ`$ of multiplicity at most $`n+1`$, where $`n`$ is the dimension of $`X`$. Take a coarsely proper function $`\alpha :R_+R_+`$ with the property that any finite cover (respectively, arbitrary cover) $`๐ฐ`$ of $`Y`$ satisfying $`L(๐ฐ,Y)\alpha (t)`$ has a refinement $`๐ฑ`$ of multiplicity at most $`n+1`$ so that $`L(๐ฑ,Y)t`$.
Given $`t>0`$ pick $`\beta (t)`$ so that $`L^g(\beta (t))>\alpha (t)`$ (see 2.9). Assume $`L(๐ฐ)>M+\beta (t)`$. Consider $`๐ฐ^{}=\{B(U_s,M)\}_{sS}`$. $`L(๐ฐ^{})>\beta (t)`$, so $`\{g^1(B(U_s,M))\}_{sS}`$ is of Lebesque number at least $`\alpha (t)`$ and it has a shrinking $`๐ฑ=\{V_s\}_{sS}`$ of multiplicity at most $`n+1`$ and $`L(๐ฑ)t`$. Suppose $`xf^1(V_s)U_s`$. Since $`d_X(x,g(f(x)))<M`$, $`g(f(x))B(U_s,M)`$. However, $`f(x)V_sg^1(B(U_s,M))`$, a contradiction. โ
###### Theorem 5.10.
The major coarse dimension of $`X`$ does not exceed the asymptotic dimension of $`X`$.
###### Proof.
Suppose $`asdim(X)=n<\mathrm{}`$ and $`๐ฐ=\{U_s\}_{sS}`$ is a coarse family in $`X`$. By Lemma 3.4 we may assume $`U`$ is coarsely proper. By induction on $`k`$ find a sequence of numbers $`M_0=1,M_1,M_2,\mathrm{}`$, and covers $`๐ฑ^k=\{V_t\}_{tT(k)}`$, $`k1`$, of multiplicity at most $`n+1`$ and satisfying the following conditions:
a. $`L(๐ฑ^k,X)M_{k1}`$ for $`k1`$.
b. The diameter of each element of $`๐ฑ^k`$ is smaller than $`M_k`$.
c. The family $`\{B(x,M_{k1})d(x,x_0)M_k\}`$ refines $`๐ฐ`$ for each $`k1`$.
d. $`M_{k+1}>2M_k`$ for all $`k1`$.
Find functions $`j(k):T(k)T(k+1)`$ so that $`V_tV_{j(k)(t)}`$. Denote $`\{x:M_kd(x,x_0)<M_{k+1}\}`$ by $`A_k`$. Given $`tT(k)`$ so that $`V_t`$ is contained in some element of $`๐ฐ`$ define $`\alpha (t)S`$ by looking at the sequence $`V_tV_{j(k)(t)}\mathrm{}`$, picking the latest element contained in some $`U_s`$ and setting $`\alpha (t)=s`$ (it is possible each element of the sequence is contained in some $`U_s`$ in which case all of them are contained in some $`U_s`$ and that $`s`$ is picked as $`\alpha (t)`$). Define $`W_s`$ as follows: it is the union of non-empty sets of the form $`V_tA_k`$ so that $`V_t๐ฑ^{k1}`$ and $`\alpha (t)=s`$. Notice that $`m(๐ฒ)n+1`$ as in the annulus $`A_k`$ the family $`๐ฒ`$ is obtained from $`๐ฑ^{k1}`$ by assembling some of its elements together.
We plan to show $`๐ฒ`$ is coarse by proving that if $`M_kd(x,x_0)<M_{k+1}`$, then $`B(x,M_{k3})`$ is contained in some $`W_s`$. Indeed, there is $`tT(k2)`$ so that $`B(x,M_{k3})V_t`$. Put $`r=j(k2)(t)`$ and $`u=j(k1)(r)`$. Points of $`B(x,M_{k3})`$ can belong to only two of the following three annuli: $`A_{k1}`$, $`A_k`$, and $`A_{k+1}`$. If $`zB(x,M_{k3})A_{k+1}`$, then $`V_uB(z,M_k)U_s`$ for some $`sS`$. We might as well put $`s=\alpha (t)=\alpha (u)=\alpha (r)`$. In this case $`B(x,M_{k3})W_s`$. If $`B(x,M_{k3})`$ misses the last annulus, then only $`\alpha (r)`$ is definitely defined ($`\alpha (u)`$ may not exist) and $`\alpha (t)=\alpha (r)`$. Now, $`B(x,M_{k3})W_s`$, where $`s=\alpha (r)`$. โ
###### Remark 5.11.
5.10 generalizes Proposition 4.5 on p.1105 of .
## 6. The large scale dimension
In this section we prove that any dimension of $`X`$ (asymptotic, major coarse, or minor asymptotic), if finite, equals the coarse dimension of $`X`$. That corresponds to results of Dranishnikov that $`asdim(X)`$ or $`asdim^{}(X)`$, if finite, are equal to the dimension of the Higson corona of any proper metric space $`X`$. Our proofs are direct and become simpler by introducing a new dimension, the large scale dimension of $`X`$. That dimension turns out to be identical with the coarse dimension.
###### Definition 6.1.
The large scale dimension $`\text{dim}_{scale}^{large}(X)`$ of $`X`$ is the smallest integer $`n`$ such that $`AL^n(๐ฐ,A)`$ is a coarsely proper function on the set of bounded subsets of $`X`$ for all finite coarse families $`๐ฐ`$ in $`X`$.
Notice $`\text{dim}_{scale}^{large}(X)=1`$ if $`X`$ is bounded.
Obviously, $`\text{dim}_{scale}^{large}(X)\text{dim}_{scale}^{large}(A)`$ for any subset $`A`$ of $`X`$.
###### Proposition 6.2.
$`ad(X)\text{dim}_{scale}^{large}(X)`$ and $`\text{dim}_{rse}^{coa}(X)\text{dim}_{scale}^{large}(X)`$.
###### Proof.
The inequality $`\text{dim}_{rse}^{coa}(X)\text{dim}_{scale}^{large}(X)`$ is almost obvious. Indeed, given $`n=\text{dim}_{rse}^{coa}(X)`$ and given a coarse family $`๐ฐ`$ in $`X`$ consisting of $`m`$ elements one has a coarse refinement $`๐ฑ`$ of $`๐ฐ`$ such that the multiplicity $`m(๐ฑ)`$ is at most $`n+1`$. In that case
$$L^n(๐ฐ,A)L(๐ฑ,A)\underset{aA}{inf}L_๐ฑ(a)$$
and is a coarsely proper function of $`A`$.
Suppose $`ad(X)=n`$ and $`๐ฐ`$ is a coarse cover of $`X`$ consisting of $`m`$ elements. Given $`t>0`$ find a bounded subset $`U`$ of $`X`$ such that $`๐ฐ|_{(XU)}`$ has a refinement $`๐ฑ`$ of multiplicity at most $`n+1`$ and Lebesque number at least $`t`$. For any bounded subset $`A`$ of $`XU`$, $`L^n(๐ฐ,A)L(๐ฑ,A)t`$ which proves $`\text{dim}_{scale}^{large}(X)n`$. โ
As shown in , the asymptotic dimension of $`R^n`$ is at most $`n`$ (see p.793). For the convenience of the reader let us reword the argument from as follows: Given $`M>0`$ consider the triangulation on the unit $`n`$-cube $`I^n`$ obtained by starring at the center of each face. It is invariant under symmetries of $`I^n`$ and the cover of $`I^n`$ by stars of vertices has a positive Lebesque number $`k`$ and is of multiplicity at most $`n+1`$. Rescale $`I^n`$ by the factor of $`M/k`$ and extend its triangulation over the whole $`R^n`$ by reflections. The cover of $`R^n`$ by stars of vertices has Lebesque number at least $`M`$ and is of multiplicity at most $`n+1`$.
Let us show how to use the large scale dimension to estimate asymptotic dimension from below.
###### Proposition 6.3.
$`\text{dim}_{scale}^{large}(R^n)n`$.
###### Proof.
Since $`\text{dim}(I^n)=n`$, there is a finite open cover $`๐ฐ`$ of $`I^n`$ with no open refinement of multiplicity at most $`n`$. Let $`I_k^nR^n`$ be a copy of $`I^n`$ enlarged $`k`$ times with the corresponding cover $`๐ฐ^k`$. We request $`I_k^n\mathrm{}`$ so that $`๐ฑ`$ obtained by adding the corresponding elements of $`๐ฐ^k`$ is a finite coarse family on $`A=\underset{k=1}{\overset{\mathrm{}}{}}I_k^n`$. Notice $`L^{n1}(๐ฑ,I_k^n)=0`$ for all $`k`$. Thus $`\text{dim}_{scale}^{large}(R^n)n`$. โ
###### Proposition 6.4.
If $`\text{dim}_{scale}^{large}(X)=0`$, then $`asdim(X)=0`$ and $`\text{dim}_{RSE}^{COA}(X)=0`$.
###### Proof.
It suffices to show $`asdim(X)=0`$ (see 5.10). Suppose $`asdim(X)>0`$. By 5.8 there exist a number $`M>0`$ and a coarsely proper sequence $`\{(x_n,y_n)\}_{n=1}^{\mathrm{}}`$ of pairs of points in $`X`$ such that $`dist(x_n,y_n)\mathrm{}`$ and the points $`x_n`$ and $`y_n`$ can be $`M`$-scale connected in $`XB(x_0,n)`$ by a chain $`P_n`$. Consider a coarse family $`๐ฐ`$ consisting of two sets: $`X\underset{n=1}{\overset{\mathrm{}}{}}\{x_n\}`$ and $`X\underset{n=1}{\overset{\mathrm{}}{}}\{y_n\}`$.
Since $`CL^0(๐ฐ,C)`$ is a coarsely proper function, there is a chain $`P_n`$ such that $`L^0(๐ฐ,P_n)>M`$. This contradicts 3.9 since $`P_n`$ is $`M`$-scale connected and the cover $`๐ฐ`$ is non-trivial on $`P_n`$. โ
###### Definition 6.5.
Given a point-finite family $`๐ฐ=\{U_s\}_{sS}`$ in $`X`$ (that means each point of $`X`$ belongs to at most finitely many elements of $`๐ฐ`$) by the canonical partition of unity of $`๐ฐ`$ we mean the family of functions $`\{f_s/f\}_{sS}`$, where $`f_s(x)=dist(x,XU_s)`$ and $`f(x)=\underset{sS}{}f_s(x)`$. If $`T`$ is a subset of $`S`$, then $`X_T`$ is defined to be $`\{xX\underset{sT}{}f_s(x)/f(x)=1\}`$ and by $`X_T`$ we mean the set of all $`xX_T`$ such that $`f_s(x)=0`$ for some $`sT`$.
Notice that $`f(x)>0`$ for all $`xX`$ such that $`L_๐ฐ(x)>0`$ and $`f`$ is a Lipschitz function if $`๐ฐ`$ is of finite multiplicity.
###### Lemma 6.6.
If the large scale dimension of $`X`$ is at most $`n`$, then any coarse family $`๐ฐ`$ in $`X`$ of finite multiplicity $`m`$ has a coarse refinement $`๐ฑ`$ of multiplicity at most $`n+1`$.
###### Proof.
Suppose $`๐ฐ`$ exists with no coarse refinement of multiplicity at most $`n+1`$. Using 3.4 we reduce the general case to that of $`๐ฐ=\{U_s\}_{sS}`$ consisting of bounded sets so that for any sequence $`x_k\mathrm{}`$ the conditions $`x_kU_{s(k)}๐ฐ`$ imply $`U_{s(k)}\mathrm{}`$. For induction on $`mn`$ it suffices to assume the multiplicity of $`๐ฐ`$ is $`n+2`$.
Pick a coarse shrinking $`๐ฒ=\{W_s\}_{sS}`$ (see 3.5) so that given $`M>0`$ there is a bounded subset $`A`$ of $`X`$ with the property that, for $`xXA`$, $`B(x,M)W_s\mathrm{}`$ implies $`B(x,M)U_s`$. Consider the canonical partition of unity $`f`$ of $`๐ฒ`$. Given a set $`T`$ in $`S`$ consisting of $`n+2`$ elements pick a shrinking $`๐ฒ^T`$ of $`๐ฒ|_{X_T}`$ of order at most $`n+1`$ and the Lebesque number at least half the maximum $`L^n(๐ฒ^T,X_T)`$ possible (if the maximum is infinity we pick a shrinking of Lebesque number twice the size of $`X_T`$). We can add $`W_sX_T`$ to $`W_s^T`$ without increasing the order of $`W^T`$ beyond $`n+1`$ (obviously, the Lebesque number does not decrease). By pasting those shrinkings for all $`T`$ one gets a refinement $`๐ฑ`$ of $`๐ฒ`$ on $`XA`$ for some bounded subset $`A`$ of $`X`$ of multiplicity at most $`n+1`$. Therefore $`๐ฑ`$ cannot be coarse and there is $`M>0`$ and a sequence of points $`x_k\mathrm{}`$ such that none of $`B(x_k,M)`$ is contained in an element of $`๐ฑ`$. In particular $`B(x_k,M)`$ is not contained in the $`n`$-skeleton of $`X`$ (the points where the order of $`f`$ is at most $`n+1`$) for large $`k`$.
Pick sets $`T(k)`$ so that $`X_{T(k)}X_{T(k)}`$ contains an element $`y_kB(x_k,M)`$. For large $`k`$, $`B(x_k,M)`$ intersecting $`W_s`$ implies $`B(x_k,M)U_s`$. Therefore the set $`T`$ of $`sS`$ so that $`B(x_k,M)`$ intersects $`W_s`$ is of cardinality at most $`n+2`$ and $`B(x_k,M)X_{T(k)}`$. For large $`k`$ the cover $`๐ฒ|_{X_{T(k)}}`$ has a refinement of order at most $`n+1`$ and Lebesque number at least $`3M`$. Therefore, $`B(x_k,M)`$ is contained in a single element of $`๐ฑ`$, a contradiction. โ
###### Corollary 6.7.
The coarse dimension of $`X`$ equals the large scale dimension of $`X`$.
###### Corollary 6.8.
If the major coarse dimension of $`X`$ is finite, then it equals the large scale dimension of $`X`$.
###### Theorem 6.9.
If the asymptotic dimension (respectively, the minor asymptotic dimension) of unbounded $`X`$ is finite, then it equals the large scale dimension of $`X`$.
###### Proof.
Suppose $`asdim(X)=n`$ (respectively, $`ad(X)=n`$) and $`\text{dim}_{scale}^{large}(X)<n`$. Notice $`n>0`$ as $`\text{dim}_{scale}^{large}(X)<0`$ is possible only for bounded $`X`$. Therefore there is $`M>0`$ and a sequence of covers (respectively, finite covers) $`๐ฐ^k`$ indexed by sets $`S(k)`$ of Lebesque number at least $`k+3M`$ and multiplicity at most $`n+1`$ so that no refinement of $`๐ฐ^k`$ of multiplicity $`n`$ has Lebesque number bigger than $`M`$. Augment each $`๐ฐ^k`$ by shrinking it to the family $`B(U,M)`$, $`U๐ฐ^k`$. Let $`f^k`$ be the canonical partition of unity of that augmentation.
Notice that for any $`k`$ and any $`xX`$ there is a subset $`T`$ of $`S(k)`$ consisting of at most $`(n+1)`$ elements so that $`B(x,M)X_T`$. We are going to show that for every $`k`$ there is $`N>0`$ such that for any $`R>N`$ there is $`T(k)S(k)`$ consisting of at most $`(n+1)`$ elements with $`X_{T(k)}XB(x_0,R)`$, $`x_0`$ a fixed point in $`X`$, so that $`Carr(f^k|_{X_{T(k)}})`$ does not admit a refinement of multiplicity at most $`n`$ and Lebesque number bigger than $`M`$.
Suppose that, for some $`k`$ and $`R>0`$, all $`Carr(f^k|_{X_T})`$ so that $`X_TXB(x_0,R)`$ do admit a refinement $`๐ฑ(T)`$ of multiplicity at most $`n`$ and Lebesque number bigger than $`M`$. By converting those refinements to shrinkings and pasting one gets a refinement $`๐ฑ`$ of $`๐ฐ^k`$ on $`XU`$ for some bounded subset $`U`$ of $`X`$ of multiplicity at most $`n`$ and Lebesque number bigger than $`M`$. More precisely, for each $`TT(k)`$ so that $`X_TXB(x_0,R)`$, we pick a shrinking $`\{V_t^T\}_{tT}`$ of $`Carr(f^k|_{X_T})`$ of multiplicity at most $`n`$ and Lebesque number bigger than $`M`$. If $`T`$ contains at most $`n`$ elements, that shrinking is picked to be exactly $`Carr(f^k|_{X_T})`$ as the multiplicity is at most $`n`$ in such case. $`๐ฑ`$ is a shrinking of $`๐ฐ^k|_{(XU)}`$, $`U`$ being the union of $`X_T`$ that are not contained in $`XB(x_0,R)`$, and $`V_s`$, $`sS(k)`$, is defined as the union of all $`V_s^T`$ with $`sT`$. The reason $`๐ฑ`$ has Lebesque number at least $`M`$ is that for any $`xX`$ there is a subset $`T`$ of $`S(k)`$ consisting of at most $`(n+1)`$ elements so that $`B(x,M)X_T`$.
Now, the cover consisting of the union of $`B(U,2M)`$ and all the elements of $`๐ฑ`$ intersecting $`B(U,2M)`$ and of all elements of $`๐ฑ`$ that do not intersect $`B(U,2M)`$ is uniformly bounded, of multiplicity at most $`n`$ (recall $`n>0`$), and of Lebesque number bigger than $`M`$, a contradiction.
Construct by induction a sequence of sets $`T(i)S(i)`$ with $`X_{T(i)}`$ being mutually disjoint and tending to infinity so that $`Carr(f^i|_{X_{T(i)}})`$ does not have a refinement of multiplicity at most $`n`$ and Lebesque number bigger than $`M`$. Paste all those carriers according to their index within each set $`T(i)`$ and get a coarse cover on a subset $`A`$ of $`X`$ that does not admit a refinement of multiplicity at most $`n`$ and Lebesque number bigger than $`M`$ on infinitely many $`X_{T(i)}`$, a contradiction. โ
## 7. Slowly oscillating functions
###### Definition 7.1.
A function $`f:XY`$ is slowly oscillating if $`f^1(๐ฐ)`$ is coarse for every cover $`๐ฐ`$ of $`Y`$ of positive Lebesque number.
###### Definition 7.2.
Given a function $`f:XY`$ of metric spaces one defines its oscillation function $`Osc(f,M):XR_+\mathrm{}`$ for every $`M>0`$ by declaring $`Osc(f,M)(a)`$ to be the supremum of $`d_Y(f(x),f(a))`$ over all $`xB(a,M)`$.
###### Proposition 7.3.
$`f`$ is slowly oscillating if and only if $`Osc(f,M)(x)0`$ as $`x\mathrm{}`$ for all $`M>0`$.
###### Proof.
Suppose $`Osc(f,M)(x)0`$ as $`x\mathrm{}`$ for all $`M>0`$. Given a cover $`๐ฐ`$ of $`Y`$ of positive Lebesque number and given $`x_n\mathrm{}`$ in $`X`$ there is $`N>0`$ such that each $`f(B(x_n,M))`$ is of diameter smaller that $`L(๐ฐ,Y)`$ for $`n>N`$. Therefore $`B(x_n,M)`$ is contained in an element of $`f^1(๐ฐ)`$ and $`f^1(๐ฐ)`$ is coarse.
Suppose $`f^1(๐ฐ)`$ is coarse for every cover $`๐ฐ`$ of $`Y`$ of positive Lebesque number. Given $`x_n\mathrm{}`$ in $`X`$ and given $`M>0`$ such that diameters of $`f(B(x_n,M))`$ are bigger than a fixed $`\delta >0`$, consider $`๐ฐ=\{B(y,\delta /2)\}_{yY}`$. Since $`f^1(๐ฐ)`$ is coarse, there is $`N>0`$ such that for all $`n>N`$ sets $`B(x_n,M)`$ are contained in an element of $`f^1(๐ฐ)`$. Therefore diameters of $`f(B(x_n,M))`$ are smaller than a $`\delta `$ for $`n>N`$, a contradiction. โ
Our basic way of constructing slowly oscillating real-valued functions is based on the following.
###### Lemma 7.4.
Suppose $`f,g:XR_+`$ and $`Osc(f,M),Osc(g,M)<ฯต`$ for some $`ฯต>0`$. If $`f(x)+g(x)>N`$ for all $`xX`$, then $`Osc(\frac{f}{f+g},M)<\frac{3ฯต}{N}`$.
###### Proof.
Let $`h=\frac{f}{f+g}`$ and $`a=\frac{3ฯต}{N}`$. If $`h(x)h(y)a`$ for some $`x,yX`$ satisfying $`d_X(x,y)<M`$, then $`\frac{f(x)}{f(x)+g(x)}\frac{f(x)ฯต}{f(x)+g(x)+2ฯต}a`$ as well. Since $`\frac{f(x)}{f(x)+g(x)}\frac{f(x)ฯต}{f(x)+g(x)+2ฯต}=\frac{f(x)2ฯต+ฯต(f(x)+g(x))}{(f(x)+g(x))(f(x)+g(x)+2ฯต)}\frac{3ฯต}{f(x)+g(x)+2ฯต}<a`$, we arrive at a contradiction. โ
###### Corollary 7.5.
If $`f`$ and $`g`$ are coarse functions from $`X`$ to $`R_+`$ such that $`f+g`$ is coarsely proper and positive, then $`f/(f+g)`$ is slowly oscillating.
Here is a simple connection between oscillation and the Lebesque number.
###### Lemma 7.6.
If $`\varphi =\{\varphi _s:XR_+\}_{sS}`$ is a family of functions with finite supremum $`sup(\varphi )`$ such that $`Osc(\varphi _s,M)<\frac{1}{2}sup(\varphi )`$ for each $`sS`$, then $`L(\varphi )M`$.
###### Proof.
Given $`aX`$ find $`sS`$ so that $`\varphi _s(a)>\frac{1}{2}sup(\varphi )(a)`$. If $`d_X(x,a)<M`$, then $`|\varphi _s(x)\varphi _s(a)|<\frac{1}{2}sup(\varphi )(a)`$, so $`\varphi _s(x)`$ cannot be $`0`$ thus affirming $`B(a,M)\varphi _s^1(0,\mathrm{})`$. โ
A partition of unity $`\varphi =\{\varphi _s:XR_+\}_{sS}`$ is called slowly oscillating if the corresponding function $`\varphi :Xl_S^1`$ is slowly oscillating.
$`\varphi `$ is called equi-slowly oscillating if the oscillation of all $`\varphi _s`$ is synchronized in the following way: for every $`M>0`$ and every $`ฯต>0`$ there is a bounded subset $`U`$ of $`X`$ such that $`Osc(\varphi _s,M)(x)<ฯต`$ for all $`xXU`$ and all $`sS`$. Obviously, every finite partition of unity into slowly oscillating functions is globally slowly oscillating and is equi-slowly oscillating. Also, every slowly oscillating partition of unity is equi-slowly oscillating.
###### Lemma 7.7.
If $`\varphi =\{\varphi _s:XR_+\}_{sS}`$ is a partition of unity of finite multiplicity $`m`$, then $`\varphi `$ is slowly oscillating if and only if it is equi-slowly oscillating.
###### Proof.
Given $`M,ฯต>0`$ we can find a bounded set $`U`$ such that $`Osc(\varphi _s,M)<ฯต/(2m)`$ for all $`xXU`$ and all $`sS`$. If $`aXU`$ and $`xB(a,M)`$, then the complement $`F`$ of set $`T=\{sS\varphi _s(x)+\varphi _s(a)=0\}`$ contains at most $`2m`$ elements. Since $`|\varphi (x)\varphi )(a)|=_{sF}|\varphi _s(x)\varphi _s(a)|<|F|ฯต/(2m)ฯต`$, $`\varphi `$ is slowly oscillating. โ
###### Lemma 7.8.
If $`\varphi =\{\varphi _s:XR_+\}_{sS}`$ is an equi-slowly oscillating partition of unity of finite multiplicity $`m`$, then its carrier family $`Carr(\varphi )`$ is coarse.
###### Proof.
Notice $`sup(\varphi )1/m`$. Given $`M>0`$ we can find a bounded set $`U`$ such that $`Osc(\varphi _s,M)<1/(2m)`$ for all $`xXU`$ and all $`sS`$. By 7.6, $`L(\varphi |_{(XU)},XU)>M`$ which proves $`Carr(\varphi )`$ is coarse. โ
###### Remark 7.9.
If one drops the assumption of $`\varphi `$ being of finite multiplicity, then the carrier family may not be coarse: Take a cloud $`C_n`$ of $`2^n+1`$ points at location $`2^n`$ with mutual distances equal $`1`$. For each $`xX`$ define $`\varphi _x`$ as taking value $`0`$ at $`x`$ and all points not in its cloud. For points $`yCloud(x)\{x\}`$ we put $`\varphi _x(y)=2^n`$.
###### Corollary 7.10.
If $`๐ฐ=\{U_s\}_{sS}`$ is a cover of $`X`$ of finite multiplicity, then the following conditions are equivalent:
* $`๐ฐ`$ is coarse.
* There is a continuous slowly oscillating partition of unity $`\varphi =\{\varphi _s\}_{sS}`$ on $`XA`$ for some bounded subset $`A`$ of $`X`$ such that $`Carr(\varphi _s)U_s`$ for each $`sS`$.
* There is a slowly oscillating partition of unity $`\varphi =\{\varphi _s\}_{sS}`$ on $`XA`$ for some bounded subset $`A`$ of $`X`$ such that $`Carr(\varphi _s)U_s`$ for each $`sS`$.
###### Proof.
1$``$2. Define $`f(x)=\underset{sS}{}dist(x,XU_s)`$ and $`f_s(x)=dist(x,XU_s)`$. Notice that $`f`$ is a coarsely proper Lipschitz function and 7.5 says that $`\{f_s/f\}_{sS}`$ is an equi-slowly oscillating partition of unity on $`XA`$, where $`A`$ is the zero-set of $`f`$. By 7.7 it is a slowly oscillating partition of unity.
2$``$3 is obvious.
3$``$1 follows from 7.8. โ
## 8. Coarse dimension and Higson corona
Given a metric space $`X`$ by the Higson compactification of $`X`$ we mean a compact Hausdorff space $`h(X)`$ containing $`X`$ as a dense subset with the property that a bounded continuous function $`f:XR_+`$ extends over $`h(X)`$ if and only if $`f`$ is slowly oscillating. If the metric on $`X`$ is proper and $`X`$ is locally compact, then $`X`$ is open in $`h(X)`$ and the remainder $`h(X)X`$ is called the Higson corona of $`X`$ and denoted by $`\nu (X)`$.
A metric space $`X`$ is called $`\delta `$-disjoint for some $`\delta >0`$ if $`d_X(x,y)\delta `$ for all $`xy`$.
###### Theorem 8.1.
If $`X`$ is a $`\delta `$-disjoint metric space for some $`\delta >0`$, then its coarse dimension equals the dimension of the Higson compactification of $`X`$.
###### Proof.
Suppose $`\text{dim}_{rse}^{coa}(X)=m<\mathrm{}`$. Given a finite open cover $`๐ฐ=\{U_s\}_{sS}`$ of the Higson compactification $`h(X)`$ of $`X`$ we find a partition of unity $`f=\{f_s\}_{sS}`$ on $`h(X)`$ such that $`cl(f_s^1(0,1])U_s`$ for each $`sS`$ (see ). As $`f|_X`$ is slowly oscillating (see 7.7), the family $`\{f_s^1(0,1]X\}_{sS}`$ is coarse in $`X`$ (see 7.8). By 7.10 there is a slowly oscillating partition of unity $`g=\{g_s\}_{sS}`$ on $`X`$ whose multiplicity is at most $`m+1`$ and $`g_s^1(0,1]f_s^1(0,1]X`$ for each $`sS`$. Extend each $`g_s`$ over $`h(X)`$ to $`k_s:h(X)[0,1]`$. The resulting family $`k=\{k_s\}_{sS}`$ is a partition of unity on $`h(X)`$. It remains to show $`m(k)m+1`$ and $`k_s^1(0,1]U_s`$ for each $`sS`$. If there is a point $`xh(X)X`$ such that $`k_s(x)>0`$ for all $`sT`$, $`T`$ containing at least $`m+2`$ elements, then the same would be true for some neighborhood $`U_x`$ of $`x`$ in $`h(X)`$. Since $`U_xX\mathrm{}`$ one arrives at a contradiction with the fact that $`m(g)m+1`$. If $`k_s^1(0,1]`$ is not a subset of $`U_s`$ for some $`sS`$, then there is $`xh(X)X`$ so that $`xk_s^1(0,1]cl(f_s^1(0,1])`$. That means there is a neighborhood $`U_x`$ of $`x`$ in $`h(X)`$ on which $`f_s`$ is identically $`0`$. Hence $`g_s|_{(U_xX)}0`$ implying $`k_s(x)=0`$, a contradiction. โ
###### Corollary 8.2.
If $`X`$ is a proper metric space, then the dimension of its Higson corona equals the coarse dimension of $`X`$.
###### Proof.
Consider a maximal $`1`$-disjoint subset $`A`$ of $`X`$. Notice $`\text{dim}_{rse}^{coa}(A)=\text{dim}_{rse}^{coa}(X)`$ and Higson coronas $`\nu (A)`$ and $`\nu (X)`$ for both $`A`$ and $`X`$ are identical. Since $`A`$ is $`1`$-disjoint, $`\text{dim}_{rse}^{coa}(A)=\text{dim}(h(A))=\text{dim}(\nu (A))=\text{dim}(\nu (X))`$. โ
###### Corollary 8.3.
If $`X=AB`$, then the coarse dimension of $`X`$ equals maximum of the coarse dimensions of $`A`$ and $`B`$.
###### Proof.
Let $`m=\mathrm{max}(\text{dim}_{rse}^{coa}(A),\text{dim}_{rse}^{coa}(B))`$. By 5.3, $`\text{dim}_{rse}^{coa}(X)m`$. By switching to maximal $`1`$-disjoint subsets of $`A`$ and $`B`$, respectively, we reduce the general case to that of $`X`$ being $`1`$-disjoint. Consider the Higson compactification $`h(X)`$ of $`X`$. Notice $`cl(A)`$ is the Higson compactification of $`A`$ as any slowly oscillating and bounded function $`f:AR_+`$ extends over $`X`$ to a bounded and slowly oscillating function. The same is true for $`B`$. Since $`h(X)=cl(A)cl(B)`$, $`\text{dim}(h(X))=\mathrm{max}(\text{dim}(cl(A)),\text{dim}(cl(B)))=\mathrm{max}(\text{dim}_{rse}^{coa}(A),\text{dim}_{rse}^{coa}(B))=m`$. โ
We plan to extend 8.3 to other dimensions as well. Our strategy is to show finiteness of the appropriate dimension of $`X`$ first, then use 8.3 as well as the fact that all other dimensions are equal to the coarse dimension of $`X`$ once they are finite (see 6.8 and 6.9).
###### Corollary 8.4.
If $`X=AB`$, then the asymptotic dimension of $`X`$ equals maximum of the asymptotic dimensions of $`A`$ and $`B`$.
###### Proof.
Let $`m=\mathrm{max}(asdim(A),asdim(B))`$. Obviously $`asdim(X)m`$. Given $`M>0`$ find uniformly bounded family $`๐ฐ_A`$ in $`A`$ covering $`A`$ and being the union of $`m+1`$ families, each of them $`3M`$-disjoint. Similarly, find uniformly bounded and $`3M`$-disjoint family $`๐ฐ_B`$ in $`B`$ covering $`B`$ and being the union of $`m+1`$ families, each of them $`3M`$-disjoint. Consider $`๐ฐ=๐ฐ_A๐ฐ_B`$ and let $`๐ฑ=\{B(U,M)\}_{U๐ฐ}`$. Notice $`๐ฑ`$ is uniformly bounded in $`X`$, is of multiplicity at most $`2(m+1)`$, and $`L(๐ฑ,X)M`$. Therefore $`asdim(X)2m+1`$ and (see 6.9) $`asdim(X)=\text{dim}_{rse}^{coa}(X)=m`$. โ
###### Remark 8.5.
8.4 was proved in (see the Finite Union Theorem there) for $`X`$ being a proper metric space by using totally different methods.
###### Corollary 8.6.
If $`X=AB`$, then the major coarse dimension of $`X`$ equals maximum of the major coarse dimensions of $`A`$ and $`B`$.
###### Proof.
Let $`m=\mathrm{max}(\text{dim}_{rse}^{coa}(A),\text{dim}_{rse}^{coa}(B))`$. By 5.3, $`\text{dim}_{RSE}^{COA}(X)m`$. Given a coarse family $`๐ฐ`$ in $`X`$ put $`f(x)=L_๐ฐ(x)`$. If $`f(x)=\mathrm{}`$ for some $`X`$, then $`๐ฐ`$ has a coarse refinement of order at most $`2`$ (see 3.2). Assume $`f(x)<\mathrm{}`$ for all $`xX`$. Pick a coarse refinement $`\{V_a\}_{aA}`$ of multiplicity at most $`m+1`$ of the family $`\{B(a,f(a)/2)\}_{aA}`$. Pick a coarse refinement $`\{V_b\}_{bB}`$ of multiplicity at most $`m+1`$ of the family $`\{B(b,f(b)/2)\}_{bB}`$. If $`V_a\mathrm{}`$ define $`e(V_a)=\{xB(a,f(a))dist(x,V_a)<dist(x,AV_a\}`$. Observe $`\underset{aT}{}e(V_a)\mathrm{}`$ implies $`\underset{aT}{}V_a\mathrm{}`$ for every finite subset $`T`$ of $`S`$. Indeed, suppose $`x\underset{aT}{}e(V_a)`$ and find $`\delta >0`$ such that $`dist(x,V_a)+\delta <dist(x,AV_a\}`$ for all $`aT`$. Pick $`yA`$ so that $`dist(x,A)+\delta >d(x,y)`$. If $`yAV_a`$ for some $`aT`$, then $`dist(x,A)+\delta dist(x,V_a)+\delta <dist(x,AV_a)d(x,y)`$, a contradiction. Therefore the multiplicity of $`\{e(V_a)\}_{aA}`$ is at most $`m+1`$. Do the same procedure for $`B`$ and produce $`\{e(V_b)\}_{bB}`$. If $`xa`$, $`M<f(a)`$ and $`B(x,M)AV_a`$, then $`B(x,M/2)e(V_a)`$. Therefore $`\{e(V_a)\}_{aA}\{e(V_b)\}_{bB}`$ is coarse in $`X`$, refines $`๐ฐ`$, and is of multiplicity at most $`2(m+1)`$. Thus $`\text{dim}_{RSE}^{COA}(X)2m+1`$ and (see 6.8) $`\text{dim}_{RSE}^{COA}(X)=\text{dim}_{rse}^{coa}(X)=m`$. โ
###### Corollary 8.7.
If $`X=AB`$, then the minor asymptotic dimension of $`X`$ equals maximum of the minor asymptotic dimensions of $`A`$ and $`B`$.
###### Proof.
Let $`m=\mathrm{max}(ad(A),ad(B))`$. Obviously $`ad(X)m`$. Suppose $`M>0`$ and find $`N>0`$ such that $`L^m(๐ฐ,A)>2M`$ for all finite covers $`๐ฐ`$ of $`A`$ satisfying $`L(๐ฐ,A)>N`$. We can use the same $`N`$ and claim $`L^m(๐ฐ,B)>2M`$ for all finite covers $`๐ฐ`$ of $`B`$ satisfying $`L(๐ฐ,B)>N`$. Given a finite family $`๐ฐ=\{U_s\}_{sS}`$ in $`X`$ satisfying $`L(๐ฐ,X)>M+N`$, consider $`\{B(U_s,M)\}_{sS}`$ and shrink it on $`A`$ to a family $`\{V_s\}_{sS}`$ of multiplicity at most $`m+1`$ and Lebesque number at least $`2M`$. Do the same for $`B`$ and shrink $`\{B(U_s,M)\}_{sS}`$ on $`B`$ to a family $`\{W_s\}_{sS}`$ of multiplicity at most $`m+1`$ and Lebesque number at least $`2M`$. If $`V_s\mathrm{}`$ define $`e(V_s)=\{xU_sdist(x,V_s)<dist(x,AV_s\}`$. Observe $`\underset{sT}{}e(V_s)\mathrm{}`$ implies $`\underset{sT}{}V_s\mathrm{}`$ for every finite subset $`T`$ of $`S`$ (see the proof of 8.6). Therefore the multiplicity of $`\{e(V_s)\}_{sS}`$ is at most $`m+1`$. Do the same procedure for $`B`$ and produce $`\{e(W_s)\}_{sS}`$. Obviously $`\{e(V_s)\}_{sS}\{e(V_s)\}_{sS}`$ refines $`๐ฐ`$ and is of multiplicity at most $`2(m+1)`$. If we show its Lebesque number is at least $`M`$ we will demonstrate $`ad(X)2m+1`$ and (see 6.9) $`ad(X)=\text{dim}_{rse}^{coa}(X)=m`$. Suppose $`xX`$. Without loss of generality we may assume $`xB`$. There is $`sS`$ such that $`B(x,2M)BW_s`$. Hence $`B(x,M)B(W_s,M)U_s`$ and, since any $`yB(x,M)`$ satisfies $`dist(y,W_s)d(y,x)<M<dist(y,BW_s)`$, we get $`ye(W_s)`$ which completes the proof. โ
## 9. Coarse dimension and absolute extensors
In (Remark 2 on p.1097) Dranishnikov pointed out that $`R_+`$ is not an absolute extensor in the category of proper metric spaces and coarse functions. He characterized proper metric spaces of coarse dimension at most $`n`$ as those for which $`R^{n+1}`$ is an absolute extensor in the category of proper approximately Lipschitz functions (Definition 4 on p.1105 and Theorem 6.6 on p.1111). That still left the door open to the possibility of characterizing coarse dimension via $`R^{n+1}`$ being an absolute extensor in the proper coarse category. The following result clarifies that issue in negative.
###### Theorem 9.1.
For a metric space $`X`$ the following conditions are equivalent:
1. The coarse dimension of $`X`$ is at most $`0`$.
2. $`Y`$ is an absolute extensor of $`X`$ in the proper coarse category for all $`Y`$.
3. $`R_+`$ is an absolute extensor of $`X`$ in the proper coarse category.
###### Proof.
1$``$2. It suffices to show that any unbounded subset $`A`$ of $`X`$ is a coarsely proper and coarse retract of $`X`$. Pick $`x_0X`$. Define by induction on $`n`$ an increasing sequence $`M_n`$ of natural numbers and covers $`๐ฐ^n`$ of $`X`$ satisfying the following properties:
a. $`M_1=1`$.
b. $`๐ฐ^n`$ is $`M_n`$-disjoint, the diameters of its elements are smaller than $`M_{n+1}`$, and $`L(๐ฐ^n,X)>M_n`$.
For each $`U๐ฐ^n`$ so that $`UA\mathrm{}`$, pick $`x_UUA`$ satisfying $`d_X(x_U,x_0)>sup\{d_X(x,x_0)xUA\}1/n`$.
By induction on $`n`$ define a sequence of subsets $`A_n`$ of $`X`$ and a sequence of functions $`r_n:A_nA`$ as follows:
i. $`A_1=A`$ and $`r_1=id_A`$.
ii. $`A_{n+1}`$ is the union of those elements of $`๐ฐ^{n+1}`$ that intersect $`A`$.
iii. If $`xUA_n`$ and $`UA\mathrm{}`$ for some $`U๐ฐ^{n+1}`$, then $`r_{n+1}(x)=x_U`$.
Notice $`X=\underset{n=1}{\overset{\mathrm{}}{}}A_n`$ and let $`r:XA`$ be obtained by pasting all $`r_n`$. Observe that $`xU๐ฐ^k`$ and $`UA\mathrm{}`$ implies $`r(x)U`$. Indeed, for each $`n`$ there is a unique element $`U_x^n๐ฐ^n`$ containing $`x`$ and $`U_x^iU_x^j`$ if $`i<j`$. Find the smallest number $`m`$ so that $`xA_m`$. In that case $`r(x)U_x^m`$ by definition and $`k`$ must be at least $`m`$ so $`U_x^mU_x^k=U`$.
We will show that $`r`$ is coarse by proving $`d_X(x,y)<M_n`$ implies $`d_X(r(x),r(y)M_{n+2}`$. Indeed, if $`d_X(x,y)<M_n`$, then one of the following cases occurs:
Case 1. $`UA_n=\mathrm{}`$, where $`U`$ is the unique element of $`๐ฐ^{n+1}`$ containing both $`x_n`$ and $`y_n`$.
Case 2. $`UA_n\mathrm{}`$, where $`U`$ is the unique element of $`๐ฐ^{n+1}`$ containing both $`x_n`$ and $`y_n`$.
In Case 1 the values $`r(x)`$ and $`r(y)`$ are identical. In Case 2 both $`r(x)`$ and $`r(y)`$ belong to $`UA`$ and the set $`UA`$ is of diameter at most $`M_{n+2}`$, so $`d_X(r(x),r(y)M_{n+2}`$ holds.
If $`r`$ is not coarsely proper, then there is a sequence $`x_n\mathrm{}`$ such that $`r(x_n)`$ is bounded. Obviously, $`x_nA`$ for almost all $`n`$. Consider an element $`U๐ฐ^k`$ containing all of $`r(x_n)`$. The way functions $`r_m`$ were defined implies that there is a sequence of elements $`U_n๐ฐ^{\alpha (n)}`$ with $`\alpha (n)\mathrm{}`$ and all $`U_n`$ containing $`U`$, such that $`U_nA`$ is of almost the same diameter as $`UA`$. That contradicts $`A`$ being unbounded.
2$``$3 is obvious.
3$``$1. Suppose $`\text{dim}_{rse}^{coa}(X)>0`$. By 5.8 there exists a number $`M>0`$ and a coarsely proper sequence $`\{(x_n,y_n)\}_{n=1}^{\mathrm{}}`$ of pairs of points in $`X`$ such that $`dist(x_n,y_n)\mathrm{}`$ and the points $`x_n`$ and $`y_n`$ can be $`M`$-scale connected in $`XB(x_0,n)`$ by long chain of length $`L_n`$ so that $`L_n\mathrm{}`$ as $`n\mathrm{}`$. We may assume $`d_X(x_{n+j},x_n)>n`$ and $`d_X(y_{n+j},y_n)>n`$ for all $`n,j1`$. Let $`B=\{x_n\}\{y_n\}`$. Define $`f:BR_+`$ by sending $`x_n`$ to $`n`$ and $`y_n`$ to $`n+nL_n`$. Notice $`f`$ is coarsely proper and coarse. Suppose $`f`$ extends to a coarse function $`g:XR_+`$. Find $`K>0`$ such that $`d_X(x,y)M`$ implies $`d(f(x),f(y))K`$. Since $`x_n`$ and $`y_n`$ can be connected by a chain of $`L_n`$ points, with consecutive points being separated by at most $`M`$, $`L_nn+nn=d(f(x_n),f(y_n))L_nK`$ which leads to a contradiction for $`n>K`$. โ
## 10. Open problems
In (Problem 1 on p.1126) it is asked if the asymptotic dimension of a proper metric space $`X`$ equals the covering dimension of its Higson corona. Here is our version of that problem.
###### Problem 10.1.
Is there a metric space $`X`$ of infinite asymptotic dimension and finite coarse dimension?
###### Problem 10.2.
Is there a metric space $`X`$ of infinite major coarse dimension and finite coarse dimension?
###### Definition 10.3.
A metric space $`X`$ is of bounded geometry if for every $`M>0`$ there is a uniformly bounded cover $`๐ฐ`$ of $`X`$ of finite multiplicity and the Lebesque number at least $`M`$.
###### Definition 10.4 (,p.1005).
Suppose $`X`$ is a metric space of bounded geometry. Given $`M>0`$ let $`d(M)=m(๐ฐ)1`$, where $`๐ฐ`$ is a uniformly bounded cover $`๐ฐ`$ of minimal multiplicity among those of the Lebesque number at least $`M`$. $`X`$ is of slow dimension growth if $`\underset{M\mathrm{}}{lim}\frac{d(M)}{M}=0`$.
Just as in (Problem 6 on p.1126) one can ask variants of problems 10.1 and 10.2 for spaces of bounded geometry or slow dimension growth.
###### Problem 10.5.
Suppose $`X`$ is of slow dimension growth and finite coarse dimension. Is asymptotic dimension of $`X`$ finite?
###### Problem 10.6.
Suppose $`X`$ is of slow dimension growth and finite coarse dimension. Is the major coarse dimension of $`X`$ finite?
The above problems remain open for minor asymptotic dimension. All of the above problems are of interest in case of $`X`$ being a finitely generated group with word metric, especially $`CAT(0)`$ groups.
###### Problem 10.7.
It is stated in that $`asdim(X\times Y)asdim(X)+asdim(Y)`$. Are the corresponding results true for other dimensions? |
warning/0506/quant-ph0506165.html | ar5iv | text | # The certainty principle
## The quantum angle and the certainty principle
### The notion of the quantum angle.
As it is known, the set of states of any quantum system forms a complex Hilbert space. We will denote it $``$.
Elements of the space $``$ we will denote as $`a,b\mathrm{}`$.
The scalar product in $``$ we will write as $`a|b`$. It is linear with respect to the second argument and anti-linear with respect to the first one.
The norm of a vector $`a`$ we will denote as $`a=a|a^{1/2}`$.
Consider two non-zero vectors $`a,b`$. Let us define between them the quantum angle by the formula:
$$\mathrm{}(a,b)=\mathrm{arccos}\frac{\left|a|b\right|}{ab}.$$
According to Cauchy-Bunyakovsky-Schwarz inequality, under the function $`\mathrm{arccos}`$ we have the value that is not greater than unity. Therefore, the quantum angle is a real number:
$$\mathrm{}(a,b),0\mathrm{}(a,b)\frac{\pi }{2}.$$
For simplification of formulas we will later on always work with normalized vectors: $`a=1`$, $`b=1`$. In this case the formula for the angle is written simply as:
$$\mathrm{}(a,b)=\mathrm{arccos}\left|a|b\right|.$$
### Geometry of quantum angle.
Let us consider the two extreme cases: $`\mathrm{}(a,b)=0`$ and $`\mathrm{}(a,b)=\pi /2`$.
According to the Parseval equality, the first case takes place when the vectors differ only by phase factor:
$$\mathrm{}(a,b)=0ab.$$
From physical point of view, we can say that the corresponding quantum states are identical.
The second case takes place when vectors are orthogonal:
$$\mathrm{}(a,b)=\frac{\pi }{2}ab.$$
In this case we can say that the corresponding quantum states are completely different.
Implying this physical terminology, which is used in the considered extreme cases, let us introduce also the following definition. Let us say that the states described by the vectors $`a`$ and $`b`$ differ not-substantially, if $`\mathrm{}(a,b)<1`$; let us also say that the states differ substantially, if $`\mathrm{}(a,b)1`$.
T h e o r e m. For any three vectors $`a,b,c`$ the inequality takes place (โthe triangle inequalityโ):
$$\mathrm{}(a,c)\mathrm{}(a,b)+\mathrm{}(b,c).$$
In order to prove this theorem, let us first notice that, so far as the inequality is proved for three vectors, we can bound ourself with the case when the Hilbert space $``$ is three dimensional: $`=^3`$.
Multiplying the vectors $`a`$, $`b`$, and $`c`$ by appropriate factors and choosing orthonormal basis in $``$ appropriately, we can achieve that components of these vectors take the form:
$$a=(\begin{array}{ccc}1& 0& 0\end{array}),b=(\begin{array}{ccc}b_1& z& b_3\end{array}),c=(\begin{array}{ccc}c_1& c_2& 0\end{array}),$$
where $`c_1,c_2,b_1,b_3[0;1]`$ are real non-negative numbers, and $`z`$ is complex.
Let us introduce also the auxiliary vector $`b^{}`$:
$$b^{}=(\begin{array}{ccc}b_1& |z|& b_3\end{array}).$$
We have:
$$\mathrm{}(b^{},c)=\mathrm{arccos}\left(b_1c_1+|z|c_2\right)\mathrm{arccos}|b_1c_1+zc_2|=\mathrm{}(b,c).$$
So far as the three vectors $`a`$, $`b^{}`$, and $`c`$ have real coordinates and unit lengths, the triangle inequality for them is the well known triangle inequality on the sphere in the real three-dimensional Euclidean space:
$$\mathrm{}(a,c)\mathrm{}(a,b^{})+\mathrm{}(b^{},c).$$
Combining the obtained inequalities we get:
$$\mathrm{}(a,c)\mathrm{}(a,b^{})+\mathrm{}(b^{},c)=\mathrm{}(a,b)+\mathrm{}(b^{},c)\mathrm{}(a,b)+\mathrm{}(b,c)\mathrm{}$$
Summarizing what was said above, we can say that the quantum angle $`\mathrm{}`$, considered as a function on the set of pairs of physical quantum states (considered to phase factor), is a metric.
T h e o r e m. The metric space of physical quantum states with the metric $`\mathrm{}`$ is complete.
The proof of this theorem is not difficult, but it requires substantial technical work. So we omit it here.
### Kinematics of quantum angle.
Let now the vector $`r`$ depend on the real parameter $`t`$: $`t`$, $`r(t)`$, $`r(t)=1`$.
Let us define the quantum velocity $`v(t)`$ by the formula:
$$v(t)=\dot{r}(t)=\underset{\delta t0}{lim}\frac{r(t+\delta t)r(t)}{\delta t}.$$
Let us define also the quantum angular speed $`\omega (t)`$ by the formula:
$$\omega (t)=\underset{\delta t0}{lim}\frac{\mathrm{}(r(t+\delta t),r(t))}{|\delta t|}.$$
In order to express $`\omega (t)`$ through $`v(t)`$ let us decompose $`v(t)`$ into the two orthogonal components:
$$v_{}(t)=r(t)r(t)|v(t),v_{}(t)=v(t)v_{}(t).$$
T h e o r e m. The quantum angular speed is equal to the norm of the orthogonal component of the quantum velocity:
$$\omega (t)=v_{}(t).$$
In order to prove this theorem let us use the Parseval equality to change $`\mathrm{arccos}`$ to $`\mathrm{arcsin}`$:
$$\mathrm{}(r(t+\delta t),r(t))=\mathrm{arccos}\left|r(t+\delta t)|r(t)\right|=$$
$$=\mathrm{arcsin}r(t+\delta t)r(t)r(t)|r(t+\delta t)=$$
$$=\mathrm{arcsin}r(t+\delta t)r(t)r(t)r(t)|r(t+\delta t)r(t)=$$
$$=\mathrm{arcsin}v(t)\delta t+o(\delta t)r(t)r(t)|v(t)\delta t+o(\delta t)=$$
$$=\mathrm{arcsin}\left(v(t)r(t)r(t)|v(t)\right)\delta t+o(\delta t)=$$
$$=\mathrm{arcsin}v_{}(t)\delta t+o(\delta t)=\mathrm{arcsin}\left(v_{}(t)|\delta t|+o(\delta t)\right)=$$
$$=v_{}(t)|\delta t|+o(\delta t)\mathrm{}$$
T h e o r e m. The quantum angle satisfy the estimate:
$$\mathrm{}(r(t_2),r(t_1))\left|_{t_1}^{t_2}\omega (t)๐t\right|.$$
(1)
For the proof let us use the triangle inequality:
$$\left|\mathrm{}(r(t+\delta t),r(t_1))\mathrm{}(r(t),r(t_1))\right|\mathrm{}(r(t+\delta t),r(t)).$$
Dividing by $`|\delta t|`$ and taking the limit $`\delta t0`$, we get:
$$\left|\frac{d}{dt}\mathrm{}(r(t),r(t_1))\right|\omega (t).$$
Performing integration from $`t_1`$ to $`t_2`$, we get the desired inequality. $`\mathrm{}`$
In fact, the estimate (1) is the best. Namely, there is the following
T h e o r e m. The quantum angle between two vectors $`r_1`$ and $`r_2`$ can be expressed by the formula:
$$\mathrm{}(r_2,r_1)=\mathrm{min}\left|_{t_1}^{t_2}\omega (t)๐t\right|,$$
where the minimum is taken among all curves $`r(t)`$ with ends in $`r_1`$ and $`r_2`$: $`r(t_1)=r_1`$, $`r(t_2)=r_2`$.
For the proof of the theorem let us twist the phase of $`r_2`$, $`r_2r_2^{}=e^{i\alpha }r_2`$, so that $`r_2^{}|r_1`$ become real and $`r_2^{}|r_1[0,1]`$.
Let us consider the linear shell of $`r_2^{}`$ and $`r_1`$, $`(r_2^{},r_1)`$. There we can choose an orthonormal basis so that
$$r_1=(\begin{array}{ccc}1& 0& \end{array}),r_2^{}=(\begin{array}{ccc}a& b& \end{array}),$$
where $`a,b[0,1]`$ are real non-negative numbers.
Considering then $`r_2^{}`$ and $`r_1`$ as real vectors on Euclidean plane we see that it is possible to stretch the circular arc between them where the estimation integral is exactly equal to the quantum angle. $`\mathrm{}`$
### Dynamics of quantum angle.
Let us have now a strongly continuous one-parameter unitary group $`U(\delta s)=e^{i\delta sA/\mathrm{}}`$, where $`A=A^{}`$ is a self-adjoint operator in the space of states $``$ (it is called the infinitesimal generator of $`U(\delta s)`$); $`\delta s`$ is the parameter of the group; $`\mathrm{}`$ is the Planckโs constant.
And suppose now that the dependence of state vector on the parameter $`\delta s`$ is defined by the formula:
$$r(\delta s)=|\delta s=U(\delta s)=e^{i\delta sA/\mathrm{}}.$$
Here $`|\delta s`$ is another notation for the state vector connected with parameter equal to $`\delta s`$; $``$ is a fixed ket-vector of state.
Let us suppose that the function $`r(\delta s)`$ is differentiable. Then the quantum velocity is expressed by the formula:
$$v(\delta s)=\frac{1}{i\mathrm{}}Ae^{i\delta sA/\mathrm{}}=\frac{1}{i\mathrm{}}A|\delta s.$$
The mean of the operator $`A`$ does not depend on time:
$$\overline{A}=\delta s|A|\delta s=e^{+i\delta sA/\mathrm{}}Ae^{i\delta sA/\mathrm{}}=Ae^{+i\delta sA/\mathrm{}}e^{i\delta sA/\mathrm{}}=A.$$
Therefore the components of the quantum velocity can be written just as:
$$v_{}(\delta s)=|\delta s\delta s|\frac{1}{i\mathrm{}}A|\delta s=\frac{1}{i\mathrm{}}\overline{A}|\delta s$$
$$v_{}(\delta s)=\frac{1}{i\mathrm{}}\left(A\overline{A}\right)|\delta s.$$
The quantum angular speed turns out to be independent of time also:
$$\omega (\delta s)=v_{}(\delta s)=\frac{1}{\mathrm{}}\delta s|\left(A\overline{A}\right)^2|\delta s^{1/2}=\frac{1}{\mathrm{}}e^{+i\delta sA/\mathrm{}}\left(A\overline{A}\right)^2e^{i\delta sA/\mathrm{}}^{1/2}=$$
$$=\frac{1}{\mathrm{}}\left(A\overline{A}\right)^2e^{+i\delta sA/\mathrm{}}e^{i\delta sA/\mathrm{}}^{1/2}=\frac{1}{\mathrm{}}\left(A\overline{A}\right)^2^{1/2}=\frac{1}{\mathrm{}}\mathrm{\Delta }_{}A.$$
Here $`\mathrm{\Delta }_{}A`$ is a short notation for the standard deviation of $`A`$ in the state $``$.
Consider now how the quantum angle $`\mathrm{}(|\delta s,)`$ behaves in this case. Using for it the estimate (1) we have:
$$\mathrm{}(|\delta s,)|_0^{\delta s}\omega (\sigma )d\sigma |=\frac{1}{\mathrm{}}|\delta s|\mathrm{\Delta }_{}A.$$
From this inequality we obtain the
T h e o r e m. So that under the action of strongly continuous one-parameter unitary group $`U(\delta s)=e^{i\delta sA/\mathrm{}}`$ the initial state vector $``$ changes substantially, it is necessary to satisfy the inequality:
$$\begin{array}{ccc}& & \\ & & \\ & |\delta s|\mathrm{\Delta }_{}A\mathrm{}& \\ & & \end{array}$$
(2)
By the example of the Schrรถdinger particle we will see that this theorem turns out to be closely connected with the Heisenberg uncertainty principle, but has other meaning. Taking into account this connection, we can name this theorem the certainty principle.
The inequality expressing the certainty principle can be written also in the following way:
$$\begin{array}{ccc}& & \\ & & \\ & \mathrm{\Delta }_{}(\delta sA)\mathrm{}& \\ & & \end{array}$$
(3)
In this form it can be naturally carried over to the case when $`\delta s`$ and $`A`$ are matrices.
## Examples
### One-dimensional Schrรถdinger particle.
Consider the one-dimensional Schrรถdinger particle with the coordinate defined by the variable $`x`$. Its state vector can be written by the wave function $`\psi (x)`$. In its space of states the group of shifts acts by the formula:
$$U(\delta x)\psi (x)=\psi (x\delta x).$$
This group can be written in the form:
$$U(\delta x)=e^{i\delta xP/\mathrm{}},P=i\mathrm{}\frac{d}{dx}.$$
Here $`P`$ is the operator of momentum.
Applying the certainty principle in the form (2), we get:
$$|\delta x|\mathrm{\Delta }_{\psi (x)}P\mathrm{}.$$
If we take as $`\psi (x)`$ a well localized packet of de Broglie waves, that turns into zero outside of some interval $`l`$, then from this inequality, in particular, we have that $`l\mathrm{}/\mathrm{\Delta }_{\psi (x)}P`$: because when the packet is moved to the distance $`l`$, the change of the quantum angle must turn out to be greater than $`1`$ (namely, $`\pi /2`$).
So, the Heisenberg uncertainty principle, if it is understood in qualitative sense, follows from the certainty principle.
But if we understand the Heisenberg uncertainty principle in quantitative sense, according to the Pauli-Weyl inequality
$$\mathrm{\Delta }_{\psi (x)}X\mathrm{\Delta }_{\psi (x)}P\frac{\mathrm{}}{2},$$
(4)
then there is no direct connection between these two principles.
Furthermore, from physical point of view, the Heisenberg uncertainty principle and the certainty principle are like two points of view on the spread of the wave packet. The Heisenberg uncertainty principle says, that the wave packet is spread, because the classical state of the particle is badly defined. On the other hand, the certainty principle states, that the wave packet is spread, because the quantum state is well defined.
To the three-dimensional case the certainty principle is easily generalized in the form (3):
$$\mathrm{\Delta }_{\psi (x)}(\delta x_iP_i)\mathrm{},$$
where summation over $`i`$ is implied.
### Schrรถdinger particle on circle.
Consider now a plane with Cartesian coordinates $`x`$ and $`y`$. Let us have the circle $`x^2+y^2=1`$ defined on the plane. On this circle as one-dimensional coordinate we can use the polar angle $`\phi `$, considered to $`2\pi `$.
The state vector of the Schrรถdinger particle on this circle can be written by the wave function $`\psi (\phi )`$. And $`\psi (\phi +2\pi )=\psi (\phi )`$.
In the space of states the action of the rotation group is naturally defined:
$$U(\delta \phi )\psi (\phi )=\psi (\phi \delta \phi ).$$
This group can be written in the form:
$$U(\delta \phi )=e^{i\delta \phi J/\mathrm{}},J=i\mathrm{}\frac{d}{d\phi }.$$
Here $`J`$ is the operator of angular momentum.
Using the certainty principle does not arise any difficulties:
$$|\delta \phi |\mathrm{\Delta }_{\psi (\phi )}J\mathrm{}.$$
As regards the uncertainty principle, a carrying over of the inequality (4) to this case is impossible<sup>2</sup><sup>2</sup>2About some attempts that were made to suggest an inequality like (4), see ..
In the three-dimensional case the certainty principle also easily gives:
$$\mathrm{\Delta }_{\psi (x,y,z)}(\delta \phi _iJ_i)\mathrm{}.$$
### A system with Hamiltonian independent of time.
Let us have now a quantum system with a Hamiltonian $`H`$ independent of time. On the space of states we have the following action of the group of time shifts:
$$U(\delta t)=e^{i\delta t(H)/\mathrm{}}.$$
The certainty principle in this case immediately gives:
$$|\delta t|\mathrm{\Delta }_{}H\mathrm{}.$$
If we apply this formula, for example, to estimation of the life time of a quasi-stationary decaying state, then it states that its typical life time is not less than Planckโs constant divided by the width of the corresponding energy level. And here all the terms can be defined with exact mathematical sense.
As regards to attempts to formulate the Heisenberg uncertainty principle for values time - energy, after Bohr has declared such a principle (in qualitative sense) so many researches and discussions were performed for clarification of its sense, that it is possible to write about them a separate review. As far as I know, a rigorous formulation of the uncertainty principle for this case have not been formulated till now<sup>3</sup><sup>3</sup>3See also the discussion of this question by J. Baez ..
### Relativistic systems.
Consider now any relativistic quantum system. On its space of states the Poincare group acts. As an example of such a system any RCQ-quantized field can serve<sup>4</sup><sup>4</sup>4All other examples that I know either come from this or are mathematically insolvent. . And let us restrict ourself to the discussion of the case when the field turns out quantized in the usual Hilbert space.
In this case the application of the certainty principle does not meet any difficulties:
$$\mathrm{\Delta }_{}(\delta x_\mu P_\mu +\frac{1}{2}\delta \omega _{\mu \nu }J_{\mu \nu })\mathrm{},$$
where $`P_\mu `$ is the vector operator of energy-momentum, $`J_{\mu \nu }`$ is the tensor operator of four-dimensional angular momentum, $`\delta x_\mu `$ and $`\delta \omega _{\nu \rho }`$ are the standard logarithmic coordinates of the Poincare group.
As regards to the application of the Heisenberg uncertainty principle, it is unlikely to be possible. In the previous paragraph we have seen that for the values time - energy it arose great difficulties.
In this connection, even from the ideas relativistic invariance it is clear that even for coordinates and momenta the situation cannot be simple. And it turns out to be true, because it is known that all attempts to introduce the notion of coordinates (as self-adjoint operators on the space of states) for relativistic systems are quite artificial<sup>5</sup><sup>5</sup>5And it turns out that for some systems introduction of good operators of coordinates is not possible at all .. |
warning/0506/cond-mat0506201.html | ar5iv | text | # Magnetic inversion symmetry breaking and ferroelectricity in ๐๐๐๐ง๐_๐
## Abstract
$`\mathrm{TbMnO}_3`$ is an orthorhombic insulator where incommensurate spin order for temperature $`T_N<41\mathrm{K}`$ is accompanied by ferroelectric order for $`T<28\mathrm{K}`$. To understand this, we establish the magnetic structure above and below the ferroelectric transition using neutron diffraction. In the paraelectric phase, the spin structure is incommensurate and longitudinally-modulated. In the ferroelectric phase, however, there is a transverse incommensurate spiral. We show that the spiral breaks spatial inversion symmetry and can account for magnetoelectricity in $`\mathrm{TbMnO}_3`$.
The coexistence of antiferromagnetism and ferroelectricity in solid materials is rare, and much rarer still is a strong coupling between these two order parameters Smolenskii ; Schmid . In non-magnetic perovskites like $`\mathrm{BaTiO}_3`$, ferroelectricity is driven by a hybridization of empty d orbitals with occupied p orbitals of the octahedrally coordinated oxygen ions Cohen . This mechanism requires empty d orbitals and thus cannot lead to magnetic ferroelectric materials. In materials such as $`\mathrm{BiMnO}_3`$, lone $`s^2`$ electron pairs can lower their energy by hybridizing with empty p orbitals Atanasov . While it leads to coexistence of magnetic order with electric polarization at low temperatures, the very different ordering temperatures show that the two order parameter lower the symmetry of the systems through distinctly different collective effects.
$`\mathrm{TbMnO}_3`$ is an antiferromagnet which contains $`\mathrm{Mn}^{3+}`$ ions with occupied d orbitals and no lone s<sup>2</sup> cation Kimura . So, as for a number of recently discovered multiferroics Lottermoser ; Hur ; Kobayashi , neither of the mechanisms described above can explain the coexistence of magnetic and electric order. In these materials, a magnetic field of a few Tesla can switch the direction of the electric polarization Kimura \- proof of a strong direct coupling between the magnetic and electric polarization. Rare earth manganese oxides show a plethora of exciting phenomena which arise from competing interactions. $`\mathrm{LaMnO}_3`$ is the parent compound to a series of materials featuring colossal magnetoresistance, and shows orbital ordering of its $`e_g`$ orbital, giving rise to layered antiferromagnetic ordering. With decreasing rare-earth size, there is a tendency towards incommensurate magnetic order. Ferroelectric ground states were recently predicted for doped manganite of the type $`\mathrm{R}_{1\mathrm{x}}\mathrm{Ca}_\mathrm{x}\mathrm{MnO}_3`$, but they have not yet been observed Efremov . The coexistence and strong coupling of ferroelectricity and antiferromagnetism in $`\mathrm{TbMnO}_3`$ suggests the presence of a non-conventional coupling mechanism involving competing spin interactions, and it is thus of both practical and fundamental importance.
To develop a microscopic theory of the new coupling mechanism, unambiguous determination of the symmetry of the magnetic order is essential. In this Letter, we present neutron diffraction measurements of orthorhombic $`\mathrm{TbMnO}_3`$, which determine the magnetic ground states and the phase diagram as a function of temperature and a magnetic field $`๐||๐`$. We find that the ferroelectric phase transition coincides with a magnetic transition from a longitudinal incommensurate structure to an incommensurate spiral structure that breaks spatial inversion symmetry. We show that a recent theory proposed for axial non-axial parity breaking Harris ; LawesRapid predicts the observed orientation of the ferroelectric polarization based on the symmetry of the magnetic structure that we report.
$`\mathrm{TbMnO}_3`$ crystals were grown using an optical floating zone furnace. A measurement of the temperature dependence of the dielectric constant of the two samples which we studied confirmed ferroelectricity below $`T`$=$`26\mathrm{K}`$ and $`T`$=$`28\mathrm{K}`$ for sample 1 and 2, respectively. The discrepancy may result from slightly differing oxygen partial pressure during annealing. The space group is #62, and in the Pbnm setting the lattice parameters are $`a`$=$`5.3\AA `$, $`b`$=$`5.86\AA `$ and $`c`$=$`7.49\AA `$. The measurements were performed on two single-crystals weighing $`40\mathrm{mg}`$ and $`220\mathrm{mg}`$, respectively, with the BT2 and SPINS spectrometers at NIST, and the TriCS 4-circle diffractometer at PSI.
The wave-vector dependence of the diffraction intensity along the $`(0,k,1)`$ direction of reciprocal space (Figs. 1 and 2) illustrates the temperature dependence of magnetic order in $`\mathrm{TbMnO}_3`$. $`\mathrm{Mn}^{3+}`$ spins develop long-range order at $`T_N`$=$`41\mathrm{K}`$ Quezel ; Kajimoto . Immediately below the transition temperature, magnetic Bragg peaks are observed at $`(0,1q,1)`$ and $`(0,q,1)`$ positions. At $`35\mathrm{K}`$ the magnetic order is described by a single Fourier component associated with a wave-vector $`(0,q,0)`$ with $`q=0.27`$, and no apparent higher order Fourier components. Upon cooling below $`T`$=$`28\mathrm{K}`$, third-order magnetic Bragg peaks $`(0,13q,0)`$ appear (Fig. 1), indicating step-like modulation of the magnetic moment. The inset to Fig. 2b shows that the wave vector continues to vary with temperature throughout the ferroelectric phase.
Below $`T`$=$`7\mathrm{K}`$, additional magnetic peaks appear that are associated with Tb moments and indicate that their interactions favor a different ordering wave-vector than for the Mn moments. At $`T`$=$`2\mathrm{K}`$, there is a strong peak at $`(0,t,1)`$ with $`t=0.425`$, and several higher-order peaks as indicated in Fig. 1. The many strong odd high-order reflections indicate that the order associated with the $`(0,t,0)`$ wave-vector is strongly distorted (bunched structure) as expected for anisotropic rare-earth magnets. The correlation length and the incommensuration of the Tb order at low-temperatures is sample dependent, as is often found in systems with phase transitions between incommensurate structures. The correlation length along the b-axis was $`58(20)`$ and $`280(20)\AA `$, and $`t`$ was $`0.41`$ and $`0.425`$ in sample $`1`$ and $`2`$, respectively.
The magnetic structures at $`T`$=$`15\mathrm{K}`$, and $`35\mathrm{K}`$ were determined from up to 922 first-order magnetic Bragg peaks each. Only first order peaks were included in the refinement as these are sufficient to determine the symmetry of the magnetic structure. Representational analysis was used to find the irreducible representations that describe the magnetic structures. The structure at $`T`$=$`35\mathrm{K}`$ in the high-temperature incommensurate (HTI) phase can be described by a single irreducible representation. The best fit with $`\chi ^2=0.86`$ was obtained for a structure described by representation $`\mathrm{\Gamma }_3`$, and the excellent agreement between the model and the data is shown in Fig. 3a). The magnetic structure is described by $`๐ฆ_3^{\mathrm{Mn}}=(0.0(8),2.90(5),0.0(5))\mu _B`$ and $`๐ฆ_3^{\mathrm{Tb}}=(0,0,0.0(4))\mu _B`$ where the subscript denotes the irreducible representation indicated in Fig. 4c. The magnetic structure is longitudinally-modulated with moment along the b-axis, as illustrated in Fig. 4a and consistent with an earlier study Quezel . The absence of observable higher order peaks indicates that the magnetic structure at $`T`$=$`35\mathrm{K}`$ is sinusoidally modulated.
Two irreducible representations are required to describe the magnetic structure at $`T`$=$`15\mathrm{K}`$ in the low-temperature incommensurate (LTI) phase. We found best agreement with $`\chi ^2=2.19`$ for magnetic ordering involving $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_3`$, as shown in Fig. 3b. Fits using the $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_3`$, or the $`\mathrm{\Gamma }_3`$ and $`\mathrm{\Gamma }_4`$ representation pairs led to $`\chi ^2=14.5`$ and higher, and can thus be excluded. Neglecting higher order reflections, the magnetic structure is given by $`๐ฆ_3^{\mathrm{Mn}}=(0.0(5),3.9(1),0.0(7))\mu _B`$, $`๐ฆ_3^{\mathrm{Tb}}=(0,0,0(1))\mu _B`$, $`๐ฆ_2^{\mathrm{Mn}}=(0.0(1),0.0(8),2.8(1))\mu _B`$ and $`๐ฆ_2^{\mathrm{Tb}}=(1.2(1),0(1),0)\mu _B`$. The experiment was not sensitive to the phase between the $`y`$ and $`z`$-component of the Mn moment. From the size of the moment, however, we deduce that the Mn moments form an elliptical spiral. The data did not favor a phase difference between the Tb and Mn moments, so these phases remain undetermined. Symmetry splits the Tb moments into two orbits which representation theory normally treats as independent. However, as suggested by Landau theory Harris , we took these two Tb amplitudes to be identical. The phase between the two orbits was found to be to $`1.3(3)\pi `$. The greatly improved fit is evidence that the Tb sublattice carries significant magnetization in the LTI phase, presumably as a consequence of the exchange field from the ordered Mn sublattice.
Fig. 5(a-b) shows the field dependence of the $`(0,q,1)`$ magnetic Bragg reflection, which arises from the Mn spin spiral. Both the position and the intensity are field independent to within errorbar - evidence that the structure remains a spiral up to at least $`6\mathrm{T}`$. Our calculations show that the intensity should drop by $`7\%`$ if the $`z`$-component of $`\mathrm{\Gamma }_2`$ were extinguished. In contrast, no decrease is observed to within an error bar of $`2\%`$ between $`0`$ and $`6\mathrm{T}`$.
The $`(0,1q,1)`$ Bragg reflection shown in Fig. 5a arises from Mn $`\mathrm{\Gamma }_3`$ magnetization along the a-axis and from Tb $`\mathrm{\Gamma }_2`$ magnetization along the a-axis. Because the $`x`$-component of the Mn moment is small, the $`(0,1q,1)`$ Bragg reflection is particularly sensitive to Tb order. For $`T<28\mathrm{K}`$ the $`(0,1q,1)`$ intensity is suppressed by a field $`๐||๐`$ confirming that the modulated Tb moment is oriented along that direction. Below the Tb ordering temperature, the field dependent magnetic Bragg intensity has a finite-field maximum (Fig. 5a), indicating a spin-flop transition.
We collected $`51`$ magnetic Bragg peaks at $`T`$=$`4\mathrm{K}`$ and $`H`$=$`4\mathrm{T}`$ along $`๐`$ to determine the magnetic structure at low temperatures above the critical field for $`(0,t,0)`$ Tb order (Fig. 5c). The magnetic structure can be described by $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_3`$ with $`\chi ^2=3.81`$ or by $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_3`$ with $`\chi ^2=4.19`$. Since a field along the a-direction disfavors antiparallel spin alignment in the same direction as in the $`\mathrm{\Gamma }_1`$-$`\mathrm{\Gamma }_3`$ structure, we infer that the structure is given by $`๐ฆ_3^{\mathrm{Mn}}=(0.3(4),4.7(3),0.0(5))\mu _B`$, $`๐ฆ_3^{\mathrm{Tb}}=(0,0,0.0(3))\mu _B`$, $`๐ฆ_2^{\mathrm{Mn}}=(0.0(2),0.0(4),3.0(3))\mu _B`$ and $`๐ฆ_2^{\mathrm{Tb}}=(0.3(2),0.0(4),0)\mu _B`$. This result suggests that the spin spiral structure is more stable for fields along the a-axis than for fields along the b-axis Kimura .
Harris et al. Harris ; LawesRapid recently showed that insulators with axial-non-axial parity breaking magnetic phase transitions must also be electrically polarized. Given the magnetic structure and the temperature dependence of the magnetic order parameter, the theory predicts the direction and temperature dependence of the electric polarization resulting from a symmetry allowed trilinear coupling term. In the following we show that this theory correctly accounts for the direction of the electric polarization in the LTI phase of $`\mathrm{TbMnO}_3`$, and the absence of electric polarization in the HTI phase. The trilinear magnetoelectric coupling term in the Landau free energy expansion is written as $`V=_{\mathrm{uv}\gamma }a_{\mathrm{uv}\gamma }\sigma _\mathrm{u}(k)\sigma _\mathrm{v}(k)P_\gamma `$. Here $`\sigma _\mathrm{u}(k)`$ is the magnetic order parameter of irreducible representation $`\mathrm{\Gamma }_u`$, $`P_\gamma `$ is the electric polarization along the the $`\gamma `$ crystallographic direction and $`a_{\mathrm{uv}\gamma }`$ parametrizes the strength of the interaction between the electric and magnetic order parameters. In the HTI phase, the magnetic order is described by only one irreducible representation, $`\mathrm{\Gamma }_3`$ and due to the high symmetry of the Mn moments, it is possible to define $`\sigma _u`$ such that under inversion $`\sigma _\mathrm{u}(k)\sigma _\mathrm{u}(k)^{}`$. The trilinear coupling thus consists only of terms such as $`V=_\gamma a_\gamma |\sigma _3(k)|^2P_\gamma `$. Since interactions in a Landau expansion must have the symmetry of the paramagnetic phase, this interaction must be invariant under inversion. This requires that $`a_\gamma `$ vanishes, so that $`P_\gamma =0`$ in the HTI phase. The conclusion remains valid for Tb order with no restriction on the phase between its two orbits.
For the LTI phase, which is described by $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_3`$, there are, however, additional terms such as $`V=_\gamma b_\gamma \sigma _2(k)\sigma _3(k)P_\gamma +c.c.`$, where c. c. denotes the complex conjugate. For $`V`$ to be an invariant, $`P_\gamma `$ must transform as the product of $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_3`$. That is, the electric polarization must be even under $`1`$ and $`m_{\mathrm{yz}}`$, and odd under $`2_\mathrm{y}`$ and $`m_{\mathrm{xy}}`$. This condition can only be satisfied for an electric polarization along the c-axis. Previous dielectric experiments have shown that the electric polarization that develops below the transition to the LTI phase is indeed oriented along the c-axis.
The coupling of the magnetic order parameter to electric polarization in $`\mathrm{TbMnO}_3`$ Kimura is similar to that in $`\mathrm{Ni}_3\mathrm{V}_2\mathrm{O}_8`$, which adopts two different incommensurate magnetic structures LawesPRL , one of them ferroelectric and described by two irreducible representations. The correct prediction of the electric polarization for both $`\mathrm{TbMnO}_3`$ and $`\mathrm{Ni}_3\mathrm{V}_2\mathrm{O}_8`$ (Ref. LawesRapid, ) suggests that magnetoelectricity resulting from a trilinear magnetoelectric coupling term may be commonplace in insulating transition metal oxides with non-collinear incommensurate structures. Accordingly we find less compelling the suggestion Kimura that the appearance of ferroelectricity is associated with an incommensurate to commensurate phase transition. Indeed, there is no evidence of a lock-in transition (see Fig. 2b) and it therefor appears to be irrelevant from the point of view of ferroelectricity whether the modulated magnetic order is commensurate or truly incommensurate.
In summary, we have determined the magnetic structure of the paraelectric and ferroelectric phases of $`\mathrm{TbMnO}_3`$. We showed that the paraelectric, magnetically incommensurate phase has sinusoidally-modulated collinear magnetic order that does not break inversion symmetry. The ferroelectric phase, however, has non-collinear incommensurate magnetic order described by two irreducible representations, which explicitly breaks inversion symmetry and thus gives rise to electric polarization. The qualitative aspects of magnetoelectric effects in $`\mathrm{TbMnO}_3`$ appear to be accounted for by a trilinear coupling term in a Landau free energy expansion as proposed by Harris et al. Harris ; LawesRapid . Understanding the magnitude of the effect will require experimental as well as theoretical work to track lattice, charge, and orbital degrees of freedom through axial-non-axial parity breaking phase transition in insulators such as $`\mathrm{TbMnO}_3`$. Apart from the fundamental challenge, improved understanding of magnetoelectricity in these systems may help to produce materials for room-temperature applications.
###### Acknowledgements.
We thank A. Aharony, O. Entin-Wohlman, and A. Ramirez for helpful discussions. This work was supported by the Swiss National Science Foundation under Contract No. PP002-102831. Work at Johns Hopkins University was supported by the DoE through DE-FG02-02ER45983. Work at University of Pennsylvania was supported by the U.S.-Israel Binational Science Foundation under Grant number 2000073. Work at Rutgers University was supported by the NSF-DMR-MRSEC-00-080008. This work is based on experiments performed at the Swiss spallation neutron source SINQ, Paul Scherrer Institute, Villigen, Switzerland. The work at SPINS is based upon activities supported by the National Science Foundation under Agreement No. DMR-9986442. |
warning/0506/hep-th0506024.html | ar5iv | text | # A Model with Interacting Composites
## 1 Introduction
Fermions are an essential ingredient in nature. It is an ever repeating idea to build a model of nature using only fermions, where all the observed bosons are constructed as composites of these entities. In solid state physics electrons, fermions in character, come together to form bosons . Heisenberg spent years to formulate a โtheory of everythingโ using only fermions . Another attempt in this direction came with the work of Gรผrsey, , where a non-polynomial Lagrangian was written to describe self-interacting fermions. Kortel found solutions to this theory which were shown to be instantonic and meronic solutions much later .
One of us, with collaborators, tried to make quantum sense of this model a while ago \[ref6-ref8\], finding that even if these attempts are justified, this model went to a trivial model as the cut-off is removed. We calculated several processes involving incoming and outgoing spinors which gave exactly the naive quark model results, missing the observed logarithmic behaviour predicted by QCD calculations.
During the last twenty years, many papers were written on making sense out of โtrivial modelsโ, interpreting them as effective theories without taking the cutoff to infinity. One of these models is the Nambu-Jona-Lasinio model . Although this model is shown to be a trivial in four dimensions \[ref11-ref13\], since the coupling constant goes to zero with a negative power of the logarithm of the ultraviolet cut-off, as an effective model in low energies it gives us important insight to several processes. In QCD, the studies of hadron mass generation through spontaneous symmetry breaking, important clues to results of the nuclear pairing interaction and the approximate validity of the interacting boson model can be cited as some examples.
There were also attempts to couple the Nambu-Jona-Lasinio model to a gauge field, the so called gauged Nambu-Jona-Lasinio model to be able to get a non-trivial field theory. It was shown that if one has sufficient number of fermion flavors, such a construction is indeed possible . Actually the Nambu-Jona-Lasinio model was constructed based on an analogy with the BCS theory of superconductivity, where fermions come together to form the bosonic interaction necessary to explain the physical phenomena.
Here we want to give a new interpretation of our old work. First we see that the polynomial form of the original model really does not correspond to it in the exact sense. The two versions have obey different symmetries. Then we go to higher orders in our calculation in the new version, beyond the one loop for the scattering processes. It is shown that by using the Dyson-Schwinger and Bethe-Salpeter equations some of the fundamental processes can be better understood. We see that while the non-trivial scattering of the fundamental fields is not allowed, bound states can scatter from each other with non-trivial amplitudes. This phenomena is another example of treating the bound states, instead of the principal fields, as physical entities, that go through physical processes such as scattering.
In our model we need an infinite renormalization in one of the diagrams. Further renormalization is necessary at each higher loop, like any other renormalizable model. The difference between our model and other renormalizable models lies in the fact that, although our model is a renormalizable one using naive dimensional counting arguments, we have only one set of diagrams which is divergent. We need to renormalize only one of the coupling constants by an infinite amount. This set of diagrams, corresponding to the scattering of two bound states to two bound states, have the same type of divergence, i.e. $`\frac{1}{ฯต}`$ in the dimensional regularization scheme for all odd number of loops. The contributions from even number of diagrams are finite, hence require no infinite renormalization. The scattering of two scalars to four, or to any higher even number of scalars is finite, as expected to have a renormalizable model, whereas production of spinors from the scattering of scalars go to zero as the cut-off is removed.
We will outline the model as is given in Refs. and in Section I and give our new results in subsequent sections .
## 2 The Model
We start with the Lagrangian of the model given as
$$L=i\overline{\psi }/\psi +g\overline{\psi }\psi \varphi +\lambda (g\overline{\psi }\psi a\varphi ^3).$$
(1)
Here the only terms with kinetic part are the spinors. Here $`\lambda `$ is a Lagrange multiplier field, $`\varphi `$ is a scalar field with no kinetic part, $`g`$ and $`a`$ are coupling constants. This expression contains two constraint equations, obtained from writing the Euler -Lagrange equations for the $`\lambda `$ and $`\varphi `$ fields.
$$g\overline{\psi }\psi 3\lambda a\varphi ^2=0,$$
(2)
and
$$g\overline{\psi }\psi a\varphi ^3=0.$$
(3)
The Lagrangian given above is just an attempt in writing the original Gรผrsey Lagrangian
$$L=i\overline{\psi }/\psi +g^{}(\overline{\psi }\psi )^{4/3},$$
(4)
in a polynomial form.
We see that the $`\gamma ^5`$ invariance of the original Lagrangian is retained in the form written in Eq. $`(1)`$. In this form, when $`\psi `$ is sent to $`\gamma ^5\psi `$, the scalar fields $`\varphi `$ and $`\lambda `$ are sent to their negatives ( minus times the field). This discrete symmetry prevents $`\psi `$ from acquiring a finite mass in higher orders.
We see that these two models are not equivalent since the latter does not obey one symmetry obeyed by the former one. If we take the original Lagrangian
$$L=\overline{\psi }_i/\psi +(\overline{\psi }\psi )^{4/3}+s(\overline{\omega }(\overline{\psi }\psi +\varphi ^3)),$$
(5)
and define a symmetry operation $`s`$ where $`s\overline{\omega }=\lambda ,s\lambda =0,s\varphi =\omega ,s\omega =0,s\psi =s\overline{\psi }=0`$, so that $`L`$ is invariant under $`s`$ . If we replace the original Lagrangian by that given in equation(1), by replacing $`(\overline{\psi }\psi )^{4/3}`$ by a combination of $`\varphi ^4,\varphi \overline{\psi }\psi `$ we see that this symmetry is not retained. We, therefore, take the second model as a model which only approximates the original one, without claiming equivalence. It is a constrained model which will replace in the original model only in a โnaiveโ sense
To quantize the latter system consistently we proceed through the path integral method. In addition to the usual spinor-Dirac primary constraints, fixing the momenta corresponding to the spinor fields $`\psi `$ and $`\overline{\psi }`$, we have two new primary constraints, setting the canonical momenta corresponding to the scalar fields $`\lambda `$ and $`\varphi `$ equal to zero. The primary Hamiltonian is obtained by adding these four constraints multiplied by arbitrary constants to the canonical Hamiltonian, obtained from the Lagrangian given in Eq. $`(1)`$. The consistency requirement of all the primary constraints, which is setting the Poisson bracket of the constraint equations with the primary hamiltonian equal to zero, results in two new, secondary constraints, given by our Eqs. $`(2)`$ and $`(3)`$ . When we calculate the Poisson bracket of these constraints with the primary Hamiltonian to check whether additional constraints are present, we see that the system is closed, determining all the arbitrary constants in the primary Hamiltonian.
Next we compute the determinant of the Poisson brackets of all the second class constraints, the so called Faddeev-Popov determinant. We see that the spinor-Dirac constraints, resulting from the canonical momenta of the spinor fields has no field dependent contribution to the Faddeev-Popov determinant. This determinant is given as
$$\mathrm{\Delta }_F=[det\{\theta _i,\theta _j\}_P]^{1/2}=det\varphi ^4.$$
(6)
the field dependent contribution coming from the constraints in Eqs.$`(2)`$ and $`(3)`$.
We write the generating functional for the Greenโs functions of the model as
$$Z=D\pi D\chi \delta (\theta _i)\mathrm{\Delta }_Fexp(i(\dot{\chi }\pi H_c)).$$
(7)
Here $`\chi `$ is the generic symbol for all the fields , $`\pi `$ is the generic symbol for all momenta and $`\theta `$ is the generic symbol for all the constraints in the model. Performing all the momenta integrals we obtain
$$Z=D\overline{\psi }D\psi D\varphi D\lambda \left(\frac{\mathrm{\Delta }_F}{3det\varphi ^2}\right)exp(iL^{}d^4x),$$
(8)
where
$$\frac{\mathrm{\Delta }_F}{det\varphi ^2}=det\varphi ^2,$$
(9)
This contribution is inserted into the Lagrangian using ghost fields $`c`$ and $`c^{}`$, and the resulting lagrangian reads
$$L^{\prime \prime }=\overline{\psi }[/+g(\varphi +\lambda )]\psi a\lambda \varphi ^3+ic^{}\varphi ^2c.$$
(10)
We can rewrite this expression by defining
$$\mathrm{\Phi }=\varphi +\lambda ,$$
(11)
$$\mathrm{\Lambda }=\varphi \lambda ,$$
(12)
as
$$L^{\prime \prime }=i\overline{\psi }[/+g\mathrm{\Phi }]\psi \frac{a}{16}(\mathrm{\Phi }^4+2\mathrm{\Phi }^3\mathrm{\Lambda }2\mathrm{\Phi }\mathrm{\Lambda }^3\mathrm{\Lambda }^4)+\frac{i}{4}c^{}(\mathrm{\Phi }^2+2\mathrm{\Phi }\mathrm{\Lambda }+\mathrm{\Lambda }^2)c.$$
(13)
By this transformation the $`\mathrm{\Lambda }`$ field is decoupled from the spinor sector of the lagrangian.
The integration over the spinor fields in the functional yields the effective action which is expressed in terms of $`\mathrm{\Phi },\mathrm{\Lambda }`$ and $`c`$, $`c^{}`$ fields only.
$`S_{eff}=Trln(i/+g\mathrm{\Phi })+{\displaystyle }d^4x[{\displaystyle \frac{a}{16}}(\mathrm{\Phi }^4+2\mathrm{\Phi }^3\mathrm{\Lambda }2\mathrm{\Phi }\mathrm{\Lambda }^3\mathrm{\Lambda }^4)`$
$`{\displaystyle \frac{i}{4}}c^{}(\mathrm{\Phi }^2+2\mathrm{\Phi }\mathrm{\Lambda }+\mathrm{\Lambda }^2)c].`$ (14)
The condition to get rid of the tadpole contribution , which is setting the first derivative of the effective action with respect to the $`\mathrm{\Phi }`$ and $`\mathrm{\Lambda }`$ fields to zero, gives us two equations
$$\frac{ig}{(2\pi )^4}Tr\frac{d^4p}{p/gv}\frac{a}{8}(2v^3+3v^2ss^3)=0,$$
(15)
and
$$\frac{a}{8}(v^33vs^22s^3)=0.$$
(16)
In these expressions -v and s are the vacuum expectation values of the fields $`\mathrm{\Phi }`$ and $`\mathrm{\Lambda }`$ respectively and the vacuum expectation value of the ghost fields are set to zero.
A consistent solution satisfying both equations is
$$s=v=0,$$
(17)
Since the $`\gamma ^5`$ symmetry is not dynamically broken, no mass is generated for the fermion dynamically. In this respect this model differs from the famous Gross-Neveu model, where this dynamical breaking takes place. It also differs from the Nambu-Jona-Lasinio model. The main reason for this behaviour is the conformal invariance present in the model . Gรผrseyโs original intention was to construct a conformal invariant model, at least classically. We find that upon quantization of our approximate model at least one phase exists which respects this symmetry .
The fermion propagator is the usual Dirac propagator in lowest order, as can be seen from the Lagrangian. The second derivative of the effective action with respect to the $`\mathrm{\Phi }`$ field gives us the induced inverse propagator for the $`\mathrm{\Phi }`$ field , with the infinite part given as
$$Inf\left[\frac{ig^2}{(2\pi )^4}\frac{d^4p}{p/(p/+q/)}\right]=\frac{g^2q^2}{4\pi ฯต}$$
(18)
Here dimensional regularization is used for the momentum integral and $`ฯต=4n`$. We see that the $`\mathrm{\Phi }`$ field propagates as a massless field.
When we study the propagators for the other fields, we see that no linear or quadratic term in $`\mathrm{\Lambda }`$ exists, so the one loop contribution to the $`\mathrm{\Lambda }`$ propagator is absent. Similarly the mixed derivatives of the effective action with respect to $`\mathrm{\Lambda }`$ and $`\mathrm{\Phi }`$ is zero at one loop, so no mixing between these two fields occurs. We can also set the propagators of the ghost fields to zero, since they give no contribution in the one loop approximation. The higher loop contributions are absent for these fields.
## 3 Dressed Fermion Propagator
In this section we calculate the above results in higher orders. To justify our result that no mass is generated for the fermion we may study the Bethe-Salpeter equation obeyed for this propagator. The Dyson-Schwinger equation for the spinor propagator is written as
$$iAp/+B=ip/+4\pi ฯต\frac{d^4q}{(iAq/+B)(pq)^2}.$$
(19)
Here $`iAp/+B`$ is the dressed fermion propagator. We use the one loop result for the scalar propagator. After rationalizing the denominator, we can take the trace of this expression over the $`\gamma `$ matrices to give us
$$B=4\pi ฯตd^4q\frac{B}{(A^2q^2+B^2)(pq)^2}.$$
(20)
The angular integral on the right hand side can be performed to give
$$B=4\pi ฯต\left[_0^{p^2}๐q^2\frac{q^2B}{p^2(A^2q^2+B^2)}+_{p^2}^{\mathrm{}}๐q^2\frac{B}{(A^2q^2+B^2)}\right].$$
(21)
If we differentiate this expression with respect to $`p^2`$ on both sides, we get
$$\frac{dB}{dp^2}=4\pi ฯต\left[_0^{p^2}๐q^2\frac{q^2B}{(p^2)^2(A^2q^2+B^2)}\right].$$
(22)
This integral is clearly finite. We get zero for the right hand side as $`ฯต`$ goes to zero. Since mass is equal to zero in the free case we get this constant equal to zero. This choice satisfies the Eq. $`(19)`$.
The similar argument can be used to show that $`A`$ is the dressed spinor propagator is a constant. We multiply Eq. $`(18)`$ by $`p/`$ and then take the trace over the spinor indices. We end up with
$$p^2A=p^2+4\pi ฯต\left[_0^{p^2}๐q^2\left(\frac{(q^4)A}{p^2(A^2q^2+B^2)}+_{p^2}^{\mathrm{}}๐q^2\frac{p^2A}{(A^2q^2+B^2)}\right)\right].$$
(23)
We divide both sides by $`p^2`$ and differentiate with respect to $`p^2`$. The end result
$$\frac{dA}{dp^2}=8\pi ฯต_0^{p^2}๐q^2\left(\frac{(q^4)A}{(p^2)^3(A^2q^2+B^2)}\right)$$
(24)
shows that $`A`$ is a constant as $`ฯต`$ goes to zero. Since the integral is finite, it equals unity for the free case, we take $`A=1`$.
## 4 Higher Orders
If our fermion field had a color index $`i`$ where $`i=1\mathrm{}N`$, we could perform an 1/N expansion to justify the use of only ladder diagrams for higher orders for the scattering processes. Although in our model the spinor has only one color, we still consider only ladder diagrams anticipating that one can construct a variation of the model with N colors.
We first see that we do not need infinite regularization for the $`<\overline{\psi }\psi \varphi >`$ vertex. The one loop correction to the tree vertex involves two fermion and one $`\varphi `$ propagator and one integration. The infinity coming from the momentum integration is cancelled by the $`ฯต`$ in the $`\varphi `$ propagator. All the higher order contributions vanish because the powers of $`ฯต`$ exceed the number of infinities coming momentum integrations. Indeed there is only one momentum integration that results in an infinity. We see that there is only a finite renormalization of the spinor-scalar coupling constant $`g`$.
We come to the same result after we write the Dyson-Schwinger equation for this vertex. We need the result of the four fermion scattering kernel to be able to perform this calculation. There is no four fermion coupling in our Lagrangian; so, this process will use at least one scalar propagator. Since the scalar particle propagator has a $`ฯต`$ factor, this process vanishes as $`ฯต`$ goes to zero. All higher orders, including the one loop contribution also vanish as $`ฯต`$ goes to zero, since they have higher powers of $`ฯต`$.
We can justify our claims also by writing the Bethe-Salpeter equations for this process. The Bethe-Salpeter equation for the fermion interaction reads as
$`G^{(2)}(p,q;P)=G_0^{(2)}(p,q;P)`$
$`+{\displaystyle \frac{1}{(2\pi )^8}}{\displaystyle d^4p^{}d^4q^{}G_0^{(2)}(p,p^{};P)K(p^{},q^{};P)G^{(2)}(q^{},q;P)}.`$ (25)
Here $`G_0^{(2)}(p,q;P)`$ is two non-crossing spinor lines, $`G^{(2)}(p,q;P)`$ is the proper four point function. $`K`$ is the Bethe-Salpeter kernel.
We note that this expression involves the four spinor kernel which we found to be of order $`ฯต`$. This expression can be written in the quenched ladder approximation , where the kernel is seperated into a scalar propagator with two spinor legs joining the proper kernel. If the proper kernel is of order $`ฯต`$ , the loop involving two spinors and a scalar propagator can be at most finite that makes the whole diagram in first order in $`ฯต`$. This fact also shows that there is no nontrivial spinor-spinor scattering.
We use this result in calculating the Dyson-Schwinger equation for the spinor-scalar vertex. This vertex involves the finite coupling constant $`g`$ plus the diagram where the scalar particle goes into two spinors which go to the four spinor Kernel. Here the $`ฯต`$ factor coming from the Kernel is cancelled with the loop integration. The loop does not involve any scalar propagators, so it diverges as $`\frac{1}{ฯต}`$. The result is finite renormalization of the three point vertex. Hence the spinor-scalar coupling constant does not run.
We see that the only infinite remormalization is needed for the four $`\varphi `$ vertex; hence the coupling constant for this process run. The first correction to the tree diagram is the box diagram. This diagram has four spinor propagators and give rise to a $`\frac{1}{ฯต}`$ type divergence. Since we included the four $`\varphi `$ term in our original lagrangian, we can renormalize the coupling constant of this vertex to incorporate this divergence. The finite part of this diagram is just a constant, renormalizing the initial coupling constant by a finite amount. There are no higher infinities for this vertex. The two loop diagram contains a $`\varphi `$ propagator which makes this diagram finite. The three-loop diagram is made out of eight spinor and two scalar lines. At worst we end up with a first order infinity of the form $`\frac{1}{ฯต}`$ using the dimensional regularization scheme. Higher order ladder diagrams give at worst the same type of divergence. This divergence for the four scalar vertex can be renormalized using standart means.
## 5 Conclusion
As a result of this analysis we end up with a model where there is no scattering of the fundamental fields, i.e. the spinors, whereas the composite fields, the scalar field, can take part in a scattering process. Here we do not give the exact expression for this amplitude, but it will be a series in $`a`$ and even powers of $`g`$, starting with $`g^4`$.
We can also have scattering processes where two scalar particles go to an even number of scalar particles. In the one loop approximation all these diagrams give finite results, like the case in the standard Yukawa coupling model. Since going to an odd number of scalars is forbidden by the $`\gamma ^5`$ invariance of the theory, we can also argue that scalar $`\varphi `$ particles can go to an even number of scalar particles only. This assertion is easily checked by diagrammatic analysis.
Any diagram which describes the process of producing spinor particles out of two scalars contains scalar propagators. The lowest of these diagrams where two scalars go to two spinors vanish since it either involves a triangle diagram made out of spinors, or a box diagram, made out of three spinors and one scalar. It vanishes due to fall of the scalar propagator in the latter case, although it is not zero unless the cut-off is removed . The diagram which involves the production of four spinors out of two scalars carries two scalar propagators and the diagram vanishes with the first power of $`ฯต`$.
As a result of our calculations we find a model which is trivial for the constituent spinor fields, whereas finite results are obtained for thee scattering of the composite scalar particles. The coupling constant for the scattering of the composite particles run, whereas the coupling constant for the spinor-scalar interaction does not run.
In the classical model, described by the Lagrangian given by Eq. $`(4)`$, we used one coupling constant $`g^{}`$, which is divided into two as $`g`$ and $`a`$ in Eq. $`(1)`$. We see that these two behave differently in the quantum case.
Our model is a toy model. We could not yet find a physical system that is effectively described by it.
Acknowledgement: We thank Ferhat Taลkฤฑn, Kayhan รlker, and an unknown correspondence for discussions and both scientific and technical assistance throughout this work. The work of M.H. is also supported by TUBA, the Academy of Sciences of Turkey. This work is also supported by TUBITAK, the Scientific and Technological Council of Turkey. |
warning/0506/hep-th0506147.html | ar5iv | text | # Carnot-Carathรฉodory metric and gauge fluctuations in Noncommutative Geometry
## I Introduction
Noncommutative geometry<sup>?</sup> enlarges differential geometry beyond the scope of Riemannian spin manifolds and gives access, as shown in various examples, to spaces obtained as the product of the continuum by the discrete. It allows one to describe in a single and coherent geometrical object the space-time of the Standard Model of elementary particles with massless neutrinos. Massive Dirac neutrinos are easily incorporated in the model<sup>?</sup> as long as one of them remain massless. Otherwise more substantial changes might be required. coupled with Euclidean general relativity <sup>?</sup>. Specifically, the diffeomorphism group of general relativity appears as the automorphism group of $`C^{\mathrm{}}\left(M\right)`$, the algebra of smooth functions over a compact Riemannian spin manifold $`M`$, while the gauge group of the strong and electroweak interactions emerges as the group $`U(๐_I)`$ of unitary elements of a finite dimensional algebra $`๐_I`$ (modulo a lift to the spinors<sup>?</sup>). Remarkably, unitaries not only act as gauge transformations but also acquire a metric significance via the so-called fluctuations of the metric. This paper aims to study in detail the analogy introduced in \[?\] between a simple kind of fluctuations of the metric, those governed by a connection $`1`$-form on a principal bundle, and the associated Carnot-Carathรฉodory metric.
A noncommutative geometry consists in a spectral triple
$$๐,,D$$
where $`๐`$ is an involutive algebra, commutative or not, $``$ a Hilbert space carrying a representation $`\mathrm{\Pi }`$ of $`๐`$ and $`D`$ a selfadjoint operator on $``$. Together with a chirality operator $`\mathrm{\Gamma }`$ and a real structure $`J`$ both acting on $``$, they satisfy a set of properties<sup>?</sup> providing the necessary and sufficient conditions for 1) an axiomatic definition of Riemannian spin geometry in terms of commutative algebra 2) its natural extension to the noncommutative framework. Points are recovered as pure states $`๐ซ(๐)`$ of $`๐`$, in analogy with the commutative case where
$`๐ซ(C^{\mathrm{}}\left(M\right))M`$ (1)
$`\omega _x(f)=f(x)`$ (2)
for any pure state $`\omega _x`$ and smooth function $`f`$. A distance $`d`$ between states $`\omega `$, $`\omega ^{}`$ of $`๐`$ is defined by
$$d(\omega ,\omega ^{})\underset{a๐}{sup}\left\{|\omega (a)\omega ^{}(a)|;[D,\mathrm{\Pi }(a)]1\right\}$$
(3)
where the norm is the operator norm on $``$. In the commutative case,
$$๐_E=C^{\mathrm{}}\left(M\right),_E=L_2(M,S),D_E=i\gamma ^\mu _\mu $$
(4)
with $`_E`$ the space of square integrable spinors and $`D_E`$ the ordinary Dirac operator of quantum field theory, $`d`$ coincides with the geodesic distance defined by the Riemannian structure of $`M`$. Thus (3) is a natural extension of the classical distance formula, all the more as it does not involve any notion ill-defined in a quantum framework such as the trajectory between points.
Carnot-Carathรฉodory metrics (or sub-Riemannian metrics)<sup>?</sup> are defined on manifolds $`P`$ equipped with a horizontal distribution, that is to say a (smooth) specification at any point $`pP`$ of a subspace $`H_pP`$ of the tangent space $`T_pP`$. The Carnot-Carathรฉodory distance $`d_H`$ between $`p`$ and $`q`$ is the length of the shortest path $`c`$ joining $`p`$ and $`q`$ whose tangent vector is everywhere horizontal,
$$d_H(p,q)=\underset{\dot{c}(t)H_{c(t)}P}{\text{Inf}}_0^1\dot{c}(t)๐t.$$
(5)
If there is no horizontal path from $`p`$ to $`q`$ then $`d_H(p,q)`$ is infinite. Any point at finite distance from $`p`$ is said accessible
$$\text{Acc}(p)\{qP;d_H(p,q)<+\mathrm{}\}.$$
(6)
Most often the norm in the integrand of (5) comes from an inner product in the horizontal subspace. The latter can be obtained in (at least) two ways: either by restricting to $`HP`$ a Riemannian structure of $`TP`$ or, when $`P\stackrel{๐}{}M`$ is a fiber bundle with a connection, by pulling back the Riemannian structure $`g`$ of $`M`$. In the latter case the horizontal distribution is the kernel of the connection $`1`$-form and any horizontal vector has norm
$$u\pi _{}(u)=\sqrt{g(\pi _{}(u),\pi _{}(u))}.$$
(7)
Note that (5) provides $`P`$ with a distance although $`P`$ may not be a metric manifold, only $`M`$ is asked to be Riemannian.
By taking the product of a Riemannian geometry (4) by a spectral triple with finite dimensional $`๐_I`$, one obtains as pure state space a $`U(๐_I)`$-bundle $`P`$ over $`M`$. A connection on $`P`$ then not only defines a Carnot-Carathรฉodory distance $`d_H`$ but also, via the process of fluctuation of the metric recalled in section II, a distance $`d`$ similar to (3) except that the ordinary Dirac operator $`D`$ is replaced by the covariant differentiation operator associated to the connection-1 form. In section III we compare the connected components for these two distances: while a connected component for $`d_H`$ is also connected for $`d`$, a connected component of $`d`$ is not necessarily connected for $`d_H`$. We investigate the importance of the holonomy group on this matter. In section IV we show that the two distances coincide when the holonomy is trivial. In the non-trivial case we work out some necessary conditions on the holonomy group that may allow $`d`$ to equal $`d_H`$. In section V we treat in detail a simple low-dimensional example in which each of the connected components of $`d_H`$ is a dense subset of a two dimensional torus $`๐`$. As a main result of this paper we show in section VI that while the Carnot-Carathรฉodory metric forgets about the fiber bundle structure of $`๐`$, the noncommutative metric deforms it in a quite intriguing way: from a specific intrinsic point of view, the fiber acquires the shape of a cardioid. Hence the classical $`2`$-torus inherits a metric which is โtrulyโ noncommutative in the sense that it cannot be described in (sub)Riemannian nor discrete terms. This is, to our knowledge, the novelty of the present work.
Notations and conventions:
* $`M`$ is a Riemannian compact spin manifold of dimension $`m`$ without boundary. Cartesian coordinates are labeled by Greek indices $`\mu ,\nu `$ and we use Einstein summation over repeated indices in alternate positions (up/down).
* $`๐ซ(๐)`$ denotes the set of pure states of $`๐`$ (positive, linear applications from $`๐`$ to $``$, with norm $`1`$ and that do not decompose as a convex combination of other states). Throughout this paper we deal with a pure state space which is a trivial bundle $`P`$ over $`M`$, with fiber $`P^{n1}`$. An element of $`P`$ is written $`\xi _x`$ where $`x`$ is a point of $`M`$ and $`\xi P^{n1}`$.
* Most of the time we omit the symbol $`\mathrm{\Pi }`$ and it should be clear from the context whether $`a`$ means an element of $`๐`$ or its representation on $``$. Unless otherwise specified a bracket denotes the scalar product on $`^n`$.
* We use the result of \[?\] according to which the supremum in (3) can be sought on positive elements of $`๐`$.
## II Fluctuations of the metric
In noncommutative geometry a connection on a geometry $`(๐,,D)`$ is defined via the identification of $`๐`$ as a finite projective module over itself (i.e. as the noncommutative equivalent of the sections of a vector bundle via the Serre-Swan theorem)<sup>?</sup>. It is implemented by replacing $`D`$ with a covariant operator
$$D_AD+A+JAJ^1$$
(8)
where $`J`$ is the real structure and $`A`$ is a selfadjoint element of the set $`\mathrm{\Omega }^1`$ of 1-forms
$$\mathrm{\Omega }^1\left\{a^i[D,b_i];a^i,b_i๐\right\}.$$
(9)
Only the part of $`D_A`$ that does not obviously commute with the representation, namely
$$๐D+A,$$
(10)
enters in the distance formula (3) and induces a so-called fluctuation of the metric. In the following we consider almost commutative geometries obtained as the product of the continuous - external - geometry (4) by an internal geometry $`๐_I,_I,D_I`$. The product of two spectral triples, defined as
$$๐=๐_E๐_I,=_E_I,D=D_E๐_I+\gamma ^5D_I$$
(11)
where $`๐_I`$ is the identity operator of $`_I`$ and $`\gamma ^5`$ the chirality of the external geometry, is again a spectral triple. The corresponding 1-forms are<sup>?,?</sup>
$$i\gamma ^\mu f_\mu ^im_i+\gamma ^5h^jn_j$$
where $`m_i๐_I`$, $`h^j,f_\mu ^iC^{\mathrm{}}\left(M\right)`$, $`n_j\mathrm{\Omega }_I^1`$. Selfadjoint $`1`$-forms decompose into an $`๐_I`$-valued skew-adjoint 1-form field over $`M`$, $`A_\mu f_\mu ^im_i`$, and an $`\mathrm{\Omega }_I^1`$-valued selfadjoint scalar field $`Hh^jn_j.`$
When the internal algebra $`๐_I`$ has finite dimension, $`A_\mu `$ takes values in the Lie algebra of unitaries of $`๐`$ and is called the gauge part of the fluctuation. In \[?\] we have computed the noncommutative distance (3) for a scalar fluctuation only ($`A_\mu =0`$). In \[?\] the distance is considered for a pure gauge fluctuation ($`H=0`$) obtained from the internal geometry
$$๐_I=M_n(),_I=M_n(),D_I=0,$$
that is to say
$$๐=i\gamma ^\mu (_\mu ๐_I+๐_EA_\mu ).$$
(12)
$`๐_E`$ being nuclear, the set of pure states of
$$๐=C^{\mathrm{}}\left(M\right)M_n\left(\right)=C^{\mathrm{}}(M,M_n\left(\right))$$
(13)
is <sup>?</sup> $`๐ซ(๐)๐ซ(๐_E)\times ๐ซ(๐_I)`$, where $`๐ซ(๐_I)`$ is the projective space $`P^{n1}`$,
$$\omega _\xi (m)=\xi ,m\xi =\text{Tr}(s_\xi m)$$
(14)
for $`m๐_I,\xi P^{n1}`$ and $`s_\xi `$ the support of $`\omega _\xi `$. The evaluation of $`\xi _x(\omega _x,\omega _\xi )`$ on $`a=f^im_i๐`$ reads
$$\xi _x(a)=\text{Tr}(s_\xi a(x))$$
(15)
where
$$a(x)f^i(x)m_i.$$
(16)
Hence $`๐ซ(๐)`$ is a trivial bundle
$$P\stackrel{๐}{}M$$
with fibre $`P^{n1}`$.
The gauge potential $`A_\mu `$ defines both a horizontal distribution $`H`$ on $`P`$, with associated Carnot-Carathรฉodory metric $`d_H`$, and a noncommutative metric $`d`$ given by formula (3) with $`๐`$ substituted for $`D`$. In the case of a zero connection, $`๐=D_E`$ and $`d`$ is the geodesic distance on $`M`$. Indeed the commutator norm condition $`[D_E,f]1`$ forces the gradient of $`f`$ to be smaller than $`1`$, so that
$$d(\omega _x,\omega _y)=\underset{\text{grad}f1}{\text{sup}}|f(x)f(y)|_0^1\dot{c}(t)๐t=d_{\text{geo}}(x,y)$$
(17)
where $`c`$, $`c(0)=x`$, $`c(1)=y`$ is a minimal geodesic from $`x`$ to $`y`$. One then easily checks that this upper bound is attained by the function
$$L(z)d_{\text{geo}}(x,z)zM$$
(18)
(or more precisely by a sequence of smooth functions converging to the continuous function $`L`$). As we shall see in the following section, in the case of a non-zero connection, one obtains without difficulty a result similar to (17) with $`d_H`$ playing the role of $`d_{\text{geo}}`$ (cf eq. (19) below). However, except in some simple cases studied in section IV, $`d_H`$ is not the least upper bound and there is no simple equivalent to the function $`L`$. In fact the main part of this paper, especially section V, is devoted to building the element $`a๐`$ that realizes the supremum in the distance formula.
## III Connected components
We say that two pure states $`\xi _x`$, $`\zeta _y`$ are connected for $`d`$ if and only if $`d(\xi _x,\zeta _y)`$ is finite.
###### Proposition III.1
For any $`\xi _x`$ in $`P`$, $`\text{Acc}(\xi _x)`$ is connected for $`d`$.
Proof. The result is obtained by showing that for any $`\zeta _y\text{Acc}(\xi _x)`$,
$$d(\xi _x,\zeta _y)d_H(\xi _x,\zeta _y).$$
(19)
Let us start by recalling how to transfer the covariant derivative<sup>?</sup> of a section $`V`$ of $`P`$,
$$_\mu V=_\mu V+A_\mu V,$$
to the algebra $`๐`$. Given $`a๐`$, the evaluation (15) is the diagonal of the sesquilinear form defined fiberwise on the vector bundle $`P^{}\stackrel{\pi ^{}}{}M`$ with fiber $`^n`$,
$$(W_x^{},V_x^{})W_x^{},a(x)V_x^{}$$
(20)
for $`W_x^{},V_x^{}\pi ^1(x)`$. Accordingly, as a $`C^{\mathrm{}}\left(M\right)`$-module, we view $`๐`$ as the sections of the bundle $`P^{\prime \prime }`$ of rank-two tensors on $`M`$
$$a=a_{ij}\overline{e^i}e^j$$
with values in $`\overline{T^{}^n}T^{}^n`$. Here $`\{e^i\}`$ is the dual of the canonical basis $`\{e_i\}`$ of $`T^n^n`$ and $`\{\overline{e^i}\}`$ its complex conjugate
$$\overline{e^i}(V)=\overline{V^i}\text{for}V=V^ie_i^n.$$
The covariant derivative of $`P`$ then naturally extends to $`P^{\prime \prime }`$ <sup>ยง</sup><sup>ยง</sup>ยง $`\{\begin{array}{c}_\mu e^i=A_{\mu k}^ie^k\\ _\mu \overline{e^i}=\overline{A_{\mu k}^i}\overline{e^k}\end{array}`$ hence $`_\mu a_\mu (a_{ij}\overline{e^i})e^j+a_{ij}\overline{e^i}_\mu e^j=(_\mu a_{ij}+[A,a]_{ij})\overline{e^i}e^j.`$
$$_\mu a=_\mu a+[A_\mu ,a].$$
(21)
Let us fix a horizontal curve of pure states $`c(t)`$, $`t[0,1]`$, between $`\xi _x`$ and $`\zeta _y`$ as defined in (15). Let $`(\pi ,V)`$ be a trivialization in $`P`$ such that
$$\pi (\xi _x)=x,V(\xi _x)=\xi \pi (\zeta _y)=y,V(\zeta _y)=\zeta $$
(22)
and define
$$V(t)V(c(t)).$$
$`c`$ is the horizontal lift starting at $`\xi _x`$ of the curve
$$c_{}(t)\pi (c(t))$$
lying in $`M`$ and tangent to
$$\pi _{}(\dot{c})=\dot{c}_{}=\dot{c}_{}^\mu _\mu .$$
(23)
Writing $`s(t)`$ for the support of the pure state $`\omega _{V(t)}`$, the curve $`ts(t)`$ is horizontal in $`P^{\prime \prime }`$ in the sense of the covariant derivative (21) in Dirac notation $`c`$ horizontal in $`P`$ is written $`\dot{|V}+\dot{c}^\mu A_\mu |V=0`$. By simple manipulations $`\{\begin{array}{c}\dot{|V}V|+\dot{c}^\mu A_\mu |VV|=0\\ |V\dot{V|}|VV|\dot{c}^\mu A_\mu =0\end{array},`$ hence $`\dot{s}=|V\dot{V|}+\dot{|V}V|=\dot{c}^\mu [|VV|,A_\mu ]=\dot{c}^\mu [s,A_\mu ]`$.
$$_{\dot{c}_{}}s\dot{c}_{}^\mu _\mu s=0.$$
(24)
Let us associates to any $`a๐`$ its evaluation $`f`$ along $`c`$,
$$f(t)\text{Tr}(s(t)a(c_{}(t))),$$
(25)
whose derivative with respect to $`t`$ is easily computed using (24)
$$\dot{f}=\text{Tr}(s_{\dot{c}_{}}a).$$
(26)
At a given $`t`$ the Cauchy-Schwarz inequality yields the bound
$$|\dot{f}(t)|df_{|t}\dot{c}_{}(t)$$
(27)
where $`df`$ is the $`1`$-form on $`c_{}`$ with components
$$_\mu f=\text{Tr}(s_\mu a).$$
(28)
$`s[๐,a]s`$ evaluated at some $`c_{}(t)`$ is an $`n^{}\times n`$ square matrix ($`n^{}=\text{dim}_E`$ is the dimension of the spin representation),
$$s[๐,a]s=i\gamma ^\mu s(_\mu a)s=i\gamma ^\mu _\mu fs,$$
(29)
with norm $`df_{|t}`$. Therefore
$$df_{|t}\underset{xM}{sup}[๐,a]_{|x}=[๐,a]$$
(30)
so, as soon as $`[๐,a]1`$,
$$|\xi _x(a)\zeta _y(a)|=|_{0}^{}{}_{}{}^{1}\dot{f}(t)๐t|_{0}^{}{}_{}{}^{1}\dot{c}_{}(t)๐t,$$
(31)
which precisely means $`d(\xi _x,\zeta _y)d_H(\xi _x,\zeta _y)`$. $`\mathrm{}`$
It would be tempting to postulate that $`d`$ and $`d_H`$ have the same connected components. Half of this way is done in the proposition above. The other half would consist in checking that $`d`$ is infinite as soon as $`d_H`$ is infinite. However this is, in general, not the case. It seems that there is no simple conclusion on that matter since we shall exhibit in section V an example in which some states that are not in $`\text{Acc}(\xi _x)`$ are at finite noncommutative distance from $`\xi _x`$ whereas others are at infinite distance. The best we can do for the moment is to work out (Proposition III.3 below) a sufficient condition on the holonomy group associated to the connection $`A_\mu `$ that guarantees the non-finiteness of $`d(\xi _x,\zeta _y)`$ for $`\zeta _y\text{Acc}(\xi _x)`$. We begin with the following elementary lemma.
###### Lemme III.2
$`d(\xi _x,\zeta _y)`$ is infinite if and only if there is a sequence $`a_n๐`$ such that
$$\underset{n+\mathrm{}}{\text{lim}}[D,a_n]0,\underset{n+\mathrm{}}{\text{lim}}|\xi _x(a_n)\zeta _y(a_n)|=+\mathrm{}.$$
(32)
Proof. The point is to show that from a sequence $`a_n`$ satisfying
$$[D,a_n]1n,\underset{n+\mathrm{}}{\text{lim}}|\xi _x(a_n)\zeta _y(a_n)|=+\mathrm{}$$
one can extract a sequence $`\stackrel{~}{a}_n`$ satisfying (32). This is done by considering
$$\stackrel{~}{a}_n\frac{a_n}{\sqrt{|\xi _x(a_n)\zeta _y(a_n)|}}.$$
$`\mathrm{}`$
###### Proposition III.3
Let $`\xi ,\zeta P^{n1}`$. If there exists a matrix $`MM_n\left(\right)`$ that commutes with the holonomy group at $`x`$, $`\text{Hol}(x)`$, and such that
$$\text{Tr}(s_\xi M)\text{Tr}(s_\zeta M),$$
(33)
then $`d(\omega ,\omega ^{})=+\mathrm{}`$ for any $`\omega \text{Acc}(\xi _x)`$, $`\omega ^{}\text{Acc}(\zeta _x)`$.
Proof. The proof is a restatement of a classical result (cf \[?\] p.113) according to which an element of $`๐`$ invariant under the adjoint action of the holonomy group is a parallel tensor, that is to say $`_\mu a=0`$ in all directions $`\mu `$. We detail this point in the following for the sake of completeness.
From now on we fix a trivialization $`(\pi ,V)`$ on $`P=๐ซ(๐)`$. Recall that given a curve from $`c_{}(0)=x`$ to $`c_{}(1)=yM`$, the end point of the horizontal lift $`c`$ of $`c_{}`$ with initial condition $`c(0)=(x,\xi )`$ is $`c(1)=(y,U_c_{}(1)\xi )`$ where
$$U_c_{}(t)=P\mathrm{exp}(_{c_{}(t)}A_\mu ๐x^\mu )$$
($`P`$ is the time-ordered product) is the solution of
$$\dot{U}=\dot{c}^\mu A_\mu U.$$
(34)
In the following we write $`U_c_{}`$ for $`U_c_{}(1)`$. Let $`MM_n\left(\right)`$ commute with $`\text{Hol}(x)`$. Define $`a_M๐`$ by
$$a_M(x)M$$
and for any $`yM`$,
$$a_M(y)U_c_{}a_M(x)U_c_{}^{}$$
(35)
where $`c_{}`$ is a curve joining $`x`$ to $`y`$. One checks that $`a_M(y)`$ commutes with any $`V_l\text{Hol}(y)`$ since
$$V_la_M(y)V_l^{}=U_c_{}U_c_{}^{}V_lU_c_{}a_M(x)U_c_{}^{}V_l^{}U_c_{}U_c_{}^{}=a_M(y)$$
where we use that $`U_c_{}^{}V_lU_c_{}`$ belongs to $`\text{Hol}(x)`$. Hence (35) uniquely defines $`a_M(y)`$ since parallel transporting $`a_M(x)`$ along another curve $`c_{}^{}`$ yields
$$a_{}^{}{}_{M}{}^{}(y)=U_c_{}^{}U_c_{}^{}a_M(y)U_c_{}U_c_{}^{}^{}=a_M(y)$$
where we used that $`U_c_{}U_c_{}^{}^{}\text{Hol}(y)`$. Using (34) one explicitly checks that
$$_{\dot{c}_{}}a_M=0.$$
Since this is true for any curve $`c_{}`$, $`a_M`$ is parallel so
$$[๐,a_M]=0.$$
Now (33) means that $`\xi _x(a_M)\zeta _x(a_M)0`$, hence $`d(\xi _x,\zeta _x)=+\mathrm{}`$ by lemma III.2, and the result follows by the triangle inequality. $`\mathrm{}`$
Proposition above only provides sufficient conditions. Whether they are necessary, i.e. whether from $`d(\xi _x,\zeta _y)=+\mathrm{}`$ on can build a matrix $`M`$ that commutes with the holonomy group and do not cancel the difference of the states is an open question. Lemma III.2 suggests that to any infinite distance is associated a tensor that commutes with the Dirac operator. Moreover it is not difficult to show that any parallel tensor commutes with the holonomy group. Therefore the question is: are the parallel tensors the only ones that commute with $`D`$ ? For the time being the answer is not clear to the author.
To close this section, let us mention a situation in which the two metrics have the same connected components.
###### Corollary III.4
If for a given $`\xi _xP`$ the vector space
$$_{hol}\text{Span}\{U\xi ;U\text{Hol}(x)\}$$
has dimension $`h<n`$, then $`\text{Acc}(\xi _x)`$ is the connected component of $`\xi _x`$ for $`d`$.
Proof. In an orthonormal basis $`\{_{hol},\}`$ of $`^n`$ with $`_{hol}`$ a basis of $`_{hol}`$, $`\text{Hol}(x)`$ is block represented, so
$$M=\left(\begin{array}{cc}0& 0\\ 0& ๐_{nh}\end{array}\right)$$
commutes with $`\text{Hol}(x)`$. Moreover $`\text{Tr}(s_\xi M)=0`$. On the contrary for any $`\zeta _x\text{Acc}(\xi _x)`$, the rank one projector $`s_\zeta `$ does not project on $`_{hol}`$ so $`\text{Tr}(s_\zeta M)0`$. Therefore, by Proposition III.3, $`d(\xi _x,\zeta _y)`$ is infinite for any $`\zeta _y\text{Acc}(\xi _x)`$, hence the result by Proposition III.1. $`\mathrm{}`$
## IV Flat case versus holonomy constraints
The preceding section suggests that the two metrics defined by a connection on the pure state space $`P`$ of the algebra (13), the Carnot-Carathรฉodory distance $`d_H`$ and the noncommutative distance $`d`$, do not coincide. It is likely that the two metrics do not have the same connected components as soon as the conditions of Proposition III.3 are not fulfilled. However nothing forbids $`d`$ from equalling $`d_H`$ on each connected component of $`d`$. We already know that $`dd_H`$ so to obtain the equality it would be enough to exhibit one positive $`a๐`$ (or a sequence of elements $`a_n`$) satisfying the commutator norm condition as well as
$$\xi _x(a)\zeta _y(a)=d_H(\xi _x,\zeta _y).$$
(36)
The existence of such an $`a`$ strongly depends on the holonomy of the connection: when the latter is trivial, e.g. by the Ambrose-Singer theorem when the connection is flat and $`M`$ simply connected, then the two metrics are equal, as shown below in Proposition IV.1. When the holonomy is non-trivial, we work out in Proposition IV.4 some necessary conditions on the shortest path that may forbid $`d`$ from equalling $`d_H`$.
###### Proposition IV.1
When the holonomy group reduces to the identity, $`d=d_H`$ on all $`P`$.
Proof. For $`\zeta _y\text{Acc}(\xi _x)`$, Corollary III.4 yields
$$d(\xi _x,\zeta _y)=+\mathrm{}=d_H(\xi _x,\zeta _y).$$
Thus we focus on the case $`\zeta _y\text{Acc}(\xi _x)`$. By Cartanโs structure equation the horizontal distribution defined by a connection with trivial holonomy is involutive, which means that the set of horizontal vector fields is a Lie algebra for the Lie bracket inherited from $`TP`$. Equivalently (Frobenius theorem) the bundle of horizontal vector fields is integrable. Hence $`\text{Acc}(\xi _x)`$ is a submanifold of $`P`$, call it $`\mathrm{\Xi }`$, such that $`Tp\mathrm{\Xi }=H_pP`$ for any $`p\mathrm{\Xi }`$. For any $`zM`$ there is at least $`1`$ point in the intersection
$$\pi ^1(z)\mathrm{\Xi }$$
(e.g. the end point of the horizontal lift, starting at $`\xi _x`$, of any curve from $`x`$ to $`z`$) and only one (otherwise there would be a horizontal curve joining two distinct points in the fiber, contradicting the triviality of the holonomy). In other words all the horizontal lifts starting at $`\xi _x`$ of curves joining $`x`$ to $`z`$ have the same end point, call it $`\sigma (z)`$, and the application
$$\sigma :z\pi ^1(z)\mathrm{\Xi }$$
defines a smooth section of $`P`$. Hence
$$\mathrm{\Xi }=\sigma (M).$$
Note that $`\zeta _y=\sigma (y)`$ is the only point in the fiber over $`y`$ which is at finite distance from $`\xi _x=\sigma (x)`$. Considering the horizontal lift of the Riemannian geodesic from $`x`$ to $`y`$, it turns out that $`d_H`$ on $`\mathrm{\Xi }`$ coincides with the geodesic distance $`d_{\text{geo}}`$ on $`M`$. The sequence of elements $`a_n`$ we are looking for in (36) is a sequence approximating the continuous $`M_2\left(\right)`$-valued function
$$L๐$$
(37)
where $`L`$ is the geodesic distance function (18). $`\mathrm{}`$
The difficulty arises when the shortest horizontal curve $`c`$ does not lie in a horizontal section. This certainly happens when the connection is not flat and/or $`M`$ not simply connected. As soon as the holonomy is non-trivial, different points $`\xi _x`$, $`\zeta _x`$ on the same fiber can be at finite non-zero Carnot-Carathรฉodory distance from one another although the Riemannian distance of their projections vanishes. The question reduces to finding the equivalent of the element (37) in the closure of $`๐`$ that attains the supremum in (36). A natural candidate to play the role of the function $`L`$ in the case of a non-trivial holonomy is the fiber-distance function which associates to any $`zM`$ the length of the shortest horizontal path joining $`\xi _x`$ to some point in $`\pi ^1(z)`$. When the holonomy is trivial this function precisely coincides with $`L`$. However there is no natural candidate to play the role of the identity matrix in (37). Possibly one might determine by purely algebraic techniques which element $`a`$ of $`๐`$ realizes the supremum in the distance formula. The best approach we found for the moment is to work out, Proposition IV.4, some conditions between the matrix part of $`a`$ and the self-intersecting points of $`c_{}`$ that are necessary for $`d`$ to equal $`d_H`$.
###### Definition IV.2
Given a curve $`c`$ in a fiber bundle with horizontal distribution $`H`$, we call a c-ordered sequence of $`K`$ self-intersecting points at $`p_0`$ a set of at least two elements $`\{c(t_0),c(t_1),\mathrm{},c(t_K)\}`$ such that
$$\pi (c(t_i))=\pi (c(t_0)),d_H(c(t_0),c(t_i))>d_H(c(t_0),c(t_i))$$
for any $`i=1,\mathrm{},K`$ (Figure 1).
###### Lemme IV.3
Let $`\xi _x`$, $`\zeta _y`$ be two points in $`P`$ such that $`d(\xi _x,\zeta _y)=d_H(\xi _x,\zeta _y)`$. Then for any $`c(t)`$ belonging to a minimal horizontal curve $`c`$ between $`c(0)=\xi _x`$ and $`c(1)=\zeta _y`$,
$$d(\xi _x,c(t))=d_H(\xi _x,c(t)).$$
(38)
Moreover, for any such curve there exists an element $`a๐`$ (or a sequence $`a_n`$) such that
$$\xi _t(a)=d_H(\xi _x,c(t))$$
(39)
for any $`t[0,1]`$, where $`\xi _t`$ denotes $`c(t)`$ viewed as a pure state of $`๐`$.
Proof. We write the proof assuming that the supremum in the distance formula is attained by some $`a๐`$. In case the supremum is not reached, the proof is identical using a sequence $`\{a_n\}`$. Assume $`a`$ does satisfies the commutator norm condition as well as (36). Let us parameterize $`c`$ by its length element $`\tau `$ and use โdotโ for the derivative $`\frac{d}{d\tau }`$. The function $`f(t)=\xi _t(a)`$ defined by (25) has constant derivative along $`c_{}`$. Indeed (36) reads
$$_0^\mathrm{\Lambda }\dot{f}(\tau )๐\tau =\mathrm{\Lambda }$$
(40)
where $`\mathrm{\Lambda }=d_H(\xi _x,\zeta _y)`$. Since $`\dot{c}_{}(\tau )=1`$ for any $`\tau [0,\mathrm{\Lambda }]`$, (27) and (30) forbid $`|\dot{f}(\tau )|`$ from being greater than $`1`$. Hence
$$\dot{f}(\tau )=1$$
(41)
for almost every $`\tau `$. Thus for any $`\lambda \mathrm{\Lambda }`$,
$$_0^\lambda \dot{f}(\tau )๐\tau =\lambda $$
(42)
which reads
$$\xi _\lambda (a)\xi _x(a)=\lambda =d_H(\xi _x,\xi _\lambda ).$$
(43)
Hence (38) by Proposition III.1, and (39) by considering $`\stackrel{~}{a}a\xi _x(a)`$. $`\mathrm{}`$
Applying lemma IV.3 to the self-intersecting points defined in IV.2 one obtains the announced necessary conditions for $`d`$ to equal $`d_H`$.
###### Proposition IV.4
The noncommutative distance between two points $`\xi _x`$, $`\zeta _y`$ in $`P`$ can equal the Carnot-Carathรฉodory one only if there exists a minimal horizontal curve $`c`$ between $`\xi _x`$ and $`\zeta _y`$ such that there exists an element $`a๐`$, or a sequence of elements $`a_n`$, satisfying the commutator norm condition as well as
$$\xi _{t_i}(a)=d_H(\xi _x,c(t_i))\text{ or }\underset{n\mathrm{}}{\text{lim}}\xi _{t_i}(a_n)=d_H(\xi _x,c(t_i))$$
(44)
for any $`\xi _{t_i}=c(t_i)`$ in any c-ordered sequence of self-intersecting points.
Given a sequence of $`K`$ self-intersecting points at $`p`$, Proposition IV.4 puts $`K+1`$ condition on the $`n^2`$ real components of the selfadjoint matrix $`a(\pi (p))`$. So it is most likely that a necessary condition for $`d(\xi _x,\zeta _y)`$ to equal $`d_H(\xi _x,\zeta _y)`$ is the existence of a minimal horizontal curve between $`\xi _x`$ and $`\zeta _y`$ such that its projection does not self-intersect more than $`n^21`$ times. We will refine this interpretation in the example of the next section. From a more general point of view it is not clear how to deal with such a condition in the framework of sub-Riemannian geometryThanks to R. Montgomery<sup>?</sup> for illuminating discussions on this matter.. It might be possible indeed that in a manifold of dimension greater than $`3`$ one may, by smooth deformation, reduce the number of self-intersecting points of a minimal horizontal curve. But this is certainly not possible in dimension $`2`$ or $`1`$. In particular, when the basis is a circle there is only one horizontal curve $`c`$ between two given points, and it is not difficult to find a connection such that $`c_{}`$ self-intersects infinitely many times. This is what motivates the following example.
## V The example $`C^{\mathrm{}}(S^1)M_2\left(\right)`$
Let us summarize our comparative analysis of $`d`$ and $`d_H`$. When the holonomy is trivial the two distances are equal by proposition IV.1. When the holonomy is non-trivial we have both:
\- a sufficient, but maybe not necessary, condition (Corollary III.4) that guarantees the two distances have the same connected components,
-a necessary condition (Proposition IV.4) for the two distances to coincide on a given connected component.
These two conditions do not seem to be related: writing $`Q^i`$ and $`Q_H^i`$ the connected components of $`d`$ and $`d_H`$ respectively, it is likely that in some situations $`Q^i=Q_H^i`$ for some $`i`$ although $`d`$ differs from $`d_H`$ on $`Q^i`$, or on the contrary $`Q_H^iQ^i`$ but $`d=d_H`$ on $`Q_h^i`$. In the present section we exhibit a simple low-dimensional example in which the $`Q^i`$โs are two dimensional tori (Proposition V.1) and the $`Q_H^i`$โs are dense subsets. $`d`$ coincides with $`d_H`$ only on some part of $`Q_H^i`$ (Corollary VI.2). The present section is technical and deals with the exact computation of the noncommutative distance (Proposition V.4). Interpretation and discussion are postponed to the following section.
Consider the trivial $`U(2)`$-bundle $`P`$ over the circle $`S^1`$ of radius one with fiber $`P^1`$, that is to say the set of pure states of $`๐=๐_E๐_I`$ with $`๐_E=C^{\mathrm{}}(S^1)`$ and $`๐_I=M_2\left(\right)`$, namely
$$๐=C^{\mathrm{}}(S^1,M_2\left(\right)).$$
Let us equip $`P`$ with a connection whose associated $`1`$-form $`A๐ฒ(2)`$ is constant. For simplicity we restrict to a matrix $`A`$ of rank one but the adaptation to a wider class of connections should be quite straightforward. Once and for all we fix a basis of $`^2`$ in which the fundamental representation of $`A`$ is written
$$A=\left(\begin{array}{cc}0& 0\\ 0& i\theta \end{array}\right)$$
(45)
where $`\theta ]0,1[`$ is a fixed real parameter. Let $`[0,2\pi [`$ parameterize the circle and call $`x`$ the point with coordinate $`0`$. Within a trivialization $`(\pi ,V)`$ the horizontal lift $`c`$ of the curve
$$c_{}(\tau )=\tau \text{ mod }[2\pi ],\tau ]\mathrm{},+\mathrm{}[$$
(46)
with initial condition
$$V(c(0))=\xi =\left(\begin{array}{c}V_1\\ V_2\end{array}\right)P^1$$
is the helix $`c(\tau )=(c_{}(\tau ),V(\tau ))`$, where
$$V(\tau )=\left(\begin{array}{c}V_1\\ V_2e^{i\theta \tau }\end{array}\right).$$
(47)
The points of $`P`$ accessible from $`\xi _x=\xi _0(\omega _{c_{}(0)},\omega _\xi )`$ are the pure states
$$\xi _\tau (\omega _{c_{}(\tau )},\omega _{V(\tau )}).$$
(48)
By the Hopf fibration the fiber $`P^1`$ is seen to be a two sphere. Explicitly $`\xi `$ is the point of $`S^2`$ with Cartesian coordinates
$$x_\xi =2\text{Re}(V_1\overline{V_2}),y_\xi =2\text{Im}(V_1\overline{V_2}),z_\xi =|V_1|^2|V_2|^2.$$
(49)
Writing
$$2V_1\overline{V_2}Re^{i\theta _0}$$
(50)
one obtains $`\xi _x`$ as the point in the fiber $`\pi ^1(x)`$ with coordinates
$$x_0=R\mathrm{cos}\theta _0,y_0=R\mathrm{sin}\theta _0,z_0=z_\xi .$$
The points in the fiber over $`c_{}(\tau )`$ that are accessible from $`\xi _x`$ are
$$\xi _\tau ^k\xi _{\tau +2k\pi },k,$$
(51)
with Hopf coordinates
$$x_\tau ^kR\mathrm{cos}(\theta _0\theta _\tau ^k),y_\tau ^kR\mathrm{sin}(\theta _0\theta _\tau ^k),z_\tau ^kz_\xi $$
(52)
where
$$\theta _\tau ^k\theta (\tau +2k\pi ).$$
All the $`\xi _\tau ^k`$โs are on the circle $`S_R`$ of radius $`R`$ located at the โaltitudeโ $`z_\xi `$ in $`\pi ^1(c_{}(\tau ))`$. Therefore
$$\text{Acc}(\xi _x)๐_\xi $$
where
$$๐_\xi S^1\times S_R$$
(53)
is the two-dimensional torus (see Figure 2).
Similarly for any $`\zeta P^1`$ such that $`z_\xi =z_\zeta `$ one has $`\text{Acc}(\zeta _x)๐_\xi `$. In fact
$$๐_\xi =\underset{\underset{z_\zeta =z_\xi }{\zeta P^1,}}{}\text{Acc}(\zeta _x).$$
(54)
Note that when $`\theta `$ is irrational $`๐_\xi `$ is the completion of $`\text{Acc}(\xi _x)`$ with respect to the Euclidean norm on each $`S_R`$.
###### Proposition V.1
$`๐_\xi `$ is the connected component of $`\xi _x`$ for $`d`$.
Proof. Let $`a_{ij}๐_E`$, $`i,j=1,2`$, be the components of a selfadjoint element of $`๐`$. (46) yields an explicit identification of $`๐_E`$ with the algebra of $`2\pi `$-periodic complex functions on $``$,
$$a_{ij}(\tau )a_{ij}(c_{}(\tau ))=a_{ij}(\tau +2k\pi )k$$
(55)
with
$$a_{ij}(0)=a_{ij}(x).$$
Let dot denote the derivative. Since $`M=S^1`$ is $`1`$-dimensional, the Clifford action reduces to multiplication by $`1`$ ($`\gamma ^\mu =\gamma ^1=1`$) and $`[D_E,a_{ij}]=i\dot{a_{ij}}`$. Therefore
$$i[๐,a]=\left(\begin{array}{cc}\dot{a_{11}}& \dot{a_{12}}+i\theta a_{12}\\ \dot{a_{21}}i\theta a_{21}& \dot{a_{22}}\end{array}\right)$$
(56)
is zero if and only if $`a_{11}=C`$, $`a_{22}=C^{}`$ are constant and $`a_{12}=\overline{a_{21}}=0`$ ($`\dot{a_{12}}=i\theta a_{12}`$ has no other $`2\pi `$-periodic solution than zero). Under these conditions
$$\xi _x(a)=|V_1|^2C+(1|V_1|^2)C^{}$$
differs from $`\zeta _y(a)`$ if and only if $`z_\zeta z_\xi `$. Hence, identifying $`a_{ij}`$ with $`\underset{n+\mathrm{}}{lim}(a_n)_{ij}`$ in Lemma III.2, one obtains that $`d(\xi _x,\zeta _y)`$ is infinite if and only if $`z_\xi z_\zeta `$, that is to say $`\zeta _y๐_\xi `$. $`\mathrm{}`$
By the proposition above the connected component $`๐_\xi `$ of $`d`$ contains, but is distinct from, the connected component $`\text{Acc}(\xi _x)`$ of $`d_H`$. This is enough to establish that the two metrics are not equal. Furthermore the results of the previous section strongly suggest that even on $`\text{Acc}(\xi _x)`$ the two metrics cannot coincide more than partially. To fix notation let us consider the distance $`d(\xi _x,\xi _\tau )`$ with $`\xi _\tau \text{Acc}(\xi _x)`$ given by (48) with $`\tau >0`$. On the one hand the function on $`\text{Acc}(\xi _x)`$
$$L(c(\tau ))d_H(\xi _x,c(\tau ))=\tau $$
(57)
is not $`2\pi `$-periodic, hence not in $`๐_E`$. Therefore it cannot be used as in (37) to realize the upper bound $`d_H`$ provided by Proposition III.1. Instead one could be tempted to use the geodesic distance on $`S^1`$,
$$F(\tau )d_{\text{geo}}(\xi _x,c_{}(\tau ))=\text{min}(\tau \text{ mod }[2\pi ],(2\pi \tau )\text{ mod }[2\pi ]),$$
(58)
but it may help in proving that $`d=d_H`$ only as long as as $`d_H`$ equals $`d_{\text{geo}}`$, that is to say as long as $`\tau \pi `$. Similarly $`L\text{ mod }[2\pi ]`$ could be efficient till $`\tau =2\pi `$ but it has infinite derivative at $`2k\pi `$ so it cannot be approximated by some $`a_n`$ satisfying the commutator norm condition. On the other hand for fixed $`k`$ the projection of the minimal horizontal curve between $`\xi _\tau ^k`$ and $`\xi _\tau `$ is a $`K`$-fold loop with
$$K=\{\begin{array}{ccc}|k|\hfill & \text{for}& \theta \text{ irrational}\hfill \\ \text{min}\{|k|,||k|q|\}\hfill & \text{for}& \theta =\frac{p}{q}\hfill \end{array}$$
where we assume that $`p`$ and $`q`$ are positive relatively prime with respect to each other and $`kp`$ is not a multiple of $`q`$ (otherwise $`\xi _\tau ^k`$ coincides with $`\xi _\tau `$). In any case when $`|k|=1`$ then $`K=1`$ and Proposition IV.4 should not forbid $`d(\xi _\tau ,\xi _\tau ^{\pm 1})`$ from equalling $`d_H(\xi _\tau ,\xi _\tau ^{\pm 1})=2\pi `$. We show below that this is indeed the case but only when $`R=1`$. On the contrary as soon as $`K>3`$ Proposition IV.4 certainly forbids $`d`$ from equalling $`d_H`$. In fact the situation is even more restrictive due to the particular choice (45) of the connection. Since the latter commutes with the diagonal part $`a_1`$ of any element $`a๐`$, $`\xi _\tau ^k(a_1)=\xi _\tau (a_1)`$ for any $`kK`$. Proposition IV.4 thus can be written as a system of $`K+1`$ equations
$`(\xi _\tau ^k\xi _\tau )(a_o)`$ $`=`$ $`2k\pi `$ (59)
$`\xi _\tau (a_o)`$ $`=`$ $`\xi _\tau (a_1)`$ (60)
where $`a_0=aa_1`$. (60) simply defines the diagonal part $`a_1`$ and one is finally left with $`K`$ equations (59) constraining the two real components of $`a_0`$. Therefore it is most likely that $`d`$ does not equal $`d_H`$ as soon as $`K>2`$.
To make these qualitative suggestions more precise, let us study the specific example of a โsea-levelโ (i.e. $`z_\xi =0`$) pure state $`\xi `$, assuming
$$|V_1|=|V_2|=\frac{1}{\sqrt{2}}.$$
(61)
All the distances on the associated connected component $`๐_\xi `$ can be explicitly computed. To do so it is convenient to isolate the part of the algebra that really enters the game in the computation of the distances. This is the objective of the following two lemmas. The first one is of algebraic nature: it deals with our explicit choice $`๐_I=M_2\left(\right)`$ and does not rely on the choice $`M=S^1`$.
###### Lemme V.2
Given $`\zeta _y`$ in $`๐_\xi `$, the search for the supremum in the computation of $`d(\xi _x,\zeta _y)`$ can be restricted to the set of elements
$$a=f๐+a_0$$
(62)
where $`๐`$ is the identity of $`M_2\left(\right)`$, $`f๐_E`$ vanishes at $`x`$ and is positive at $`y`$, while $`a_0`$ is an element of $`๐`$ whose diagonal terms are both zero and such that
$$\zeta _y(a_0)\xi _x(a_0)0.$$
(63)
Proof. Let $`\stackrel{~}{}`$ denote the operation that permutes the elements on the diagonal. By (56)
$$[๐,a]=\underset{\pm }{\mathrm{max}}\frac{(\dot{a_{11}}+\dot{a_{22}})\pm \sqrt{(\dot{a_{11}}\dot{a_{22}})^2+4|\dot{a_{12}}+i\theta a_{12}|^2}}{2}$$
(64)
is invariant under the permutation of $`a_{11}`$ and $`a_{22}`$. Thus $`[๐,a]=[๐,\stackrel{~}{a}]`$ so
$$[๐,\frac{\stackrel{~}{a}+a}{2}][๐,a].$$
(65)
Meanwhile
$$\xi _x(\frac{\stackrel{~}{a}+a}{2})=\xi _x(a)\text{ and }\zeta _y(\frac{\stackrel{~}{a}+a}{2})=\zeta _y(a)$$
(66)
therefore the supremum in the distance formula can be sought on
$$๐+\stackrel{~}{๐}=C^{\mathrm{}}(S^1)๐+๐_0$$
where $`๐_0`$ is the set of selfadjoint elements of $`๐`$ whose diagonal terms are zero. This fixes eq.(62). Now if $`a=f๐+a_0`$ attains the supremum then so does $`af(x)๐`$, hence the vanishing of $`f`$ at $`x`$. Moreover
$`[๐,f๐]=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)[๐,a]\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)`$ $``$ $`[๐,a]`$ (71)
$`[๐,a_0]\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)[๐,a]\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$ $``$ $`[๐,a],`$ (76)
so when $`a`$ satisfies the commutator norm condition so do $`f๐`$ and $`a_0`$. This implies that $`|\xi _x(a_0)\zeta _y(a_0)|`$ and $`|\xi _x(f๐)\zeta _y(f๐)|=|f(y)|`$ are smaller than
$$|\xi _x(a)\zeta _y(a)|=|f(y)+\zeta _y(a_0)\xi _x(a_0)|.$$
(77)
In particular, $`f(y)`$ and $`\zeta _y(a_0)\xi _x(a_0)`$ have the same sign, which we assume positive (if not, consider $`a`$ instead of $`a`$). $`\mathrm{}`$
Other simplifications come from the choice of $`S^1`$ as the base manifold. Especially the following lemma makes clear the role played by the functions $`L`$ and $`F`$ discussed in (57,58).
###### Lemme V.3
Let $`a=f๐+a_0`$ satisfy the commutator norm condition, then
$$\dot{f}1\text{ and }|f(\tau )|\dot{f}F(\tau )$$
(78)
where $`F(\tau )`$ is the $`2\pi `$-periodic function defined on $`[0,2\pi [`$ by
$$F(\tau )\text{min}(\tau ,2\pi \tau ).$$
(79)
Meanwhile
$$a_0=\left(\begin{array}{cc}0& ge^{i\theta L}\\ \overline{g}e^{i\theta L}& 0\end{array}\right)$$
(80)
where $`L(\tau )=\tau `$ for all $`\tau `$ in $``$ and $`g`$ is a smooth function on $``$ given by
$$g(\tau )=g(0)+_0^\tau \rho (u)e^{i\varphi (u)}๐u$$
(81)
with $`\rho C^{\mathrm{}}(,^+)`$, $`\rho 1`$, and $`\varphi C^{\mathrm{}}(,)`$ satisfying
$$\rho (u+2\pi )e^{i\varphi (u+2\pi )}=\rho (u)e^{i(\varphi (u)+2\theta \pi )}$$
(82)
while the integration constant is
$$g(0)=\frac{1}{e^{2i\theta \pi }1}_0^{2\pi }\rho (u)e^{i\varphi (u)}๐u.$$
(83)
Proof. (78) comes from the commutator norm condition (71) together with the $`2\pi `$-periodicity of $`f`$ (55), namely
$$f(\tau )=_0^\tau \dot{f}(u)๐u=_\tau ^{2\pi }\dot{f}(u)๐u.$$
The explicit form of $`a_0`$ is obtained by noting that any complex smooth function $`a_{12}๐_E`$ can be written $`ge^{i\theta L}`$ where $`ga_{12}e^{i\theta L}C^{\mathrm{}}()`$ satisfies
$$g(\tau +2\pi )=g(\tau )e^{2i\theta \pi }.$$
(84)
Hence any selfadjoint $`a_0`$ can be written as in (80), which yields for the commutator
$$[๐,a_0]=i\left(\begin{array}{cc}0& \dot{g}e^{i\theta L}\\ \dot{\overline{g}}e^{i\theta L}& 0\end{array}\right).$$
(85)
By (76) the commutator norm condition implies $`\dot{g}1`$, that is to say
$$g(\tau )=g(0)+_0^\tau \rho (u)e^{i\varphi (u)}๐u$$
(86)
where $`\rho C^{\mathrm{}}(,^+)`$, $`\rho 1`$, $`\varphi C^{\mathrm{}}(,)`$. The integration constant is fixed by (84),
$$g(0)=\frac{1}{e^{2i\theta \pi }1}\left(_0^{\tau +2\pi }\rho (u)e^{i\varphi (u)}๐u_0^\tau \rho (u)e^{i(\varphi (u)+2\theta \pi )}๐u\right),$$
(87)
and one extracts (82) from $`\frac{d}{d\tau }g(0)=0`$. Reinserted in (87) it finally yields (83). $`\mathrm{}`$
These lemmas yield the main result of the section: the computation of all distances on $`๐_\xi `$.
###### Proposition V.4
Let $`P`$ be the $`P^1`$ trivial bundle over the circle $`S^1`$ of radius one with connection (45). Let $`\xi _x`$ be a point in $`P`$ at altitude $`z_\xi =0`$ and $`๐_\xi `$ its connected component for the noncommutative geometry distance $`d`$. For any $`\zeta _y๐_\xi `$ there exists an equivalence class of real couples $`(\tau ,\theta ^{})(\tau +2\pi ,\theta ^{}2\theta \pi )`$ such that
$$\zeta _y=\left(\begin{array}{cc}1& 0\\ 0& e^{i\theta ^{}}\end{array}\right)\xi _\tau $$
(88)
where $`\xi _\tau `$ is given in (48,47). Without loss of generality one may assume that $`\tau `$ is positive (if not, permute the role played by $`\xi _x`$ and $`\zeta _y`$) so that
$$\tau =2k\pi +\tau _0$$
(89)
with $`k`$ and $`0\tau _02\pi `$. Then
$$d(\xi _x,\zeta _y)=\{\begin{array}{cc}\mathrm{max}(X;X+\tau _0Y)\hfill & \text{when }\tau _0\pi \hfill \\ \mathrm{max}(X;X+(2\pi \tau _0)Y)\hfill & \text{when }\pi \tau _0\hfill \end{array}$$
(90)
in which
$`X`$ $``$ $`RW_{k+1}\tau _0+RW_k(2\pi \tau _0)`$ (91)
$`Y`$ $``$ $`1RW_{k+1}RW_k`$ (92)
with $`R`$ defined in (50) and
$$W_k\frac{|\mathrm{sin}(k\theta \pi +\frac{\theta ^{}}{2})|}{|\mathrm{sin}\theta \pi |}$$
(93)
do not depend on the choice of the representantive of the equivalence class $`(\tau ,\theta ^{})`$.
Proof. The form (88) of $`\zeta _y`$ comes from the definition (54) of $`๐_\xi `$. It gives, for an element $`a`$ of Lemma V.3,
$$|\xi _x(a)\zeta _y(a)|=f(\tau )+\mathrm{}\left(Re^{i\theta _0}(g(\tau )e^{i\theta ^{}}g(0))\right)$$
(94)
where we use the definition (50) of $`\theta _0`$, the vanishing of $`f`$ at $`x`$, the positivity of $`f(y)=f(\tau )`$ as well as (63). The explicit form (81) of $`g`$ allows us to rewrite (94) as
$$f(\tau )+R_0^\tau \rho (u)\mathrm{cos}(\varphi ^{}(u))+\mathrm{}\left(Re^{i\theta _0}g(0)\left(e^{i\theta ^{}}1\right)\right)$$
(95)
where
$$\varphi ^{}(u)\varphi (u)\theta _0+\theta ^{}.$$
The point is to find the maximum of (95) on all the $`2\pi `$-periodic $`f`$ satisfying (78), the positive $`\rho `$, $`\rho 1`$ and the $`\varphi `$ satisfying (82). To do so we will first find an upper bound (eqs. (113) and (114) below) and prove that it is the lowest one.
Fixing a pure state $`\zeta _y`$ means fixing two values $`\theta ^{}`$ and $`\tau `$ or, equivalently by (89), fixing $`\theta ^{}`$, $`k`$ and $`\tau _0`$. The integral term in (95) then splits into
$$\mathrm{}_0^{2k\pi }\rho (u)e^{i\varphi ^{}(u)}๐u=\mathrm{}\left(\underset{n=0}{\overset{k1}{\mathrm{\Sigma }}}e^{2in\theta \pi }_0^{2\pi }\rho (u)e^{i\varphi ^{}(u)}๐u\right)$$
(96)
and
$$\mathrm{}_{2k\pi }^{2k\pi +\tau _0}\rho (u)e^{i(\varphi ^{}(u)})du=\mathrm{}(e^{2ik\theta \pi }_0^{\tau _0}\rho (u)e^{i\varphi ^{}(u)}du)$$
(97)
that recombine as
$$S_{k+1}_0^{\tau _0}\rho (u)\mathrm{cos}\varphi _k(u)๐u+S_k_{\tau _0}^{2\pi }\rho (u)\mathrm{cos}\varphi _{k_1}(u)๐u$$
(98)
where
$$S_k\frac{\mathrm{sin}k\theta \pi }{\mathrm{sin}\theta \pi }\text{ and }\varphi _k(u)\varphi ^{}(u)+k\theta \pi .$$
(99)
To compute the real-part term of (95) one uses the definition (83) of $`g(0)`$ and obtain
$$S_{\frac{1}{2}}_0^{2\pi }\rho (u)\mathrm{cos}\varphi _{\frac{1}{2}}(u)๐u$$
(100)
where
$$S_{\frac{1}{2}}\frac{\mathrm{sin}\theta ^{}/2}{\mathrm{sin}\theta \pi }\text{ and }\varphi _{1/2}(u)\varphi ^{}(u)\frac{\theta ^{}}{2}\theta \pi .$$
(95) is rewritten as
$$|\xi _x(a)\zeta _y(a)|=f(\tau )+R_0^{\tau _0}\rho (u)G_{k+1}(u)๐u+R_{\tau _0}^{2\pi }\rho (u)G_k(u)๐u$$
(101)
with
$$G_kS_k\mathrm{cos}\varphi _{k1}+S_{\frac{1}{2}}\mathrm{cos}\varphi _{\frac{1}{2}}.$$
(102)
The split of the integral makes the search for the lowest upper bound easier. Calling $`W_k`$ the maximum of $`|G_k(u)|`$ on $`[0,2\pi [`$, the positivity of $`\rho `$ makes (101) bounded by
$$f(\tau )+RW_{k+1}_0^{\tau _0}\rho (u)๐u+RW_k_{\tau _0}^{2\pi }\rho (u)๐u.$$
(103)
Now (64) with $`a_{11}=a_{22}=f`$ and $`|\dot{a_{21}}+i\theta a_{12}|=\rho `$ yields
$`|\dot{f}(u)+\rho (u)|1`$ whenever $`|\dot{f}(u)|0`$ (104)
$`|\dot{f}(u)\rho (u)|1`$ whenever $`|\dot{f}(u)|0`$ (105)
for any $`u`$, that is to say
$$\rho 1|\dot{f}|.$$
(106)
Therefore
$$_0^{\tau _0}\rho (u)๐u\tau _0_0^{\tau _0}|\dot{f}|.$$
(107)
Moreover $`f(\tau )=f(\tau _0)`$ ($`2\pi `$-periodicity of $`f`$) is positive by Lemma V.2 so
$$f(\tau _0)=|f(\tau _0)|_0^{\tau _0}|\dot{f}(u)|๐u.$$
(108)
Hence (107) gives
$$_0^{\tau _0}\rho (u)๐u\tau _0f(\tau _0).$$
(109)
Similarly
$$_{\tau _0}^{2\pi }|\dot{f}(u)|๐u|_{\tau _0}^{2\pi }\dot{f}(u)๐u|=|_0^{\tau _0}\dot{f}(u)๐u|=f(\tau _0)$$
hence
$$_{\tau _0}^{2\pi }\rho (u)๐u2\pi \tau _0f(\tau _0).$$
(110)
Back to (103), equations (109) and (110) yield the bound
$$f(\tau _0)Y+X$$
(111)
where $`X`$ is defined in (91) and $`Y`$ in (92). By (78) and in case
$$Y0,$$
(112)
(111) yields
$$|\xi _x(a)\zeta _y(a)|\{\begin{array}{cc}X+\tau _0Y\hfill & \text{for }0\tau _0\pi \hfill \\ X+(2\pi \tau _0)Y\hfill & \text{for }\pi \tau _02\pi \hfill \end{array}.$$
(113)
When $`Y0`$,
$$|\xi _x(a)\zeta _y(a)|X.$$
(114)
These are the announced lowest upper bounds. To convince ourselves let us build a sequence $`a_n`$ that realizes (113) or (114) at the limit $`n+\mathrm{}`$. As a preliminary step note that an easy calculation from (102) yields
$$G_k=A_k\mathrm{cos}\varphi ^{}+B_k\mathrm{sin}\varphi ^{}$$
where
$`A_k`$ $``$ $`S_{\frac{1}{2}}\mathrm{cos}({\displaystyle \frac{\theta ^{}}{2}}+\theta \pi )+S_k\mathrm{cos}(k1)\theta \pi `$ (115)
$`B_k`$ $``$ $`S_{\frac{1}{2}}\mathrm{sin}({\displaystyle \frac{\theta ^{}}{2}}+\theta \pi )S_k\mathrm{sin}(k1)\theta \pi .`$ (116)
$`G_k`$ attains its maximum value
$$W_k|A_k|\sqrt{1+\frac{|B_k|^2}{|A_k|^2}}=\frac{|\mathrm{sin}(k\theta \pi +\frac{\theta ^{}}{2})|}{|\mathrm{sin}(\theta \pi )|}$$
(117)
when<sup>\**</sup><sup>\**</sup>\**The ambiguity in the explicit form of $`\mathrm{\Phi }_k`$ is not relevant. Depending on the respective signs of $`A_k`$ and $`B_k`$, one choice yields $`W_k`$ whereas the other one yields $`W_K`$. What is important is the existence of a well defined value $`\mathrm{\Phi }_k`$ such that $`A_k\mathrm{cos}\mathrm{\Phi }_k+B_k\mathrm{sin}\mathrm{\Phi }_k=W_k`$.
$$\varphi ^{}=\mathrm{\Phi }_k\text{Arctan}\frac{B_k}{A_k}\text{ or }\text{Arctan}\frac{B_k}{A_k}+\pi .$$
(118)
Let then
$$a_n=\left(\begin{array}{cc}f_n& g_ne^{i\theta L}\\ \overline{g}_ne^{i\theta L}& f_n\end{array}\right)$$
be a sequence of elements of $`๐`$ that depend on the fixed value $`\tau =2k\pi +\tau _0`$ in the following way: in case (112) is fulfilled and $`\tau _0\pi `$, $`f_n`$ approximates from below the $`2\pi `$-periodic function
$$f_{}(t)=\{\begin{array}{cc}t\hfill & \text{ for }\mathrm{\hspace{0.17em}0}t\tau _0\hfill \\ \tau _0C(t\tau _0)\hfill & \text{ for }\tau _0t2\pi \hfill \end{array}$$
(119)
with
$$C\frac{\tau _0}{2\pi \tau _0}.$$
In case $`\tau _0\pi `$, $`f_n`$ approximates
$$f_+(t)=\{\begin{array}{cc}\frac{t}{C}\hfill & \text{ for }\mathrm{\hspace{0.17em}0}t\tau _0\hfill \\ 2\pi t\hfill & \text{ for }\tau _0t2\pi \hfill \end{array}.$$
(120)
When (112) is not fulfilled, $`f_n=f_0`$ is simply the zero function. In any case and whatever $`\tau _0`$, $`g_n`$ is defined via (81) and (83), replacing $`\varphi `$ with a sequence $`\varphi _n`$ approximating the step function $`\mathrm{\Phi }`$ of width $`2\pi `$ and height $`2\theta \pi `$ defined on $`[0,2\pi [`$ by
$$\mathrm{\Phi }(u)=\{\begin{array}{cc}\mathrm{\Phi }_{k+1}+\theta _0\theta ^{}\hfill & \text{ for }0u<\tau _0\hfill \\ \mathrm{\Phi }_k+\theta _0\theta ^{}\hfill & \text{ for }\tau _0<u<\pi \hfill \end{array},$$
(121)
and replacing $`\rho `$ with a sequence $`\rho _n`$ approximating the $`2\pi `$-periodic function
$$\mathrm{\Gamma }_I=1|\dot{f_I}|$$
(122)
where $`I=+,`$ or $`0`$. By construction the $`a_n`$โs satisfy the commutator norm condition. In particular the fact that $`\underset{n+\mathrm{}}{\text{lim}}\rho _n`$ and $`\mathrm{\Phi }`$ are step functions is not problematic since their derivatives are not constrained by the commutator. For technical details on how to approximate step functions by sequences of smooth functions, the reader is invited to consult classical textbooks such as \[?\]. The last point is to check that
$$\underset{n+\mathrm{}}{\text{lim}}|\xi _x(a_n)\zeta _y(a_n)|=(\text{113}).$$
(123)
This is a simple notation exercise: (99) gives
$`\varphi ^{}(u)`$ $`=`$ $`\mathrm{\Phi }_{k+1}\text{ for }\mathrm{\hspace{0.17em}0}u<\tau _0`$ (124)
$`\varphi ^{}(u)`$ $`=`$ $`\mathrm{\Phi }_k\text{ for }\tau _0<u<2\pi .`$ (125)
Therefore, by (101) together with (122),
$`\underset{n+\mathrm{}}{\text{lim}}|\xi _x(a_n)\zeta _y(a_n)|`$ $`=`$ $`f_I(\tau _0)+RW_{k+1}{\displaystyle _0^{\tau _0}}\mathrm{\Gamma }_I(u)๐u+RW_k{\displaystyle _{\tau _0}^{2\pi }}\mathrm{\Gamma }_I(u)๐u`$ (126)
$`=`$ $`f_I(\tau _0)RW_{k+1}{\displaystyle _0^{\tau _0}}|\dot{f_\pm }|๐uRW_k{\displaystyle _{\tau _0}^{2\pi }}|\dot{f_I}|๐u`$
$`+RW_{k+1}\tau _0+RW_k(2\pi \tau _0).`$
When $`(\text{112})`$ is fulfilled and $`\tau _0\pi `$, the subscript of $`f`$ is minus and (119) makes (126) equal to
$$\tau _0RW_{k+1}\tau _0RW_k(2\pi \tau _0)C+RW_{k+1}\tau _0+RW_k(2\pi \tau _0)$$
which is exactly the first line of (113). Similarly for $`\tau _0\pi `$, the subscript turns to $`+`$ and (120) yields for (126)
$$(2\pi \tau _0)RW_{k+1}\frac{\tau _0}{C}RW_k(2\pi \tau _0)+RW_{k+1}\tau _0+RW_k(2\pi \tau _0),$$
which is nothing but the second line of (113). Finally, when (112) is not fulfilled, $`f_I=\dot{f}_I=0`$ and (126) equals (114). $`\mathrm{}`$
Let us check the coherence of our result by noticing that for $`\tau _0=\pi `$ both formulas of (90) agree and yield
###### Check 1
$`d(\xi _x,\zeta _\pi )=\mathrm{max}(X,X+\pi Y)=\mathrm{max}(\pi R(W_{k+1}+W_k);\pi ).`$
Similarly for a given $`k`$ and $`\tau _0=2\pi `$, the second line of (90) agrees with the first line with $`k+1`$ and $`\tau _0=0`$, namely
###### Check 2
$`d(\xi _x,\zeta _{2k\pi +2\pi })=2\pi RW_{k+1}=d(\xi _x,\zeta _{2(k+1)\pi +0}).`$
This is nothing but the restriction of $`d`$ to the fiber over $`x`$. Its extreme simplicity (no โmaxโ is involved) indicates that the noncommutative metric is better understood fiberwise. We shall see in the next section that this is the main difference from the Carnot-Carathรฉodory metric. Another check, and certainly the best guarantee that Proposition V.4 is true, is to directly verify that formula (90) does define a metric: the vanishing of $`d`$ when $`\zeta _y=\xi _x`$ is obvious; the invariance under the exchange $`\xi _x\zeta _y`$ is not testable since the symmetry $`\tau \tau `$ is broken from the beginning by the specification that $`\tau `$ is positive. There remains the triangle inequality.
###### Check 3
For any $`\zeta _1,\zeta _2๐_\xi `$, $`d(\xi _x,\zeta _2)d(\xi _x,\zeta _1)+d(\zeta _1,\zeta _2).`$
Proof. Let $`\zeta _{\tau _i}`$, $`i=1,2`$, be two pure states defined by $`\tau _i=2\pi k_i+t_i`$ and $`\theta _i^{}`$, labeled in such a way that $`\tau _1\tau _2`$. The point is to check that
$$\mathrm{\Delta }d(\xi _x,\zeta _{\tau _1})+d(\zeta _{\tau _1},\zeta _{\tau _2})d(\xi _x,\zeta _{\tau _2})$$
(127)
is positive. Proposition V.4 is invariant under translation (i.e. a reparameterization of the circle $`\tau \tau +\text{constant}`$), which means that $`d(\zeta _{\tau _1},\zeta _{\tau _2})`$ is given by formula (90) with $`W_k`$ replaced by
$$W_{k_{12}}\frac{|\mathrm{sin}(k_{12}\theta \pi +\frac{\theta _2^{}\theta _1^{}}{2})|}{|\mathrm{sin}\theta \pi |}$$
and $`\tau _0`$ replaced by $`t_{12}`$. Here $`k_{12}`$ and $`t_{12}`$ are such that $`\tau _{12}\tau _2\tau _12k_{12}\pi +t_{12}`$. Explicitly
$`k_{12}=k_2k_1,`$ $`t_{12}=t_2t_1`$ $`\text{if }t_1t_2`$ (128)
$`k_{12}=k_2k_11,`$ $`t_{12}=2\pi +t_2t_1`$ $`\text{if }t_2t_1.`$ (129)
Let $`X_i,Y_i`$, $`i\{1,2,12\},`$ denote (91) and (92) in which $`k`$ is replaced by $`k_i`$. The only difficulty in checking that (127) is positive is the quite large number of possible expressions for $`\mathrm{\Delta }`$: one for each combination of the signs of the $`Y_i`$โs and $`t_i\pi `$. A simple way to reduce the number of cases under investigation is to decorate $`\mathrm{\Delta }`$ with three arrows indicating whether $`Y_1`$, $`Y_{12}`$ and $`Y_2`$ respectively are positive (upper arrow) or negative (lower arrow). For instance $`\mathrm{\Delta }_{}`$ denotes the value of $`\mathrm{\Delta }`$ when $`Y_10`$, $`Y_{12}0`$, and $`Y_20`$. Let us also use $`\stackrel{~}{\mathrm{\Delta }}`$ decorated with arrows to denote the formal expression (127) in which $`d(\xi _x,\xi _{\tau _1})`$, $`d(\zeta _{\tau _1},\zeta _{\tau _2})`$ and $`d(\xi _x,\zeta _{\tau _2})`$ are replaced either by $`X_i+t_i^mY_i`$ (upper arrow) or by $`X_i`$ (lower arrow). Here $`t_i^m\mathrm{min}(t_i,2\pi t_i)`$. For instance
$`\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`\stackrel{~}{\mathrm{\Delta }}_{}+t_1^mY_1`$ (130)
$`=`$ $`\stackrel{~}{\mathrm{\Delta }}_{}+t_2^mY_2`$ (131)
$`=`$ $`\stackrel{~}{\mathrm{\Delta }}_{}+t_1^mY_1+t_2^mY_2.`$ (132)
Now suppose that $`Y_1`$, $`Y_{12}`$, $`Y_2`$ are all positive, then
$$\mathrm{\Delta }=\mathrm{\Delta }_{}=\stackrel{~}{\mathrm{\Delta }}_{}\{\begin{array}{c}\stackrel{~}{\mathrm{\Delta }}_{}\\ \stackrel{~}{\mathrm{\Delta }}_{}\\ \stackrel{~}{\mathrm{\Delta }}_{}\end{array}.$$
Changing the sign of $`Y_10`$ and $`Y_{12}`$ yields
$$\mathrm{\Delta }=\mathrm{\Delta }_{}=\stackrel{~}{\mathrm{\Delta }}_{}\{\begin{array}{c}\stackrel{~}{\mathrm{\Delta }}_{}\\ \stackrel{~}{\mathrm{\Delta }}_{}\\ \stackrel{~}{\mathrm{\Delta }}_{}\end{array}.$$
Therefore, if one is able to show without using the sign of $`Y_1`$ or the sign of $`Y_{12}`$ that $`\stackrel{~}{\mathrm{\Delta }}_{}`$ is positive, one proves that both $`\mathrm{\Delta }_{}`$ and $`\mathrm{\Delta }_{}`$ are positive. In fact showing that one of the $`\stackrel{~}{\mathrm{\Delta }}_{}`$โs is positive is enough to prove that all the $`\mathrm{\Delta }_{}`$โs are positive (here $``$ means either $``$ or $``$). Of course the same is true with $`\stackrel{~}{\mathrm{\Delta }}_{}`$ so that, at the end, one just has to check the inequality of the triangle for one of the $`\stackrel{~}{\mathrm{\Delta }}_{}`$ and one of the $`\stackrel{~}{\mathrm{\Delta }}_{}`$.
Let us begin by $`\stackrel{~}{\mathrm{\Delta }}_{}`$, assuming first $`t_1t_2`$. With $`W_iW_{k_i}`$, $`W_{i+1}W_{k_i+1}`$, (128) yields
$`R^1\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`W_{1+1}t_1+W_1(2\pi t_1)+W_{12+1}t_{12}+W_{12}(2\pi t_{12})`$
$`W_{2+1}t_2W_2(2\pi t_2)`$
$`=`$ $`(2\pi t_2)(W_1+W_{12}W_2)+t_{12}(W_1+W_{12+1}W_{2+1})`$
$`+t_1(W_{1+1}+W_{12}W_{2+1})`$
which is positive since <sup>โ โ </sup><sup>โ โ </sup>โ โ this comes from $`|\mathrm{sin}(a+b)||\mathrm{sin}a|+|\mathrm{sin}b|`$ with $`a=(k_2k_1)\theta \pi +\theta _2^{}\theta _1^{}`$ and $`b=k_1\theta \pi +\theta _1^{}`$
$$W_{k_2}W_{k_1}+W_{k_{12}}$$
(133)
and similar equations for the other indices. Assuming now $`t_2t_1`$, (129) yields
$`R^1\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`t_2(W_{1+1}+W_{12+1}W_{2+1})+(2\pi t_1)(W_1+W_{12+1}W_2)`$
$`+(2\pi t_{12})(W_{12}+W_{1+1}W_2)`$
which is also positive by equations similar to (133) (be careful to use the definition (129) of $`k_{12}`$ and no longer definition (128)). Thus, whatever $`t_1`$ and $`t_2`$, $`\stackrel{~}{\mathrm{\Delta }}_{}`$ is positive and the triangle inequality is checked for all the configurations $``$ of the $`Y_i`$โs.
Things are slightly more complicated for the configurations $``$ for one also has to deal with the signs of $`t_i\pi `$. First assume $`t_1t_2`$:
* $`t_1t_2\pi `$ (implies $`t_{12}\pi `$),
$`(2R)^1\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`W_1(\pi t_1)+W_{12}(\pi t_{12})W_2(\pi t_2)`$
$``$ $`(\pi t_2)(W_1+W_{12}W_2).`$
* $`\pi t_1t_2`$ (implies $`t_{12}\pi `$),
$`\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`2R(W_{1+1}(t_1\pi )+W_{12}(\pi t_{12})W_{2+1}(t_2\pi ))+2(t_2t_1)`$
$``$ $`2R(t_1\pi )(W_{1+1}+W_{12}W_{2+1})+2(t_2t_1)(1RW_{2+1}).`$
* $`t_1\pi t_2`$ and $`t_{12}\pi `$,
$$\stackrel{~}{\mathrm{\Delta }}_{}=2R(W_1(\pi t_1)+W_{12}(\pi t_{12}))+2(t_2\pi )(1RW_{2+1}).$$
* $`t_1\pi t_2`$ and $`t_{12}\pi `$,
$`\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`2R(W_1(\pi t_1)+W_{12+1}(t_{12}\pi )W_{2+1}(t_{12}\pi ))+t_1(12RW_{2+1})`$
$``$ $`2R(t_{12}\pi )(W_1+W_{12+1}W_{2+1})+2t_1(1RW_{2+1}).`$
These five expressions are positive by (133) and the positivity of $`Y_2`$. Similarly, in case $`t_2t_1`$:
* $`t_2t_1\pi `$ (implies $`t_{12}\pi `$),
$`\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`2R(W_1(\pi t_1)+W_{12+1}(t_{12}\pi )W_2(\pi t_1))+2(t_1t_2)(1RW_2)`$
$``$ $`2R(\pi t_1)\left(W_1+W_{12+1}W_2\right)+2(t_1t_2)\left(1RW_2\right).`$
* $`\pi t_2t_1`$ (implies $`t_{12}\pi `$),
$`(2R)^1\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`W_{1+1}(t_1\pi )+W_{12+1}(t_{12}\pi )W_{2+1}(t_2\pi )`$
$``$ $`(t_2\pi )(W_{1+1}+W_{12+1}W_{2+1}).`$
* $`t_2\pi t_1`$ and $`t_{12}\pi `$,
$`\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`2R(W_{1+1}(t_1\pi )+W_{12}(\pi t_{12})W_2(\pi t_{12}))+2(2\pi t_1)(1RW_2)`$
$``$ $`2R(\pi t_{12})(W_{1+1}+W_{12}W_2)+2(2\pi t_1)(1RW_2).`$
* $`t_2\pi t_1`$ and $`t_{12}\pi `$,
$`\stackrel{~}{\mathrm{\Delta }}_{}`$ $`=`$ $`2RW_{1+1}(t_1\pi )+2RW_{12+1}(t_{12}\pi )+2(\pi t_2)(1RW_2).`$
$`\mathrm{}`$
The proof above is long but we believe it is important to convince oneself that formula V.4 does define a metric, which is not obvious at first sight. As a final test, let us come back to the beginning of this section and verify Lemma III.1.
###### Check 4
$`d(\xi _x,\zeta _y)d_H(\xi _x,\zeta _y)\text{ for any }\zeta _y\text{Acc}(\xi _x).`$
Proof. Let $`\zeta _y=\xi _\tau `$. Then $`d_H(\xi _x,\xi _\tau )=2k\pi +\tau _0`$ so that
$$d(\xi _x,\xi _\tau )d_H(\xi _x,\xi _\tau )\{\begin{array}{ccc}=\hfill & 2RW_k(\pi \tau _0)2k\pi \hfill & \\ & 2\pi (W_kk)\hfill & \text{ when }Y0,\tau _0\pi ,\hfill \\ =\hfill & 2RW_{k+1}(\tau _0\pi )2(\tau _0\pi )2k\pi \hfill & \\ & 2(\tau _0\pi )W_k2k\pi \hfill & \text{ when }Y0,\tau _0\pi ,\hfill \\ =\hfill & \tau _0(RW_{k+1}RW_k1)\hfill & \\ & +2\pi (RW_kk)\hfill & \text{ when }Y0.\hfill \end{array}$$
(134)
These three expressions are negative by (133) and <sup>โกโก</sup><sup>โกโก</sup>โกโกobvious for $`k1`$, then by induction $`|\mathrm{sin}k\theta \pi |k|\mathrm{sin}\theta \pi |.`$ $`\mathrm{}`$
## VI Interpretation: a smooth cardio-torus
This section aims at analyzing the result of Proposition V.4. We first compare $`d`$ to $`d_H`$ on $`\text{Acc}(\xi _x)`$ (corollaries VI.1 and VI.2), then study the restriction of $`d`$ to the fiber over $`x`$ and to the base $`M=S^1`$. The reader may wonder why we do not systematically replace $`R`$ by its value $`1`$. The point is that for two states on the same fiber ($`y=x`$) the diagonal part of $`a`$ does not play any role so that Proposition V.4 is valid also for non vanishing $`z_\xi `$. Also, for $`yx`$ some calculations<sup>?</sup> show that V.4 is still valid for non-zero $`z_\xi `$ as long as $`2V_i^2R(W_{k+1}+W_k)`$ is positive for both $`i=1,2`$. This is the reason why, in the following discussion, we keep writing $`R`$.
### VI.1 The shape of $`๐_\xi `$
Taking $`\zeta _y`$ in $`\text{Acc}(\xi _x)`$ amounts to setting $`\theta ^{}=0`$. $`W_k`$ is replaced by
$$S_k\frac{|\mathrm{sin}k\theta \pi |}{|\mathrm{sin}\theta \pi |}$$
and proposition V.4 is rewritten in a somehow more readable fashion.
###### Corollary VI.1
Let $`\zeta _y=\xi _\tau \text{Acc}(\xi _x)`$, with $`\tau =2k\pi +\tau _0.`$ For $`k`$ such that $`S_{k+1}+S_k\frac{1}{R}`$,
$$d(\xi _x,\xi _\tau )=\{\begin{array}{cc}2RS_k(\pi \tau _0)+\tau _0\hfill & \text{when }\tau _0\pi \hfill \\ 2RS_{k+1}(\tau _0\pi )+2\pi \tau _0\hfill & \text{when }\pi \tau _0\hfill \end{array}.$$
(135)
For $`k`$ such that $`S_{k+1}+S_k\frac{1}{R}`$,
$$d(\xi _x,\xi _\tau )=RS_{k+1}\tau _0+RS_k(2\pi \tau _0).$$
It is easy to see on which part of $`P`$ the noncommutative geometry metric and the Carnot-Carathรฉodory one coincide.
###### Corollary VI.2
For any $`R`$, $`d(\xi _x,\xi _\tau )=d_H(\xi _x,\xi _\tau )`$ for $`\tau [0,\pi ]`$. Moreover if $`R=1`$ the two metrics are also equal for $`\tau [\pi ,2\pi ]`$. These are the only situations in which $`d=d_H`$.
Proof. $`S_0=0`$, $`S_1=1`$ and by construction $`R1`$. Therefore for $`k=0`$, $`S_{k+1}+S_k=1\frac{1}{R}`$ so
$$d(\xi _x,\xi _\tau )=\{\begin{array}{cc}\tau _0=d_H(\xi _x,\xi _{\tau _0})\hfill & \text{when }\tau _0\pi \hfill \\ 2\pi (1R)+\tau _0(2R1)\hfill & \text{when }\pi \tau _0\hfill \end{array}$$
(136)
which yields the equality of $`d`$ and $`d_H`$ for the indicated values of $`\tau `$ and $`R`$. From check 4 in the preceding section, $`d`$ may equal $`d_H`$ only if $`S_k=k`$, i.e. $`k=0`$ or $`1`$. When $`k=1`$, $`S_k+S_{k+1}1`$ and the last line of (134) gives the difference $`\delta `$ between $`d`$ and $`d_H`$,
$$\delta =\tau _0(RS_2R1)+2\pi (R1).$$
(137)
$`S_22`$ so $`\delta (R1)(2\pi +\tau )`$. $`\delta `$ may vanish only if $`R=1`$ and, going back to (137), only if $`\tau _0=0`$. $`\mathrm{}`$
This result is more restrictive that what was expected from Proposition IV.4 revisited in (59), namely that $`d`$ may equal $`d_H`$ as long as $`c`$ does not have sequences of more than $`2`$ self-intersecting points, i.e. up to $`\tau _{\text{max}}=4\pi +\tau _0`$. It seems that Proposition IV.4 alone is not sufficient to show that $`\tau _{\text{max}}2\pi `$. At best one can obtain
$$\tau _{\text{max}}<4\pi .$$
(138)
Although (138) is not in se an interesting result but simply a weaker formulation of Corollary VI.2, we believe it is interesting to see how far Proposition IV.4 can lead. This could be the starting point for a generalization of the results of this paper to manifolds other than $`S^1`$. Let $`G,\overline{G}`$ be the off-diagonal components of $`a`$. (59) is rewritten as
$$\mathrm{}\left(G(\tau )e^{i(\theta \tau \theta _0)}e^{ik\theta \pi }\right)=\frac{k\pi }{R\mathrm{sin}k\theta \pi }$$
(139)
for any $`k=1,\mathrm{},K`$. For $`K=2`$ this system has a unique solution
$$G(\tau )=Ce^{i\theta \tau }e^{i(\theta _0\frac{\pi }{2})}$$
(140)
where
$$C\frac{2\pi }{R\mathrm{sin}2\theta \pi }$$
is a constant. Therefore $`\xi _\tau (a_0)=\mathrm{}\left(e^{i(\theta \tau \theta _0)}G(\tau )\right)=0`$ so that, by (60), $`\xi _\tau (a)=0`$. By Proposition IV.4 this is possible only for $`\tau =0`$. Hence there cannot be more than one sequence of 2 self-intersecting points, hence (138).
In any case, when $`\tau `$ is greater than $`2\pi `$, $`d`$ strongly differs from $`d_H`$. While the latter is unbounded, the former is bounded,
$$d(\xi _x,\zeta _y)\text{max}(\frac{2\pi R}{|\mathrm{sin}\theta \pi |},\pi ).$$
As illustrated in figure 4, $`\text{Acc}(\xi _x)`$ viewed as a $`1`$-dimensional object looks like a straight line when it is equipped with $`d_H`$, whereas it looks rather chaotic when it is equipped with $`d`$.
### VI.2 The shape of the fiber
From a fiberwise point of view the situation drastically changes. Parameterizing the fiber $`S_x`$ over $`x`$ by
$$\varphi 2k\theta \pi +\theta ^{}\text{ mod }[2\pi ],$$
one obtains a very simple expression for the noncommutative distance,
$$d(0,\varphi )=\frac{2\pi R}{|\mathrm{sin}\theta \pi |}\mathrm{sin}\frac{\varphi }{2}.$$
(141)
For those points of $`S_x`$ which are accessible from $`\xi _x`$, namely for $`\theta ^{}=0`$, the Carnot-Carathรฉodory metric is
$$d_H(0,\varphi )=2k\pi .$$
Hence, when $`\theta `$ is irrational and in any neighborhood of $`\xi _x=0`$ in the Euclidean topology of $`S_x`$, it is always possible to find some
$$\varphi _k\xi _0^k=2k\theta \pi \text{ mod }[2\pi ]$$
which are arbitrarily Carnot-Carathรฉodory-far from $`\xi _x`$. In other terms $`d_H`$ destroys the $`S^1`$ structure of the fiber. On the contrary $`d`$ keeps it in mind in a rather intriguing way. Let us compare $`d`$ to the Euclidean distance $`d_E`$ on the circle of radius
$$\frac{2R}{|\mathrm{sin}\theta \pi |}.$$
(142)
At the cut-locus $`\varphi =\pi `$, the two distances are equal but whereas $`d_E(0,.)`$ is not smooth, the noncommutative geometry distance is smooth (cf Figure 5).
In this sense, if we imagine an observer localized at $`\xi _x`$ and whose only information about the geometry of the surrounding world is the measurement of the function $`d(0,\varphi )`$, $`S_x`$ looks โsmoother than a circleโ. More rigourously, (141) turns out to be the length $`L(\varphi )`$ of the minimal arc joining the origin to a point $`\varphi `$ on the cardioid with polar equation
$$r=\frac{\pi }{4}(1+\mathrm{cos}\phi ).$$
(143)
Indeed restricting to $`0\varphi \pi `$ (since $`L(\varphi )=L(2\pi \varphi )`$),
$$L(\varphi )=_0^\varphi \sqrt{r^2+(\frac{dr}{d\phi })^2}๐\phi =_0^\varphi \frac{\pi }{2}\mathrm{cos}\frac{\phi }{2}d\phi =\pi \mathrm{sin}\frac{\varphi }{2}=d(0,\varphi ).$$
(144)
One has to be careful with the interpretation of equation (144). The noncommutative geometry distance does not turn the loop $`S_x`$ into a cardioid. What the noncommutative metric does is to turn $`S_x`$ into an object that looks like a cardioid for an observer localized at $`x`$ who is measuring the distance between him and a point of $`S_x`$. Corollary VI.1 being invariant under a re-parameterization of the basis $`S^1`$ ($`\tau \tau +\text{const.}`$), the same analysis is true for an observer localized at $`yx`$. In this sense the cardioid point of view is an intrinsic point of view.
Things are clearer in analogy with the circle (Figure 7): consider $`2`$ observers $`๐ช_i`$, $`i=1,2`$, located at distinct points $`\varphi _i`$ on a loop $`S`$. Assume each of them measures its own distance function
$$d_i:zSd_i(x_i,z).$$
If both find that $`d_i=d_E`$, then they will agree that $`S`$ is a circle. On the contrary if both find that $`d_i=d`$, then each of them will pretend to be localized at the point opposite to the cut locus of the cardioid and they will disagree on the nature of $`S`$. In fact their disagreement is only due to their belief that $`S`$ is a manifold. What the present work shows is precisely that the loop $`S_x`$ equipped with the noncommutative metric $`d`$ is not a manifold. This example nicely illustrates how the distance formula (3) allows one to define on very simple objects (like tori) a metric which is not accessible from classical differential geometry.
### VI.3 The shape of the basis
From an intrinsic point of view the fiber looks like a cardioid. What does the base $`M=S^1`$ look like ? Let $`S_\xi `$ denote the set of points of $`๐_\xi `$ corresponding to the same vector $`\xi P^{n1}`$,
$$S_\xi \{pP,V(p)=\xi \}.$$
We parameterize $`S_\xi `$ by $`\phi [0,2\pi [`$ with $`\xi _x=0`$. Any point in $`S_\xi `$ can be obtained as a $`\zeta _\tau `$ where $`\tau =2k\pi +\phi `$ and $`\zeta `$ defined by (88) with
$$\theta ^{}=\theta \tau .$$
(145)
In order to compute $`d_H`$, note that $`\zeta _\tau `$ is accessible from $`\xi _x`$ if and only if $`\zeta _0`$ is accessible, that is to say iff $`\theta ^{}=2k^{}\theta \pi \text{ mod }[2\pi ]`$ for some integer $`k^{}`$. In other words $`\text{Acc}(\xi _x)S_\xi `$ is the subset of $`[0,2\pi [`$ given by the numbers $`\phi `$ that write
$$\phi =2p\theta ^1\pi +2p^{}\pi $$
for some integers $`p,p^{}`$. When $`\theta `$ is irrational $`\text{Acc}(\xi _x)S_\xi `$ is dense in $`S_\xi `$ and to a given $`\phi `$ corresponds one and only one couple of integers $`p,p^{}`$. By $`(\text{145})`$ one obtains
$$\zeta _0=\xi _0^{(k+p^{})},$$
where we used the notation (51). Hence $`\zeta _\tau =\xi _{2p\theta ^1\pi }`$, so that
$$d_H(0,\phi )d_H(\xi _x,\zeta _\tau )=2p\theta ^1\pi .$$
As in the case of the fiber $`S_x`$, one finds close to $`0S_\xi `$ in the Euclidean topology some points that are infinitely Carnot-Carathรฉodory far from $`0`$. Hence $`d_H`$ not only forgets the shape of the fiber but also the shape of the base.
On the contrary the noncommutative distance $`d`$ is finite on $`S_\xi `$ and preserves the shape of the base, although the latter is deformed in a slightly more complicated way than the fiber. Note that, via (145),
$$W_k=\frac{|\mathrm{sin}(\frac{\theta }{2}\phi )|}{|\mathrm{sin}\theta \pi |},W_{k+1}=\frac{|\mathrm{sin}(\frac{\theta }{2}(2\pi \phi ))|}{|\mathrm{sin}\theta \pi |}$$
are independent of $`k`$. The same is true for $`X`$ and $`Y`$ so that $`d(0,\varphi )=d(\xi _x,\zeta _\tau )`$ only depends on $`\phi `$ as expected. Explicitly, defining $`\lambda \frac{\phi }{2\pi }`$, Proposition V.4 writes
$$d(0,\phi )=\pi \left(\lambda \mathrm{sin}(\theta \pi (1\lambda ))+(1\lambda )\mathrm{sin}(\theta \pi \lambda )\right)$$
(146)
when $`Y`$ is negative and
$$d(0,\phi )=\{\begin{array}{cc}2\pi \left((\frac{1}{2}\lambda )\mathrm{sin}\theta \pi \lambda +\lambda \right)\hfill & \text{when }\lambda \frac{1}{2}\hfill \\ 2\pi \left((\lambda \frac{1}{2})\mathrm{sin}\theta \pi (1\lambda )+1\lambda \right)\hfill & \text{when }\lambda \frac{1}{2}.\hfill \end{array}$$
(147)
when $`Y`$ is positive. Even for a fixed value of $`R`$, $`Y`$ may change sign when $`\phi `$ runs from $`0`$ to $`2\pi `$ so it seems difficult to find for $`S_\xi `$ a picture like the cardioid for $`S_x`$. However, assuming that $`Y`$ is always negative, one can view the first line of (146) as a kind of convex deformation of a cardioid. In particular when $`\theta 1`$ or $`\theta 0`$, $`Y`$ is indeed negative for any $`\phi `$ so that
$$\underset{\theta 1}{lim}d(0,\phi )=\pi \mathrm{sin}\frac{\phi }{2}$$
which corresponds to the length on a cardioid of infinite radius (since $`\underset{\theta 1}{lim}=+\mathrm{}`$), while
$$\underset{\theta 0}{lim}d(0,\varphi )=2R\phi (1\frac{\phi }{2\pi }).$$
This is the arc length of the curve $`r(\phi )`$, solution of
$$r^2+\dot{r}^2=(1\frac{\phi }{\pi })^2.$$
(148)
(148) has no global solution. Gluing the solution of $`\dot{r}=\sqrt{(1\frac{\phi }{\pi })^2r^2}`$ on $`[\pi ,2\pi ]`$ with the solution of $`\dot{r}=\sqrt{(1\frac{\phi }{\pi })^2r^2}`$ on $`[0,\pi ]`$ with initial condition $`r(\pi )=0`$, one obtains that at the limit $`\theta 1`$ the base $`S_\xi `$, seen for $`\xi `$, has the shape of a heart (figure 8). Hence, still from the intrinsic point of view developed from $`S_x`$, $`\theta `$ is a deformation parameter for the base of $`P`$ from an infinite cardioid to a heart. The shape of $`S_\xi `$ for intermediate values of $`\theta `$ is deserving of further study.
## VII Conclusion and outlook
The $`2`$-torus $`๐_\xi `$ inherits from noncommutative geometry a metric smoother than the Euclidean one (the associated distance function is smooth at the cut locus). It gives to both the fiber and the base the shape of a cardioid or a heart. Such a โsmooth cardio-torusโ (shall we denote it $`\mathrm{}_\xi `$ ?) offers a concrete example in which the distance (3) is โtrulyโ noncommutative, in the sense that is not a Riemannian geodesic distance (as in the commutative case), nor a combination of the latter with a discrete space (as in the two-sheet model), not even the Carnot-Carathรฉodory one. The noncommutative distance combines some aspects of the Euclidean metric on the torus (preservation of the fiber structure) with some aspects of the Sub-Riemannian metric (dependance on the connection).
From a geometrical point of view several questions remain to be studied: what is the metric when both the scalar and the gauge fluctuations are non-zero ? How to extend the present result to manifolds other than $`S^1`$ ? In particular it could be interesting to separate in the holonomy conditions the role of the curvature from the role of the non-connectedness. For instance could it be that, in a certain โlocalโ sense, $`d`$ equals $`d_H`$ ? Let us also underline that the present work is intended to be the first step in the computation of the metric aspect of the noncommutative torus where the bundle of pure states $`P`$ is no longer trivial.
From a physics point of view, it would be interesting to reexamine in the light of the present results some interpretations that were given to sub-Riemannian-geodesics as effective trajectories of particles (Wongโs equations). This should be the object of further work.
Acknowledgments A preliminary version of Proposition III.1 was established by T. Krajewski. B. Iochum suggested to study the example $`M=S^1`$, and pointed out that the conditions of Corollary III.4 were not equivalent to the holonomy being trivial. Thanks to all the NCG group of CPT and IML for numbers of useful discussions, and to the director for hosting. Warm thanks to P. Almeida and R. Montgomery for illuminating remarks. Many thanks to A. Greenspoon for careful reading of the manuscript.
Work partially supported by EU network geometric analysis. CPT is the UMR 6207 from CNRS and Universities de Provence, Mรฉditerranรฉe and Sud Toulon-Var, affiliated to the FRUMAM. |
warning/0506/hep-th0506048.html | ar5iv | text | # 1 Introduction
## 1 Introduction
It has been conjectured by Peter West that the still elusive M-Theory possesses a rank eleven Kac-Moody symmetry algebra, called $`E_{11}`$, that is the triple extended or very extended $`E_8`$ algebra. Very extended algebras can be defined for any finite-dimensional Lie algebra $`๐ข`$ . So it is tempting to argue that other theories, which are associated with other triple extensions of Lie algebras, may exist. Indeed the same analysis was applied to a conjectured extension of the eleven dimensional supergravity , . More generally, it has been proposed that the closed bosonic string in D dimensions and type I supergravity and pure gravity theories exhibit a Kac-Moody symmetry algebra, respectively identified as the triple extensions of the $`D`$ and $`A`$ series , , . This conjecture is supported by dimensional reduction and by the so called cosmological billiards , , . Then it is natural to look for a more general symmetry algebra which can include all these Kac-Moody algebras as particular cases. So we address the question of how to go beyond $`๐ข^{+++}`$ algebras; we find that the adjoint of a new simple root $`\alpha _{+4}`$ introduces multiple links and loops in the structure of the 4-extended algebra, if $`\alpha _{+4}`$ is an ordinary Kac-Moody simple root, while the โsimple-linksโ structure is preserved if we allow $`\alpha _{+4}`$ is a Borcherds (imaginary) simple root.
The (first) extension of a finite-dimensional Lie algebra is the construction of (untwisted) affine Kac-Moody algebras, which are obtained adding to the simple roots of any finite-dimensional Lie algebra $`๐ข`$ a root $`\alpha _0`$ that is the opposite of the highest root (h.r.) plus a light-like vector $`k_+`$, in order to make $`\alpha _0`$ linearly independent from the system of the simple roots of $`๐ข`$, keeping unchanged its length, and are denoted as $`๐ข^+`$ or $`\widehat{๐ข}`$ or $`๐ข^1`$. This procedure for the simply laced algebras of the $`D_N`$ and $`E_N`$ series can be formulated as the addition to the simple root system of $`๐ข`$ of another root $`\alpha _0`$, that is the opposite of the unique fundamental weight of length 2, which is in the root lattice of the algebra, plus a light-like vector $`k_+`$. The light-like vector can be considered to belong to a 2-dim. Lorentzian lattice, usually denoted II<sup>1,1</sup>, and the double extension or overextension of $`๐ข`$, denoted $`๐ข^{++}`$, is obtained by adding a new simple root, of length 2, which is formed by the sum of the two light like vectors $`k_\pm `$, $`(k_+,k_{})=1`$, spanning II<sup>1,1</sup>. The triple extended or very extended $`๐ข`$, denoted $`๐ข^{+++}`$, is obtained adding a new simple root of length 2, which belongs to a new copy of the lattice II<sup>1,1</sup>, plus $`k_+`$. In this way, an indefinite Kac-Moody or Lorentzian algebra of rank $`r+3`$ is obtained whose roots belong to a Lorentzian lattice of dimension $`r+4`$, so it is natural to wonder if an indefinite Kac-Moody algebra of rank $`r+4`$ can be obtained by a further extension. Moreover, let us notice that in the lattice II<sup>1,1</sup> vectors of negative length do exist (see Appendix A). From this remark one can make an extension of $`๐ข`$ adding a new root, that is the opposite of any fundamental weight, that can be written as linear combination with integer coefficients of the simple roots of the algebra, plus a suitable element of the lattice II<sup>1,1</sup> in order to have an independent new simple root of length 2. This construction will be discussed below. It has been pointed out in that the structure of subalgebras of hyperbolic Kac-Moody, in general of $`2`$-extended (overextended) Lie algebras, is very rich and surprising. Some of the results of that paper can be generalized to more general extensions and we comment on this point below. This paper is organized as it follows. In Sec. 2, to make the paper self contained, we recall the well-known construction of the 3-extended Lie algebras. We show that the 4-extended algebras are described by Dynkin-Kac diagrams with loops and multiple links (so their structure is quite different from that of the 1-, 2- and 3-extended algebras). In particular, we show that we can not have a situation in which the new (fourth) simple root $`\alpha _{+4}`$ is simply linked to $`\alpha _{+3}`$, unless we let $`\alpha _{+4}`$ be a Borcherds simple root (with squared norm zero or negative). So we also study the possibility to have a Borcherds extension of the $`๐ข^{+++}`$ algebras, but this extension has sense only when the algebra $`๐ข`$ is simply-laced (see Appendix B for the definition of a Borcherds algebra ). In Sec. 3, we show that all the non-simply laced 3-extended Lie algebras can be obtained by *folding* the simply-laced ones. In Sec. 4, we discuss non standard extension procedures, discussing in detail a few examples which may be of physical interest. In Sec. 5, we show that the algebras obtained by some general non standard procedure, but not for all the procedures, are indeed subalgebras, of the same rank, of the standard triple extended algebras. In particular we prove that any non standard $`1`$-extension of $`E_9`$, with a root simply linked to a simple roots of $`E_9`$, is a subalgebra of $`E_{10}`$. The simple root systems of a set of rank 11 subalgebras of $`E_{11}`$, containing as subalgebra $`E_{10}`$, are explicitly written. Finally we present a few conclusions and perspectives. To make the paper self-contained two very short Appendices are added to recall the main features of the $`2`$-dim Lorentzian lattice II<sup>1,1</sup> and of Borcherds algebras.
## 2 On extensions of $`๐ข^{+++}`$ algebras
An excellent discussion of the mechanism of standard extensions of Lie algebras can be found in , here we briefly recall the essential points, mainly to introduce the notation and to make the paper self-consistent <sup>1</sup><sup>1</sup>1Let us remember that the standard extension for G<sub>2</sub> does not work, because G<sub>2</sub> is the only finite Lie algebra for which it is not possible to normalize the highest root $`\theta `$ such that $`(\theta ,\theta )=2`$. However, itโs possible to extend G<sub>2</sub> with the following choice of the extended roots: $`\alpha _{+1}=k_+\theta `$, $`\alpha _{+2}=(k_++3k_{})`$ and $`\alpha _{+3}=(l_++3l_{})+k_+`$. Anyway, G$`{}_{}{}^{+++}{}_{2}{}^{}`$ can be obtained by folding D$`{}_{}{}^{+++}{}_{4}{}^{}`$ (see Section 3).. Let $`๐ข`$ a simple Lie algebra of rank $`r`$, with simple root system $`\alpha _i`$ ($`i=1,\mathrm{},r`$) and root lattice $`\mathrm{\Lambda }_๐ข=_{i=1}^r\alpha _i`$ (for the roots and fundamental weights we use the notation of ). Let us consider also two copies of the lattice II<sup>1,1</sup>, which we indicate with II$`{}_{k_\pm }{}^{}{}_{}{}^{1,1}`$ and II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$. In the following, we consider indefinite Kac-Moody algebras with root lattice included in the direct sum $`\mathrm{\Lambda }_๐ข`$ II$`{}_{k_\pm }{}^{1,1}`$ II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$.
Let $`๐ข^+`$ the extended Lie algebra (or affine Kac-Moody algebra), with simple root system {$`\alpha _i`$, $`\alpha _0\alpha _{r+1}\alpha _{+1}=h.r.+k_+\}`$ ($`i=1,\mathrm{},r`$), where $`h.r.`$ denotes the highest root of $`๐ข`$ (which is $`\theta h.r.=_{i=1}^ra_i\alpha _i`$ where $`a_i`$ are Kac marks, see ), and
$$(k_+,k_+)=(k_+,\alpha _i)=0.$$
(1)
Let $`๐ข^{++}`$ the double extended or overextended Lie algebra (actually a Lorentzian Kac-Moody algebra), with simple root system {$`\alpha _j`$, $`\alpha _{r+2}\alpha _{+2}=(k_++k_{})\}`$ ($`j=1,\mathrm{},r+1`$), where
$$(k_+,k_{})=1,(k_{},k_{})=(k_{},\alpha _i)=0.$$
(2)
Letโs note that the root lattice of $`๐ข^+`$ is properly contained in the direct sum $`\mathrm{\Lambda }_๐ข`$ II$`{}_{k_\pm }{}^{}{}_{}{}^{1,1}`$, whereas the root lattice of $`๐ข^{++}`$ coincides with the direct sum. To see this, it is enough to obtain $`k_+`$ and $`k_{}`$ from the other roots $`\alpha _j`$ ($`j=1,\mathrm{},r+2`$). In fact,
$$k_+=\theta +\alpha _{+1}=\underset{i=1}{\overset{r}{}}a_i\alpha _i+\alpha _{+1}$$
(3)
$$k_{}=k_+\alpha _{+2}=\theta \alpha _{+1}\alpha _{+2}$$
(4)
so $`k_+`$ and $`k_{}`$ are linear combinations (with integer coefficients) of the $`r+2`$ simple roots of the $`๐ข^{++}`$ algebra. This means that with a change of basis (in the lattice), we can pass from $`\{\alpha _j(j=1,\mathrm{},r+2)\}`$ to $`\{\alpha _i(i=1,\mathrm{},r);k_+,k_{}\}`$, that is the lattice $`\mathrm{\Lambda }_๐ข\alpha _{+1}\alpha _{+2}`$ coincides with $`\mathrm{\Lambda }_๐ข`$ II$`{}_{k_\pm }{}^{}{}_{}{}^{1,1}`$.
Now, let $`๐ข^{+++}`$ the triple extended Lie algebra (still a Lorentzian Kac-Moody algebra), with simple root system {$`\alpha _k`$, $`\alpha _{r+3}\alpha _{+3}=(l_++l_{})+k_+`$} ($`k=1,\mathrm{},r+2`$), where
$$(\alpha _i,l_\pm )=(k_\pm ,l_\pm )=(k_{},l_\pm )=(l_\pm ,l_\pm )=0,(l_{},l_\pm )=1.$$
(5)
In this way starting from the $`r`$dim. Euclidean lattice $`\mathrm{\Lambda }_๐ข`$, we have build up a $`(r+4)`$dim. Lorentzian lattice $`๐ฒ\mathrm{\Lambda }_๐ข`$ II$`{}_{k_\pm }{}^{1,1}`$ II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$, with signature $`+\mathrm{}+`$ ($`r+2`$ plus signs). The simple root system of $`๐ข^{+++}`$ clearly spans a 3-dim. sub-lattice of $`๐ฒ`$ (because $`\alpha _{+3}`$ only takes the direction $`l_++l_{}`$ in the second lattice II<sup>1,1</sup>, so the orthogonal direction is lacking), so it is natural to wonder if it is possible to extend further the $`๐ข^{+++}`$ algebra and fill in $`๐ฒ`$.
Motivated by the previous steps, one would just add another node at the Dynkin diagram of a $`๐ข^{+++}`$ algebra, with a simple link to the root $`\alpha _{+3}`$; in this way, as it happens in the case of the $`๐ข^{++}`$ algebra, one would expect that this construction fills the root lattice $`\mathrm{\Lambda }_๐ข`$II$`{}_{k_\pm }{}^{1,1}`$II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$. In the following, we show that actually this is not the case <sup>2</sup><sup>2</sup>2Our discussion does not take account for the introduction of another lattice II<sup>1,1</sup>, which solves the question only partially, because it moves the problem to fill this new lattice..
As we already mentioned, the simple root $`\alpha _{+3}`$ contains only the combination $`l_++l_{}`$, so if we want to span completely also the second lattice II<sup>1,1</sup>, we need a new simple root $`\alpha _{+4}`$ which allows us to obtain $`l_+`$ and $`l_{}`$ separately. In this way, we are guaranteed that the root lattice of the new algebra contains also all vectors which are integer multiples of $`l_+`$ or $`l_{}`$, so this lattice coincides with $`๐ฒ`$. In the following, letโs choose for $`\alpha _{+4}`$ the more general form:
$$\alpha _{+4}al_++bl_{}+ck_++dk_{}$$
(6)
which allows us to make many considerations about the possible extensions of the $`๐ข^{+++}`$ algebras <sup>3</sup><sup>3</sup>3As it happens also for $`\alpha _{+2}`$ and $`\alpha _{+3}`$, here we do not consider the situation in which $`\alpha _{+4}`$ is also linked with the simple roots $`\alpha _i`$ of $`๐ข`$. However this simple choice does not really change the conclusions of our analysis .. First of all, letโs obtain the expression for $`l_+`$ and $`l_{}`$ using the definition of $`\alpha _{+3}`$ and $`\alpha _{+4}`$:
$$(ab)l_+=(dbc)\theta +(dbc)\alpha _{+1}+d\alpha _{+2}+b\alpha _{+3}+\alpha _{+4},$$
(7)
$$(ba)l_{}=(dac)\theta +(dac)\alpha _{+1}+d\alpha _{+2}+a\alpha _{+3}+\alpha _{+4}.$$
(8)
We see then that it is possible to obtain $`l_+`$ and $`l_{}`$ as linear combinations (over $``$) of the simple roots if and only if $`|ab|=1`$. Only in this case we are guaranteed that we actually catch all the vectors in the lattice II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$. With this in mind, letโs find the other relations $`a,b,c,d`$ have to satisfy, that is let us compute the norm of $`\alpha _{+4}`$ and the scalar products with the other simple roots (here we write only the relevant elements $`a_{ij}`$ of the generalized Cartan matrix):
$$(\alpha _{+4})^2=2(ab+cd),$$
(9)
$$a_{+4,+1}2\frac{(\alpha _{+4},\alpha _{+1})}{(\alpha _{+4},\alpha _{+4})}=\frac{d}{ab+cd},a_{+4,+2}2\frac{(\alpha _{+4},\alpha _{+2})}{(\alpha _{+4},\alpha _{+4})}=\frac{cd}{ab+cd},$$
(10)
$$a_{+4,+3}2\frac{(\alpha _{+4},\alpha _{+3})}{(\alpha _{+4},\alpha _{+4})}=\frac{dab}{ab+cd}.$$
(11)
Let us limit ourselves to consider the case $`(\alpha _{+4})^2=2`$ (to look for a โnatural extensionโ) and let all the four coefficients be different from zero; we must have:
$$ab+cd=1,$$
(12)
$$d<0,cd>0,da+b,$$
(13)
because only with these constraints $`\alpha _{+4}`$ is an acceptable simple root (ร la Kac-Moody). In particular, the product $`cd`$ must be negative, and as eq. (12) holds, the product $`ab`$ must be positive (actually, it must be at least 2), that is $`a`$ and $`b`$ must have the same sign. A simple solution <sup>4</sup><sup>4</sup>4We thank C. Helfgott for pointing us this solution. to eqs. (12)-(13) is the following:
$$a=2,b=1,c=1,d=1,$$
(14)
where $`\alpha _{+4}`$ has norm 2 and the scalar products are $`(\alpha _{+4},\alpha _{+1})=1`$, $`(\alpha _{+4},\alpha _{+2})=0`$ and $`(\alpha _{+4},\alpha _{+3})=4`$. Furthermore, this solution verifies also the condition $`ab=1`$, so its root lattice is precisely $`๐ฒ`$, but it presents one loop and a multiple link with the (last) root $`\alpha _{+3}`$, so its structure is very different from the other previous extensions <sup>5</sup><sup>5</sup>5We thank A. Kleinschmidt for suggesting us another solution $`\alpha _{+4}=\alpha +3k_+2k_{}2l_+3l_{}`$ in which $`\alpha _{+4}`$ is also linked to the only simple root $`\alpha `$ of su(2); this solution too presents loops and multiple links.. This situation is common to all solutions of eqs. (12)-(13): it is not possible (we stress: without adding any II<sup>1,1</sup> or different lattice) to have a simple link with $`\alpha _{+4}`$ and no loops (for example the equation $`dab=1`$ is never satisfied if $`a`$ and $`b`$ are both positive, because d is at least -1, and goes in contradiction with the other relations if $`a`$ and $`b`$ are both negative). In this sense, the procedure of standard extension stops at the third step ($`๐ข^{+++}`$).
It is clear that there exist infinite solutions to eqs. (12)-(13) (for example, it is enough to take $`a,b`$ both positive, $`d=1`$ and an opportune value of $`c`$ to satisfy eq. (12) and the second of eqs. (13)), but only those with $`|ab|=1`$ have the property that their root lattice coincides with $`๐ฒ`$. We can also try (in $`\alpha _{+4}`$) to put one coefficient equal to zero and to explore more specific cases. An investigation, case by case, shows that:
* $`d=0`$ : the scalar products imply $`ab=1`$ and $`a+b>0`$, that is $`a=b=1`$, but the case $`a=b`$ is not acceptable because, in this case, $`\alpha _{+4}`$ is a linear combination of the other simple roots;
* $`c=0`$ : the scalar products with $`\alpha _{+1}`$ and $`\alpha _{+2}`$ are respectively proportional to $`d`$ and $`d`$, so one of the two is positive and has the wrong sign, or $`d=0`$ and then $`\alpha _{+4}`$ is a linear combination of the other simple roots;
* $`a=0`$ : this case implies $`c=1`$, but then $`d=1`$ too and the scalar product with $`\alpha _{+2}`$ is positive;
* $`b=0`$ : it has the same problem as the case $`a=0`$.
If we put equal to zero two coefficients, the only possibility is to have $`c=d=0`$ (we can not take equal to zero a coefficient of $`l_+`$ or $`l_{}`$ and a coefficient of $`k_+`$ or $`k_{}`$, because $`\alpha _{+4}`$ would have vanishing norm), but in this case, as already discussed, we have problems.
So far, we have seen that it is never possible to add a (Kac-Moody) simple root in order to have a simple link with $`\alpha _{+3}`$; besides this, the possible solutions which span the whole lattice $`๐ฒ`$ are restricted to the condition $`|ab|=1`$. Yet, if we abandon the condition that $`\alpha _{+4}`$ is a Kac-Moody simple root and let it be a Borcherds (imaginary) simple root, the situation changes. In fact, we can easily find an imaginary root of the kind:
$$\alpha _{+4}=al_++bl_{}$$
(15)
which has scalar product equal to $`1`$ with $`\alpha _{+3}`$: it is only necessary that $`b=1a`$, that is $`\alpha _{+4}=al_++(1a)l_{}`$, with norm $`2a(1a)`$. Then if $`a=0`$ or $`a=1`$, $`\alpha _{+4}`$ has norm zero, in all the other cases its norm is negative, that is $`\alpha _{+4}`$ is a good Borcherds simple root. If we want to fill in $`๐ฒ`$ with the introduction of a Borcherds simple root, we have always to fulfill the condition $`|ab|=1`$ (this condition is independent from the norm of $`\alpha _{+4}`$), but now it is possible to have both $`(\alpha _{+4},\alpha _{+3})=1`$ and, at the same time, to span the whole lattice $`๐ฒ`$. In fact the values $`a=1,b=0`$ (and viceversa $`a=0,b=1`$), which correspond to $`\alpha _{+4}=l_+`$ ($`\alpha _{+4}=l_{}`$), are the only ones which allow to have:
$$\mathrm{\Lambda }_{roots}(๐ข{}_{}{}^{+++})\mathrm{\Lambda }_{roots}(๐ข{}_{}{}^{++})\text{II}_{l_\pm }^{1,1}=๐ฒ,$$
(16)
$$(\alpha _{+4},\alpha _{+3})=1,$$
(17)
as opposite to the Kac-Moody case, where, as we already said, itโs never possible to satisfy eq. (17) and only some solutions allow to fill in $`๐ฒ`$ (we have called $`๐ข^{+++}`$ the Borcherds extension of $`๐ข^{+++}`$ corresponding to $`\alpha _{+4}=l_+`$ or $`\alpha _{+4}=l_{}`$). Denoting by a crossed dot the Borcherds simple root $`\alpha _{+4}`$, we can draw the Dynkin diagram of this Borcherds algebra in the following way:
Actually, this construction makes sense if $`๐ข`$ is simply-laced, because otherwise the Cartan matrix is not well-defined. In fact, if $`\alpha _{+4}`$ is an imaginary (isotropic) root, we cannot define the extended Cartan matrix as $`2\frac{(\alpha _i,\alpha _j)}{(\alpha _i,\alpha _i)}`$ for all $`\alpha _i`$ because $`\alpha _{+4}^2=0`$. We have the same problem if we add an imaginary simple root with negative squared norm (as previously recalled, in II<sup>1,1</sup> there are infinite vectors whose squared norm is negative), because we shall have positive elements out of the principal diagonal. The solution is then to consider a $`๐ข`$ simply-laced and define the extended Cartan matrix by the scalar products between all the simple roots: $`b_{i,j}:=(\alpha _i,\alpha _j)`$. So, while the extension of $`๐ข`$ is possible up to $`๐ข^{+++}`$ for any finite $`๐ข`$, in the case of a Borcherds extension it is necessary to choose for $`๐ข`$ a simply-laced algebra (in fact, a Borcherds algebra is defined only on a symmetric Cartan matrix). To summarize our result, we have proven the following
###### Proposition 1
The extension of $`๐ข^{+++}`$ algebra, whose simple root system spans completely $`๐ฒ`$ and whose Dynkin-Kac diagram has no loops and only simple links between the dots, is the Borcherds algebra $`๐ข^{+++}`$ (with $`๐ข`$ simply-laced).
The particular solution $`a=1,b=0`$ (or viceversa) looks like the same construction of $`๐ข^+`$ and $`๐ข^{++}`$, as the role of $`k_\pm `$ is now played by $`l_\pm `$, with the difference that $`\alpha _{+1}`$ contains also a root of a simple Lie algebra, while $`\alpha _{+4}`$ doesnโt. This observation suggests that itโs possible to fuse together two (finite) simple Lie algebras, letโs say $`๐ข`$ and $`๐ข^{}`$, adding the highest root $`\theta ^{}`$ of $`๐ข^{}`$ to $`\alpha _{+4}`$. In this way, $`\alpha _{+4}^{}\alpha _{+4}\theta ^{}`$ is again a Kac-Moody simple root, indeed the affine root of $`๐ข^+`$. Anyway, we donโt insist on this point, because there are many possible ways to fuse together two (or more) finite dimensional simple Lie algebras (with or without the introduction of intermediate Kac-Moody or Borcherds algebras).
Remark: In this section, we have seen that triple extended Lie algebras $`๐ข^{+++}`$ have their root lattice properly included in $`๐ฒ=\mathrm{\Lambda }_๐ข`$ II$`{}_{}{}^{1,1}`$ II<sup>1,1</sup> (in particular, their root lattice is Lorentzian with just one negative eigenvalue). Among the 4-extensions we considered in this Section, we have shown that there are many algebras (of Kac-Moody or generalized Kac-Moody type) whose root lattice coincides with $`๐ฒ`$, which is Lorentzian in a more general sense (it has two negative eigenvalues). These seem to be the first algebras obtained in literature with this kind of lattice, together with some similar algebras studied by Harvey and Moore in , (strictly speaking, their algebras are not Borcherds algebras, because they could not satisfy some grading conditions in the characterization of generalized Kac-Moody algebras, while our Borcherds algebras are *true* Borcherds algebras because we have constructed them from an acceptable generalized Cartan matrix). Our result seems to be in contrast with the statement of , according to which algebras whose root lattices are of the kind $`\mathrm{\Gamma }^{p,q}`$ (that is, with a signature with more that one negative sign) cannot be described in terms of generators and relations similar to Kac-Moody or Borcherds algebras and belong to a new class of Lie algebras. That result has been found looking for Lie algebras of the physical states of a vertex algebra constructed on a general even self-dual lattice $`\mathrm{\Gamma }^{p,q}`$; maybe the condition of physical states is too strong and is not compatible with the algebras constructed by us.
## 3 Folding of triple extended Lie algebras
The *folding* technique is a simple and powerful method to find class of singular subalgebras of finite Lie algebras as well as of affine or indefinite Kac-Moody algebras. The starting point is to use the symmetry $`\tau `$ of the Dynkin diagram, corresponding to an exterior automorphism of the algebra $`๐ข`$. In the finite case all the Dynkin diagrams of simply-laced algebras show an automorphism of order $`k=2`$, except the case of D<sub>4</sub> where the order is 3. Let $`\alpha _i`$ be a simple root of $`๐ข`$. Using the automorphism $`\tau `$ of order $`k`$, we obtain
$$\beta _i:=\alpha _i+\tau (\alpha _i)+\mathrm{}+\tau ^{k1}(\alpha _i)$$
(18)
which form the simple root system of a singular subalgebra $``$ of $`๐ข`$. The generators of $`๐ข`$, corresponding to the simple roots $`\alpha _i`$, left unchanged by $`\tau `$, become generators of $``$, while the other ones transform according to a relation analogous to eq.(18). In the following we apply the folding method to the triple extended Lie algebras, obtaining all the non simply laced triple extended algebras, as in the finite case. The automorphism of the $`3`$-extended Dynkin diagram acts on the standard way upon the roots of the finite classical subalgebras and trivially on the extended roots $`\alpha _{+1},\alpha _{+2},\alpha _{+3}`$. This property has to hold if we want to preserve the structure of the triple extension of the non-simply laced algebras <sup>6</sup><sup>6</sup>6Indeed, also the non simply-laced $`G^{++}`$ algebras can be obtained with the same folding technique from the simply-laced $`๐ข^{++}`$, while a different kind of folding applied to the $`๐ข^+`$ algebras allows to get all the twisted affine algebras.. Let us enumerate all the cases.
### 3.1 A$`{}_{2N1}{}^{+++}`$ C$`{}_{}{}^{+++}{}_{N}{}^{}`$
The Cartan matrix of A$`{}_{}{}^{+++}{}_{2N1}{}^{}`$ can be written, in block form, as
$$A=(a_{ij})=(\alpha _i,\alpha _j)=\left(\begin{array}{ccccccc}& & & 1& 0& 0& \\ & A_{A_{2N1}}& & \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & 1& 0& 0& \\ 1& \mathrm{}& 1& 2& 1& 0& \\ 0& \mathrm{}& 0& 1& 2& 1& \\ 0& \mathrm{}& 0& 0& 1& 2& \end{array}\right)$$
(19)
The not trivial action of $`\tau `$ on the simple roots gives
$`\beta _1`$ $`:=`$ $`\alpha _1+\tau (\alpha _1)=\alpha _1+\alpha _{2N1}`$ (20)
$`\beta _2`$ $`:=`$ $`\alpha _2+\tau (\alpha _2)=\alpha _2+\alpha _{2N2}\mathrm{}`$ (21)
$`\beta _{N1}`$ $`:=`$ $`\alpha _{N1}+\tau (\alpha _{N1})=\alpha _{N1}+\alpha _{N+1}`$ (22)
$`\beta _N`$ $`:=`$ $`\alpha _N+\tau (\alpha _N)=2\alpha _N`$ (23)
$`\beta _{+1}`$ $`:=`$ $`\alpha _{+1}+\tau (\alpha _{+1})=2\alpha _{+1}`$ (24)
$`\beta _{+2}`$ $`:=`$ $`\alpha _{+2}+\tau (\alpha _{+2})=2\alpha _{+2}`$ (25)
$`\beta _{+3}`$ $`:=`$ $`\alpha _{+3}+\tau (\alpha _{+3})=2\alpha _{+3}`$ (26)
the length of the simple roots is
$`\beta _1^2`$ $`=`$ $`\mathrm{}=\beta _{N1}^2=4`$ (27)
$`\beta _{+3}^2`$ $`=`$ $`\beta _{+2}^2=\beta _{+1}^2=\beta _N^2=8.`$ (28)
The corresponding Cartan matrix is
$$B=(b_{ij})_{i,j}=2\frac{(\beta _i,\beta _j)}{(\beta _i,\beta _i)}=\left(\begin{array}{ccccccc}& & & 2& 0& 0& \\ & A_{C_N}& & \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & 0& 0& 0& \\ 1& \mathrm{}& 0& 2& 1& 0& \\ 0& \mathrm{}& 0& 1& 2& 1& \\ 0& \mathrm{}& 0& 0& 1& 2& \end{array}\right)$$
(29)
So we get the $`3`$-extended Lie algebra C$`{}_{}{}^{+++}{}_{N}{}^{}`$ with Dynkin diagram
### 3.2 D$`{}_{N}{}^{+++}`$ B$`{}_{}{}^{+++}{}_{N1}{}^{}`$
The Cartan matrix can be written as
$$A=(a_{ij})_{i,j}=(\alpha _i,\alpha _j)=\left(\begin{array}{ccccccc}& & & & 0& 0& 0\\ & & & & 1& \mathrm{}& \mathrm{}\\ & & A_{D_N}& & \mathrm{}& \mathrm{}& \mathrm{}\\ & & & & 0& 0& 0\\ 0& 1& \mathrm{}& 0& 2& 1& 0\\ 0& \mathrm{}& \mathrm{}& 0& 1& 2& 1\\ 0& \mathrm{}& \mathrm{}& 0& 0& 1& 2\end{array}\right)$$
(30)
The non trivial action of $`\tau `$ is
$$\tau (\alpha _{N1})=\alpha _N,\tau (\alpha _N)=\alpha _{N1}$$
(31)
The new simple roots are
$`\beta _{N1}`$ $`:=`$ $`\alpha _{N1}+\tau (\alpha _{N1})=\alpha _{N1}+\alpha _N`$ (32)
$`\beta _{N2}`$ $`:=`$ $`\alpha _{N2}+\tau (\alpha _{N2})=2\alpha _{N2}\mathrm{}`$ (33)
$`\beta _1`$ $`:=`$ $`\alpha _1+\tau (\alpha _1)=2\alpha _1`$ (34)
$`\beta _{+1}`$ $`:=`$ $`\alpha _{+1}+\tau (\alpha _{+1})=2\alpha _{+1}`$ (35)
$`\beta _{+2}`$ $`:=`$ $`\alpha _{+2}+\tau (\alpha _{+2})=2\alpha _{+2}`$ (36)
$`\beta _{+3}`$ $`:=`$ $`\alpha _{+3}+\tau (\alpha _{+3})=2\alpha _{+3}`$ (37)
where
$$\beta _{N1}^2=4,\beta _{+3}^2=\beta _{+2}^2=\mathrm{}=\beta _{N2}^2=8.$$
(38)
The Cartan matrix and the corresponding Dynkin diagram of B$`{}_{}{}^{+++}{}_{N1}{}^{}`$ are
$$B=(b_{ij})_{i,j}=2\frac{(\beta _i,\beta _j)}{(\beta _i,\beta _i)}=\left(\begin{array}{ccccccc}& & & & 0& 0& 0\\ & & & & 1& \mathrm{}& \mathrm{}\\ & & A_{B_{N1}}& & \mathrm{}& \mathrm{}& \mathrm{}\\ & & & & 0& 0& 0\\ 0& 1& \mathrm{}& 0& 2& 1& 0\\ 0& \mathrm{}& \mathrm{}& 0& 1& 2& 1\\ 0& \mathrm{}& \mathrm{}& 0& 0& 1& 2\end{array}\right)$$
(39)
### 3.3 E$`{}_{6}{}^{+++}`$ F$`{}_{}{}^{+++}{}_{4}{}^{}`$
The Cartan matrix of E$`{}_{}{}^{+++}{}_{6}{}^{}`$ is
$$A=(a_{ij})_{i,j}=(\alpha _i,\alpha _j)=\left(\begin{array}{cccccc}& & & 0& 0& 0\\ & A_{E_6}& & \mathrm{}& \mathrm{}& \mathrm{}\\ & & & 1& 0& 0\\ 0& \mathrm{}& 1& 2& 1& 0\\ 0& \mathrm{}& 0& 1& 2& 1\\ 0& \mathrm{}& 0& 0& 1& 2\end{array}\right)$$
(40)
The simple roots are given by
$`\beta _1`$ $`:=`$ $`\alpha _1+\tau (\alpha _1)=\alpha _1+\alpha _5`$ (41)
$`\beta _2`$ $`:=`$ $`\alpha _2+\tau (\alpha _2)=\alpha _2+\alpha _4`$ (42)
$`\beta _3`$ $`:=`$ $`\alpha _3+\tau (\alpha _3)=2\alpha _3`$ (43)
$`\beta _4`$ $`:=`$ $`\alpha _6+\tau (\alpha _6)=2\alpha _6`$ (44)
$`\beta _{+1}`$ $`:=`$ $`2\alpha _{+1},\beta _{+2}:=2\alpha _{+2}`$ (45)
$`\beta _{+3}`$ $`:=`$ $`2\alpha _{+3}`$ (46)
and
$$\beta _1^2=\beta _2^2=4,\beta _{+3}^2=\beta _{+2}^2=\mathrm{}=\beta _3^2=8.$$
(47)
One gets the $`3`$-extended algebra F$`{}_{}{}^{+++}{}_{4}{}^{}`$, with Cartan matrix and Dynkin diagram
$$B=(b_{ij})_{i,j}=2\frac{(\beta _i,\beta _j)}{(\beta _i,\beta _i)}=\left(\begin{array}{cccccc}& & & 0& 0& 0\\ & A_{F_4}& & \mathrm{}& \mathrm{}& \mathrm{}\\ & & & 1& 0& 0\\ 0& \mathrm{}& 1& 2& 1& 0\\ 0& \mathrm{}& 0& 1& 2& 1\\ 0& \mathrm{}& 0& 0& 1& 2\end{array}\right)$$
(48)
### 3.4 $`D_4^{+++}G_2^{+++}`$
The Cartan matrix of D$`{}_{}{}^{+++}{}_{4}{}^{}`$ is
$$A=(a_{ij})_{i,j}=(\alpha _i,\alpha _j)=\left(\begin{array}{ccccccc}2& 1& 0& 0& 0& 0& 0\\ 1& 2& 1& 1& 1& 0& 0\\ 0& 1& 2& 0& 0& 0& 0\\ 0& 1& 0& 2& 0& 0& 0\\ 0& 1& 0& 0& 2& 1& 0\\ 0& 0& 0& 0& 1& 2& 1\\ 0& 0& 0& 0& 0& 1& 2\end{array}\right)$$
(49)
The action of $`\tau `$ on the simple roots gives
$`\beta _1`$ $`:=`$ $`\alpha _1+\tau (\alpha _1)+\tau ^2(\alpha _1)=\alpha _1+\alpha _3+\alpha _4`$ (50)
$`\beta _2`$ $`:=`$ $`3\alpha _2`$ (51)
$`\beta _{+1}`$ $`:=`$ $`3\alpha _{+1},\beta _{+2}:=3\alpha _{+2}`$ (52)
$`\beta _{+3}`$ $`:=`$ $`3\alpha _{+3}`$ (53)
and
$$\beta _1^2=6,\beta _{+3}^2=\mathrm{}=\beta _2^2=18.$$
(54)
One gets the extended algebra G$`{}_{}{}^{+++}{}_{2}{}^{}`$ with Cartan matrix and Dynkin diagram:
$$B=(b_{ij})_{i,j}=2\frac{(\beta _i,\beta _j)}{(\beta _i,\beta _i)}=\left(\begin{array}{ccccc}2& 3& 0& 0& 0\\ 1& 2& 1& 0& 0\\ 0& 1& 2& 1& 0\\ 0& 0& 1& 2& 1\\ 0& 0& 0& 1& 2\end{array}\right)$$
(55)
Let us remark that, in general, by applying the folding procedure to an algebra $`๐ข`$, defined by the Cartan matrix $`A`$ and Dynkin diagram $`S(A)`$, we obtain the Cartan matrix $`B`$ and the corresponding Dynkin diagram $`S(B)`$ of another algebra $``$. However we have to check that the new generators, defined in function of the generators of of $`๐ข`$, satisfy all the defining relations of the algebra $``$. Let us see how the new generators are obtained. Let $`\alpha _i`$ be a simple root of of $`๐ข`$ and $`h_i,e_i,f_i`$ the associated generators. Let us denote by $`\beta _i`$, respectively $`h_i^{},e_i^{},f_i^{}`$, the roots and the associated generators transformed under the action of $`\tau `$, which we identify, respectively, as the simple roots and the associated generators of $``$. If the action of $`\tau `$ is trivial, that is $`\beta _i=k\alpha _i`$ (where $`k`$ is the order of the automorphism $`\tau `$, $`\tau ^k=1`$) then the generators are not transformed $`h_i^{}=h_i`$, $`e_i^{}=e_i`$ and $`f_i^{}=f_i`$. If the action of $`\tau `$ is not trivial, that is $`\beta _i=\alpha _i+\tau (\alpha _i)+\mathrm{}+\tau ^{k1}(\alpha _i)`$, we obtain $`h_i^{}=h_i+h_{\tau (\alpha _i)}+\mathrm{}+h_{\tau ^{k1}(\alpha _i)}`$, $`e_i^{}=e_i+e_{\tau (\alpha _i)}+\mathrm{}+e_{\tau ^{k1}(\alpha _i)}`$ and $`f_i^{}=f_i+f_{\tau (\alpha _i)}+\mathrm{}+f_{\tau ^{k1}(\alpha _i)}`$. We have to verify that the generators $`h_i^{},e_i^{},f_i^{}`$ satisfy the defining relations
$$[e_i^{},f_j^{}]=\delta _{ij}h_j^{},$$
$$[h_i^{},e_j^{}]=b_{ij}e_j^{},[h_i^{},f_j^{}]=b_{ij}f_j^{},$$
$$[h_i^{},h_j^{}]=0,$$
$$(\text{ad}e_i^{})^{1b_{ij}}e_j^{}=0,(\text{ad}f_i^{})^{1b_{ij}}f_j^{}=0,(ij).$$
We do no report here the explicit calculations, but everything works nicely. Finally it should be remarked that the folding procedure for indefinite Kac-Moody algebra, when applicable, always gives rise to indefinite Kac-Moody algebra, as it happens for the finite, affine, hyperbolic Kac-Moody algebras. On the contrary other reduction procedures, as the orbifolding, do not preserve the kind of algebras. Indeed as remarked in , the orbifolding of $`E_{10}`$, to which the folding procedure cannot be applied, gives rise to non Kac-Moody algebras.
## 4 Non-standard extensions of Lie algebras
In this section we present a non-standard construction of extended Lie algebras; as stated in Sec. 1, the idea of the non-standard extension is to add to the simple root system {$`\alpha _i`$} of a simple Lie algebra $`๐ข`$ new roots, which are formed by those fundamental weights of the algebra that are linear combinations with integer coefficients of $`\alpha _i`$, plus a suitable combinations of vectors belonging to the Lorentzian lattice II$`{}_{k_\pm }{}^{}{}_{}{}^{1,1}`$ and/or II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$. The new roots have to satisfy the requirements that their squared norms are equal to 2 and that are suitably linked with the previous ones. Let us remark that the roots of the non-standard extension do not generally span the whole lattice $`๐ฒ=\mathrm{\Lambda }_๐ข\text{II}_{k_\pm }^{1,1}\text{II}_{l_\pm }^{1,1}`$ and that, moreover, the structure of the added simple root is, by no way, unique. Of course one can add more than two $`2`$-dim. Lorentzian lattices, but these extensions will not be considered in the present paper, where we add at most three new roots. Also we shall not discuss the case where the squared norm of the added roots is not equal to 2. So, given a simple Lie algebra $`๐ข`$, we add to the root lattice $`\mathrm{\Lambda }_๐ข`$ a new simple root $`\alpha _{r+1}\alpha _{+1}`$, which is formed by the opposite of a fundamental weight $`\mathrm{\Lambda }_i`$ and by a suitable linear combination, with integer coefficients, of the vectors $`k_\pm `$, in order to have $`\alpha _{+1}^2=2`$, as $`\mathrm{\Lambda }_i^2`$ is not necessarily 2. Letโs remember that the fundamental weights have the property $`2\frac{(\mathrm{\Lambda }_i,\alpha _j)}{(\alpha _j,\alpha _j)}=\delta _{i,j}`$; they span the weight lattice $`P=_{i=1}^r\mathrm{\Lambda }_i`$, which is dual to the coroot lattice $`\mathrm{\Lambda }_๐ข^{}=_{i=1}^r\alpha _i^{}`$ where $`\alpha _i^{}=2\frac{\alpha _i}{(\alpha _i,\alpha _i)}`$ are the coroots. So while it is always true that $`P=(\mathrm{\Lambda }_๐ข^{})^{}`$, in general we have $`\mathrm{\Lambda }_๐ขP`$. This means that for each $`๐ข`$, only some $`\mathrm{\Lambda }_i`$ belong to the root lattice; so, in defining $`\alpha _{+1}`$ we choose the $`\mathrm{\Lambda }_i\mathrm{\Lambda }_๐ข`$. Since $`\mathrm{\Lambda }_i`$ is only linked with the simple root $`\alpha _i`$, we have $`(\alpha _{+1},\alpha _j)=\delta _{+1,i}`$. So the first extension is made by adding the root $`\alpha _{r+1}\alpha _{+1}:=\mathrm{\Lambda }_i+k_+ak_{}`$, where $`a๐_+`$ is fixed by the condition $`\alpha _{+1}^2=\mathrm{\Lambda }_i^22a=2.`$ At this point, we add the simple root $`\alpha _{r+2}\alpha _{+2}:=\theta +bk_{}l_{}`$, where $`\theta `$ is the highest root of $`๐ข`$ and $`b๐`$ is a coefficient chosen in order to have $`(\alpha _{+2},\alpha _{+1})=(\mathrm{\Lambda }_i,\theta )b=0`$. In this way, we have $`\alpha _{+2}^2=2`$ and $`\alpha _{+2}`$ behaves like an affine root (that is, it is linked with the simple roots of $`๐ข`$ in a way completely analogous as the affine root of the algebras $`\widehat{๐ข}`$). At the end, we add the third simple root $`\alpha _{r+3}\alpha _{+3}:=l_++l_{}`$ with the property that $`(\alpha _{+3},\alpha _{+2})=1`$ and $`(\alpha _{+3},\alpha _i)=0`$ for $`i=1,\mathrm{},r,r+1`$. As $`(\mathrm{\Lambda }_i,\mathrm{\Lambda }_j)๐_>`$ for the Lie algebra $`๐ข`$ below considered, this procedure is completely general and the extended algebra contains as subalgebra the affine extension of $`๐ข`$ (so sometimes we shall call this a non standard affine extension). Clearly the lightlike vector $`l_{}`$ can be hanged up to any other simple root, producing an other indefinite Kac-Moody. We shall comment on this point in Sec.5. This construction leads to indefinite Kac-Moody algebras, 1-, 2- and 3-extended, whose (symmetric) Cartan matrix $`2\frac{(\alpha _i,\alpha _j)}{(\alpha _i,\alpha _i)}`$ (for $`i,j=1,\mathrm{},r+3`$) has Lorentzian signature $`(+\mathrm{}+)`$ with $`r+2`$ plus signs and 1 minus sign. The root lattice of the 3-extended algebra is properly contained in $`\mathrm{\Lambda }_๐ข`$ II$`{}_{k_\pm }{}^{1,1}`$ II$`{}_{l_\pm }{}^{}{}_{}{}^{1,1}`$. For these algebras, one can make similar discussions as those in Sec. 2 on eventual further extensions. Now we want to discuss another possible extension, which cannot be performed for any fundamental weight $`\mathrm{\Lambda }_i`$ belonging to the root lattice of $`๐ข`$. The first extension is performed as before, but as second extension we add the root $`\alpha _{+2}:=\mathrm{\Lambda }_j+k_+bk_{}l_{}`$ ($`ij`$), where $`b๐_+`$ is such that $`\alpha _{+2}^2=2`$ and $`(\alpha _{+1},\alpha _{+2})=(\mathrm{\Lambda }_i,\mathrm{\Lambda }_j)ab=0`$. Below we shall show that, for $`๐ขE_6`$, for any $`i`$ ($`\mathrm{\Lambda }_i\theta `$), at least one $`j`$ exists which satisfies the above condition. In the following, we discuss only some examples of the general construction, we called affine extension; in particular we concentrate on the simply laced algebras, but it is possible to consider also the other cases paying attention at the choice of the fundamental weight.
Looking at the fundamental weights of simply laced-Lie algebras, see , one realizes that the fundamental weights which can be written as
$$\mathrm{\Lambda }_i=\underset{n}{}c_n\alpha _nc_n๐$$
(56)
are
1. for $`D_N=so(2N)`$ ($`N4`$), the weights $`\mathrm{\Lambda }_i`$ with $`i`$ even number ($`N2i2`$)
2. for $`E_6`$, only the weights $`\mathrm{\Lambda }_i`$ ($`i=3,6`$)
3. for $`E_7`$, only the weights $`\mathrm{\Lambda }_i`$ ($`i=1,2,3,5`$)
4. for $`E_8`$, all the weights $`\mathrm{\Lambda }_i`$, which is just a consequence of the $`E_8`$-lattice being a self-dual one.
In the following we discuss some of the possible non-standard extensions, with the aim to illustrate the procedure in a few examples which may be relevant for their subalgebras content. Let us emphasize that the discussed extensions as well their subalgebras content are not at all exhaustive, being the choice of the extended simple roots not unique, in general.
### 4.1 $`D_N=so(2N)`$
In order to illustrate the general procedure, we discuss in some detail the case of $`D_6=so(12)`$, which is the first algebra of the even orthogonal series which admits a non-standard extension. We add the simple root
$$\alpha _{+1}:=\mathrm{\Lambda }_4+k_+k_{},\alpha _{+1}^2=2,$$
(57)
where $`\mathrm{\Lambda }_4`$ is the fundamental weight <sup>7</sup><sup>7</sup>7the $`\epsilon _i`$ are unit ortho-normal vectors in $`^6`$.
$$\mathrm{\Lambda }_4=\epsilon _1+\epsilon _2+\epsilon _3+\epsilon _4,\mathrm{\Lambda }_3^2=4,$$
(58)
Clearly we have
$$(\alpha _{+1},\alpha _i)=\delta _{4,i}.$$
(59)
We add now the root
$$\alpha _{+2}:=h.r.2k_{}l_{}=(\epsilon _1+\epsilon _2)2k_{}l_{},$$
(60)
$$(\alpha _{+2},\alpha _j)=\delta _{2,j},$$
(61)
and
$$\alpha _{+3}:=l_++l_{},$$
(62)
$$(\alpha _{+3},\alpha _k)=\delta _{+2,k},$$
(63)
with Dynkin diagram:
Letโs observe that the first extension of $`D_6`$ is the same algebra as $`D_4^{+++}`$, so folding $`\alpha _{+1}`$, $`\alpha _5`$ and $`\alpha _6`$ we re-obtain $`G_2^{+++}`$. Clearly the choice of the extended simple roots is not unique. One can easily see that:
* a non-standard extension of $`so(4N)`$ admits as subalgebra the affine extension of $`so(4N)`$ and $`so(4(N1))`$. Indeed one adds to the roots of $`so(4N)`$ the root of the affine extension
$$\alpha _{+1}:=\mathrm{\Lambda }_2+k_+$$
(64)
and the new non-standard root
$$\alpha _{+2}:=\mathrm{\Lambda }_4+l_+l_{}$$
(65)
Taking away the roots $`\alpha _{+1},\alpha _j`$ ($`j=1,2`$) one gets the algebra $`so(4(N1))^{(1)}`$. Let us remark that if we add the root
$$\alpha _{+2}:=\mathrm{\Lambda }_6+l_+2l_{}$$
(66)
and then we take away the roots $`\alpha _{+1},\alpha _j`$ ($`j=1,\mathrm{},4`$) one gets the algebra $`so(4(N2))^{(1)}`$.
* the non standard extension of $`so(24)`$ is the smallest extension of the orthogonal series which contains as subalgebra $`E_{11}`$. Indeed adding to the roots of $`so(24)`$ the non standard root
$$\alpha _{+1}:=\mathrm{\Lambda }_8+k_+3k_{}$$
(67)
and deleting $`\alpha _{11},\alpha _{12}`$ one gets $`E_{11}`$.
Let us call $`\widehat{\mathrm{\Lambda }}_{2n}=\mathrm{\Lambda }_{2n}+k_+(n1)k_{}`$, where $`n_+`$ and $`\mathrm{\Lambda }_{2n}=_{i=1}^{2n}\epsilon _i`$ is a fundamental weight. Clearly we have
$$\widehat{\mathrm{\Lambda }}_{2n}^2=2(\widehat{\mathrm{\Lambda }}_{2n},\widehat{\mathrm{\Lambda }}_{2n+2})=0$$
(68)
### 4.2 $`E_6`$
Let us add to the simple root system of $`E_6`$ the root
$$\alpha _{+1}:=\mathrm{\Lambda }_3+k_+2k_{},\alpha _{+1}^2=2,$$
(69)
where $`\mathrm{\Lambda }_3`$ is the fundamental weight
$$\mathrm{\Lambda }_3=\epsilon _3+\epsilon _4+\epsilon _5+\epsilon _8\epsilon _7\epsilon _6,\mathrm{\Lambda }_3^2=6.$$
(70)
Clearly we have
$$(\alpha _{+1},\alpha _i)=\delta _{3,i}.$$
(71)
We add now the root
$$\alpha _{+2}:=h.r.3k_{}l_{}=\frac{1}{2}(\epsilon _1+\epsilon _2+\epsilon _3+\epsilon _4+\epsilon _5\epsilon _6\epsilon _7+\epsilon _8)3k_{}l_{},$$
(72)
$$(\alpha _{+2},\alpha _j)=\delta _{6,j},$$
(73)
and
$$\alpha _{+3}:=l_++l_{},(\alpha _{+3},\alpha _k)=\delta _{+2,k}.$$
(74)
Alternatively we can add the roots
$`\alpha _{+2}`$ $`:=`$ $`k_{}+k_+l_{},(\alpha _{+2},\alpha _j)=\delta _{+1,j},`$ (75)
$`\alpha _{+3}`$ $`:=`$ $`l_++l_{},(\alpha _{+3},\alpha _k)=\delta _{+2,k},`$ (76)
and we obtain the same Dynkin diagram:
where in the second construction the roles of $`\alpha _{+1}`$ and $`\alpha _6`$ are exchanged. Letโs observe that many non-standard (simply-laced) Dynkin diagrams can be folded to obtain other (non simply-laced) Dynkin diagrams. For example, in the case of $`E_6`$, we can identify the roots $`\alpha _1`$ and $`\alpha _2`$ with $`\alpha _5`$ and $`\alpha _4`$ respectively (so the new simple roots are $`\beta _1=\alpha _1+\alpha _5`$, $`\beta _2=\alpha _2+\alpha _4`$ and $`\beta _i=2\alpha _i`$ for $`i=3,+1,+2,+3`$), obtaining the following folded Dynkin diagram:
Actually, if we consider only the first extension of $`E_6`$, then we can identify also $`\alpha _{+1}`$ with $`\alpha _6`$ and we obtain:
### 4.3 $`E_7`$
In this case we could use the fundamental weights $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$, $`\mathrm{\Lambda }_3`$ and $`\mathrm{\Lambda }_5`$. To illustrate the procedure, letโs consider the weight $`\mathrm{\Lambda }_5`$ ($`(\mathrm{\Lambda }_5,\mathrm{\Lambda }_5)=4`$) and add the simple roots:
$$\alpha _{+1}:=\mathrm{\Lambda }_5+k_+k_{},\alpha _{+1}^2=2,(\alpha _{+1},\alpha _i)=\delta _{5,i},$$
(77)
with
$$\mathrm{\Lambda }_5=\epsilon _5+\epsilon _6\epsilon _7+\epsilon _8,\mathrm{\Lambda }_5^2=4.$$
(78)
Then
$$\alpha _{+2}:=h.r.2k_{}l_{}=\epsilon _7\epsilon _82k_{}l_{}\alpha _{+3}=l_++l_{},$$
(79)
with $`(\alpha _{+2},\alpha _k)=\delta _{k,\mathrm{\hspace{0.17em}1}}`$ and $`(\alpha _{+3},\alpha _j)=\delta _{j,+2}`$. In this way we obtain the Dynkin diagram:
Let us call $`\widehat{\mathrm{\Lambda }}_i:=\mathrm{\Lambda }_i+k_+ak_{}`$, where $`a_+`$ and $`\mathrm{\Lambda }_i`$ , $`i=1,2,3,5`$, is a fundamental weight. We have
$$\widehat{\mathrm{\Lambda }}_i^2=2(\widehat{\mathrm{\Lambda }}_3,\widehat{\mathrm{\Lambda }}_5)=0$$
(80)
### 4.4 $`E_8`$
$`E_8`$ root lattice is self-dual, so it coincides with the weight lattice. The non-standard extension can be made adding to the simple root system a root equal to the opposite of any weight $`\mathrm{\Lambda }_i`$ ($`i=1,\mathrm{},8`$) plus some combination of $`k_+`$ and $`k_{}`$. In this way we have 8 different extensions of $`E_8`$, with the nodes $`+2`$, $`+3`$ always in the same position (that is $`(\alpha _{+2},\alpha _7)=1`$ and $`(\alpha _{+3},\alpha _{+2})=1`$), while the node $`+1`$ moves from the node 1 to the node 8, when $`i`$ runs from 1 to 8 respectively. Actually, this situation is general for the non-standard extensions.
Let us emphasize again that the choice of the extended simple roots is not unique at all. Motivated by this consideration, we observe that classically we have $`E_6E_7E_8`$, while this inclusion is lost when we consider the corresponding affine algebras (and the same thing is true for the double and the triple extensions). So we look for an algebra that may contain all the $`E`$ series and the $`E^{(1)}`$ series. This is possible considering the following non-standard extension of $`E_8`$ (in which the new simple roots are linked to those of $`E_8`$ using different fundamental weights). We add the new simple root:
$$\alpha _{+1}:=h.r.+k_++l_+=(\epsilon _7+\epsilon _8)+k_++l_+,\alpha _{+1}^2=2,$$
(81)
Then we add the two simple roots:
$$\mathrm{\Lambda }_1=2\epsilon _8,\mathrm{\Lambda }_8=\frac{1}{2}\left(\underset{i=1}{\overset{7}{}}\epsilon _i+5\epsilon _8\right),$$
(82)
$$\alpha _{+2}:=\mathrm{\Lambda }_1k_{}+l_+l_{},\alpha _{+1}^2=2,$$
(83)
and
$$\alpha _{+3}:=\mathrm{\Lambda }_8+k_++l_+3l_{},\alpha _{+3}^2=2,$$
(84)
so that the only non-zero scalar products are: $`(\alpha _{+1},\alpha _7)=1`$, $`(\alpha _{+2},\alpha _1)=1`$ and $`(\alpha _{+3},\alpha _8)=1`$ and we have the following Dynkin diagram:
This algebra contains $`E_8`$ (then also $`E_7`$ and $`E_6`$) and all the affinizations $`E_{6,7,8}^{(1)}`$, so it seems to be interesting for its content in sub-algebras. Let us call $`\widehat{\mathrm{\Lambda }}_i=\mathrm{\Lambda }_i+k_+ak_{}`$, where $`a_+`$ and $`\mathrm{\Lambda }_i`$, is any fundamental weight. We have ($`\widehat{\mathrm{\Lambda }}_i^2=2`$)
$`(\widehat{\mathrm{\Lambda }}_1,\widehat{\mathrm{\Lambda }}_2)=(\widehat{\mathrm{\Lambda }}_1,\widehat{\mathrm{\Lambda }}_5)=0`$
$`(\widehat{\mathrm{\Lambda }}_4,\widehat{\mathrm{\Lambda }}_8)=0`$ $`(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_3)=(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_4)=(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_5)=(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_6)=0`$ (85)
$`(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_3)=(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_4)=(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_5)=(\widehat{\mathrm{\Lambda }}_2,\widehat{\mathrm{\Lambda }}_6)=0`$
We have proposed a procedure to build non standard triple extended Lie algebras, which we have illustrated with a number of relevant examples. As already recalled, Kac-Moody or Borcherds extensions can be defined for these Lie algebras too. On the light of the remarks of , it is natural to wonder if these or some of these algebras are not really subalgebras of the standard triple extended Lie algebras. This point will be discussed in the next section, where we discuss also a few examples of subalgebras which point out the intriguing and surprising structure of the subalgebras.
## 5 Subalgebras of extended Lie algebras
First of all, let us discuss another non-standard procedure to extend a Lie algebras of rank $`r`$. If one adds to the simple roots of $`๐ข`$ the opposite of the h.r. $`\alpha _0\alpha _{r+1}`$ and, to recover the linear independence of the simple root system, one glues the light like vector $`k_+`$ to a simple root $`\alpha _i`$ ($`1ir`$), one gets exactly the affine $`๐ข`$. As next step one adds the new root, which belongs to II$`{}_{k}{}^{}{}_{}{}^{1,1}`$, $`\alpha _{r+2}=(k_++k_{})`$. From the Feingold-Nicolai theorem , it is easy to realize that in this way one obtains a generalized Kac-Moody algebra which is really a subalgebra of the standard overextended Lie algebra $`๐ข^{++}`$. Things may be different if one considers simple root systems, obtained by analogous procedure, in the lattice $`\mathrm{\Lambda }_๐ข\text{II}_k^{1,1}\text{II}_l^{1,1}`$. Indeed indefinite Kac-Moody algebras are obtained which, in general, are described by Dynkin-Kac diagrams not equivalent to the ones obtained by the standard and not standard procedure described in the previous section. A general discussion of these algebras is beyond the aim of this paper and we limit ourselves to state a few properties and to present some examples. Let us start with the following
###### Proposition 2
The roots $`\widehat{\mathrm{\Lambda }}_i`$ defined in Sec. 4 are roots of the standard overextended $`๐ข^{++}`$.
Proof: We shall explicitly write $`\widehat{\mathrm{\Lambda }}_i`$ in terms of the simple roots of $`๐ข^{++}`$. Let us remark that, by construction, $`๐ข^{++}`$ contains two affine $`๐ข`$, i.e. $`๐ข^+`$, whose real root system is formed by the roots of $`๐ข`$ plus $`nk_+`$, respectively $`nk_{}`$, ($`n๐`$). Clearly the following decomposition in roots is not all unique.
1. $`so(2N)`$
$`\widehat{\mathrm{\Lambda }}_{2m}`$ $`=`$ $`(\epsilon _1\epsilon _2+k_+)+(\epsilon _3\epsilon _4k_{})+\mathrm{}+(\epsilon _{2m1}\epsilon _{2m}k_{})`$ (86)
$`=`$ $`\mathrm{\Lambda }_{2m}+k_+(m1)k_{}`$
2. $`E_6`$
$`\widehat{\mathrm{\Lambda }}_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \underset{i6,7}{}}\epsilon _i+\epsilon _6+\epsilon _7+k_+]+{\displaystyle \frac{1}{2}}[\epsilon _1+\epsilon _2\epsilon _3\epsilon _4\epsilon _5+\epsilon _6+\epsilon _7\epsilon _82k_{}]`$ (87)
$`=`$ $`\mathrm{\Lambda }_3+k_+2k_{}`$
3. $`E_7`$
$`\widehat{\mathrm{\Lambda }}_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\epsilon _1\epsilon _2\epsilon _3\epsilon _4\epsilon _5\epsilon _6+\epsilon _7\epsilon _8+k_+]+(\epsilon _7\epsilon _82k_{})`$ (88)
$`=`$ $`\mathrm{\Lambda }_2+k_+2k_{}`$
$`\widehat{\mathrm{\Lambda }}_3`$ $`=`$ $`(\epsilon _3+\epsilon _7+k_+)+(\epsilon _4+\epsilon _72k_{})+(\epsilon _5\epsilon _8k_{})+(\epsilon _6\epsilon _82k_{})`$ (89)
$`=`$ $`\mathrm{\Lambda }_3+k_+5k_{}`$
$`\widehat{\mathrm{\Lambda }}_5`$ $`=`$ $`(\epsilon _3\epsilon _6+k_+)+(\epsilon _7+\epsilon _8k_{})=\mathrm{\Lambda }_5+k_+k_{}`$ (90)
4. $`E_8`$
$`\widehat{\mathrm{\Lambda }}_1`$ $`=`$ $`(\epsilon _i\epsilon _8+k_+)+(\epsilon _i\epsilon _8k_{})=\mathrm{\Lambda }_1+k_+k_{}(i8)`$
$`\widehat{\mathrm{\Lambda }}_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(\epsilon _1+\epsilon _2+\epsilon _3+\epsilon _4\epsilon _5\epsilon _6\epsilon _7\epsilon _8)+k_+]+(\epsilon _3\epsilon _8k_{})`$
$`+`$ $`(\epsilon _4\epsilon _82k_{})+(\epsilon _5\epsilon _83k_{})=\mathrm{\Lambda }_2+k_+6k_{}`$
$`\widehat{\mathrm{\Lambda }}_3`$ $`=`$ $`(\epsilon _3\epsilon _8+k_+)+(\epsilon _4\epsilon _82k_{})+(\epsilon _5\epsilon _83k_{})`$
$`+`$ $`(\epsilon _6\epsilon _84k_{})+(\epsilon _7\epsilon _85k_{})=\mathrm{\Lambda }_3+k_+14k_{}`$
$`\widehat{\mathrm{\Lambda }}_4`$ $`=`$ $`(\epsilon _4\epsilon _8+k_+)+(\epsilon _5\epsilon _82k_{})+(\epsilon _6\epsilon _83k_{})+(\epsilon _7\epsilon _84k_{})`$
$`=`$ $`\mathrm{\Lambda }_4+k_+9k_{}`$
$`\widehat{\mathrm{\Lambda }}_5`$ $`=`$ $`(\epsilon _5\epsilon _8+k_+)+(\epsilon _6\epsilon _82k_{})+(\epsilon _7\epsilon _83k_{})`$
$`=`$ $`\mathrm{\Lambda }_5+k_+5k_{}`$
$`\widehat{\mathrm{\Lambda }}_6`$ $`=`$ $`(\epsilon _6\epsilon _8+k_+)+(\epsilon _7\epsilon _82k_{})=\mathrm{\Lambda }_6+k_+2k_{}`$
$`\widehat{\mathrm{\Lambda }}_8`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(\epsilon _1\epsilon _2\epsilon _3\epsilon _4\epsilon _5\epsilon _6+\epsilon _7+\epsilon _8)+k_+]+(\epsilon _6\epsilon _8k_{})`$
$`+`$ $`(\epsilon _7\epsilon _82k_{})=\mathrm{\Lambda }_8+k_+3k_{}`$
Therefore, the considered combinations of the fundamental weights with the vectors of the Lorentzian lattices II$`{}_{k}{}^{}{}_{}{}^{1,1}`$ do belong to the root system of the overextended Lie algebras, even if the considered fundamental weights do not belong to the root system of the native Lie algebra <sup>8</sup><sup>8</sup>8One can recover the generators by the change $`\alpha E_\alpha `$ and $`\alpha +\beta [E_\alpha ,E_\beta ]`$.. It follows that the double non standard extensions (called previously affine extensions) are all subalgebras, of the same rank, of $`๐ข^{++}`$. Indeed it is easy to verify that the differences $`\beta _i\beta _j`$ ($`i,j;ij`$) of the simple roots $`\beta _i`$ of the non standard double extension of $`๐ข`$ do not belong to the root system , therefore satisfying the conditions of the Feingold-Nicolai theorem. In particular we have
###### Proposition 3
The indefinite Kac-Moody algebras of rank 10 described by the Dynkin-Kac diagrams, obtained by adding to the diagram of the affine algebra $`E_8`$, i.e. $`E_9`$, a dot, connected with a simple link to the $`j`$-th dot of $`E_9`$ ($`1j9`$), is a subalgebra of $`E_{10}`$.
The simple root systems of these subalgebra is formed by the simple roots of $`E_8`$, by a root $`\widehat{\mathrm{\Lambda }}_i`$ given by eq.(LABEL:eq.e8) and by the root $`\theta +ak_{}`$ where $`a`$ is a positive integer such that $`(\widehat{\mathrm{\Lambda }}_i,\theta +ak_{})=0`$.
One can naturally ask if an analogous theorem holds for $`E_{11}`$, that is if the triple non standard extensions of $`E_8`$ form Lorentzian algebras of rank 11, subalgebra of $`E_{11}`$. We have:
###### Proposition 4
The indefinite Kac-Moody algebras of rank 11 obtained by adding to the simple root system of the algebra $`E_8`$ three roots: $`\alpha _{+1}=\widehat{\mathrm{\Lambda }}_j`$, connected with a simple link to the $`j`$-th dot of $`E_8`$ ($`1j8,j7`$); $`\alpha _{+2}=\theta ak_{}`$, where $`a`$ is a positive integer such that ($`\alpha _{+2}\alpha _{+1}`$) = 0 and $`\alpha _{+3}=l_++l_{}k_{}`$, simply linked with $`\alpha _{+1}`$, is a subalgebra of $`E_{11}`$.
The proof is straightforward using the explicit expressions of $`\widehat{\mathrm{\Lambda }}_j`$, given in eq.(LABEL:eq.e8). Let us remark that these subalgebras do not have as subalgebra $`E_{10}`$. Below we give the simple roots systems of a set (not exhaustive) of rank eleven subalgebra $`E_{11}`$, which contains as subalgebra $`E_{10}`$. Let us consider the following simple root system of $`E_{11}`$: $`\alpha _i`$ ($`1i8`$) are the simple roots of $`E_8`$, $`\alpha _9=\alpha _0k_+`$, $`\alpha _{10}=k_++k_{}`$ and $`\alpha _{11}=(l_++l_{})k_+`$. We follow the convention of and, for the reader convenience, we explicitly write here the $`E_8`$ simple root system and the root system $`\mathrm{\Delta }`$:
$`\alpha _1={\displaystyle \frac{1}{2}}(\epsilon _1+\epsilon _8{\displaystyle \underset{j=2}{\overset{7}{}}}\epsilon _j)\alpha _i=\epsilon _i\epsilon _{i1}i=2,\mathrm{},7`$
$`\alpha _8=\epsilon _1+\epsilon _2h.r.:=\alpha _0=\epsilon _7+\epsilon _8`$ (92)
$$\mathrm{\Delta }=\{\frac{1}{2}(\pm \epsilon _1\pm \epsilon _2\pm \epsilon _3\pm \epsilon _4\pm \epsilon _5\pm \epsilon _6\pm \epsilon _7\pm \epsilon _8),\pm \epsilon _i\pm \epsilon _j\}$$
(93)
where the total number of $`+`$ signs (or $``$ signs) in the first expression is an even number. Let us consider the $`E_{10}`$ Dynkin-Kac diagram obtained by the $`E_{11}`$ diagram, deleting the dot corresponding to $`11`$-th simple root. Let us denote by $`E_{10}^{(j)}`$ the algebra of rank 11 whose Dynkin-Kac diagram is obtained by the $`E_{10}`$ diagram adding a dot (in the following denoted by +1) with a simple link to the $`j`$-th dot ($`1j10`$) and by $`\beta _i^{(j)}`$ ($`i=1,\mathrm{},10,+1`$) the simple roots of $`E_{10}^{(j)}`$. In the following we do not explicitly write the upper label $`j`$ in the roots. Clearly $`E_{10}^{(10)}=E_{11}`$ and, in this case, $`\beta _{+1}=\alpha _{11}`$ . We make the following (not unique) choice for the simple root system of $`E_{10}^{(j)}`$ ($`1j9`$):
* j = 9)
$`\beta _1=\alpha _1k_+\beta _i=\alpha _ii=2,\mathrm{},7,8`$
$`\beta _9=\epsilon _7+\epsilon _8+k_{}\beta _{10}=\alpha _0k_+\beta _{+1}=\alpha _{11}`$ (94)
* j = 7)
$`\beta _i=\alpha _ii=1,\mathrm{},6,8\beta _7=\alpha _7k_+`$
$`\beta _9=\alpha _{10}\beta _{10}=\alpha _{11}\beta _{+1}=\alpha _0`$ (95)
* j = 6)
$`\beta _i=\alpha _ii=1,\mathrm{},5,7,8\beta _6=\alpha _6+k_{}`$
$`\beta _9=\alpha _0\beta _{10}=\epsilon _7+\epsilon _6k_+\beta _{+1}=\alpha _{11}`$ (96)
* j = 5)
$`\beta _i=\alpha _ii=1,2,3,4,6,7,8\beta _5=\alpha _5+k_{}`$
$`\beta _9=\alpha _0\beta _{10}=\epsilon _5+\epsilon _8k_+\beta _{+1}=\alpha _{11}`$ (97)
* j = 4)
$`\beta _i=\alpha _ii=1,2,3,5,6,7,8\beta _4=\alpha _4+k_{}`$
$`\beta _9=\alpha _0\beta _{10}=\epsilon _4+\epsilon _8k_+\beta _{+1}=\alpha _{11}`$ (98)
* j = 3)
$`\beta _i=\alpha _ii=1,2,\mathrm{},7,8\beta _3=\alpha _3+k_{}`$
$`\beta _9=\alpha _0\beta _{10}=\epsilon _3+\epsilon _8k_+\beta _{+1}=\alpha _{11}`$ (99)
* j = 2)
$`\beta _i=\alpha _ii=1,3,\mathrm{},7,8\beta _2=\alpha _2+k_{}`$
$`\beta _9=\alpha _0\beta _{10}=\epsilon _1+\epsilon _8k_+\beta _{+1}=\alpha _{11}`$ (100)
* j = 1)
$`\beta _1={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=i}{\overset{8}{}}}\epsilon _i+k_{}\beta _2=\alpha _8\beta _i=\alpha _ii=3,\mathrm{},6\beta _8=\alpha _2`$
$`\beta _7=\alpha _7\beta _9=\epsilon _7+\epsilon _8\beta _{10}=\alpha _1+k_+\beta _{+1}=\alpha _{11}`$ (101)
* j = 8)
$`\beta _i=\alpha _ii=1,\mathrm{},7\beta _8=\alpha _8+k_{}`$
$`\beta _9=\alpha _0\beta _{10}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{8}{}}}\epsilon _ik_+\beta _{+1}=\alpha _{11}`$ (102)
It is easy to verify that the roots $`\beta _i`$ belong to the roots sistems of $`E_{11}`$,<sup>9</sup><sup>9</sup>9Actually they belong, except for $`\alpha _{11}`$, to the root system of affine $`E_8=E_9`$ which is a regular sub-algebra of $`E_{11}`$ while the differences $`\beta _i\beta _j`$ ($`i,j;ij`$) ($`|\beta _i\beta _j|^24`$) do not belong, therefore satisfying the conditions of the Feingold-Nicolai theorem. The Dynkin-Kac diagrams describing these algebras contain loops except for $`j=7`$. This algebra has been considered in , where it has been denoted $`EE_{11}`$, and it has been shown to be a subalgebra of $`E_{11}`$ by explicitly constructing the generators by commutation of the $`E_{11}`$ generators. The Dynkin-Kac diagrams for these algebras are easily drawn by adding to the $`E_{10}`$ diagram a dot simply connected with the $`j`$-th dot and then connecting, with a simple link, the following dots: $`j=9)\mathrm{\hspace{0.33em}\hspace{0.33em}7}10;j=6)\mathrm{\hspace{0.33em}\hspace{0.33em}1}10;j=5)\mathrm{\hspace{0.33em}\hspace{0.33em}6}10;j=4)\mathrm{\hspace{0.33em}\hspace{0.33em}5}10;j=3)\mathrm{\hspace{0.33em}\hspace{0.33em}4}10;j=2)\mathrm{\hspace{0.33em}\hspace{0.33em}8}10;j=1)\mathrm{\hspace{0.33em}\hspace{0.33em}7}10,810;j=8)\mathrm{\hspace{0.33em}\hspace{0.33em}1}10`$. Of course, these subalgebras do not exhaust the set of 11 dimensional indefinite Kac-Moody subalgebras of $`E_{11}`$. In particular we have not considered the subalgebras which do appear as invariant algebras with respect to an involution of the generators of $`E_{11}`$, see and .
## 6 Conclusions and future developments
In studying 4-extended Lie algebras, we have seen that Borcherds algebras seem to emerge naturally. This remark rises the question: which are the fingerprints of a theory which exhibits a symmetry under a Borcherds algebra? This question is indeed interesting on the light of the remark that many dualities have a group-theoretical origin in the Weyl group of the algebra. The Weyl group of the Borcherds algebra has peculiar properties as the reflection with respect to the imaginary vanishing roots is not defined. Some particular properties related to this kind of algebras have already been discussed in . The non-standard extension introduced in this paper have peculiar features, which deserve further investigation, on both their mathematical structure and their possible physical relevance. A classification of these algebras is beyond the aim of this paper, where we present only a few representative examples. As, however, very little is known on Lorentzian Kac-Moody algebras, we believe that any new information is interesting. In the cited literature on the physical role of the very extended Lie algebras, non-linear realizations of the indefinite Kac-Moody algebras are used. How a Chevalley realization of this algebra looks like? In a procedure to build up vertex realization of Lorentzian algebra with only a lattice II<sup>1,1</sup> has been proposed and applied to the very simple case of the overextended $`A_1`$ algebra. It seems possible to generalize that procedure to the triple extended Lie algebras. Moreover, it has also been argued by P. West that $`sl(32)`$ is contained in the Cartan invariant sub-algebra of $`E_{11}`$. At first sight the rank of $`sl(32)`$ is too large to be a sub-algebra, so it seems that very extended algebras, at least in the non linear realization, admit finite dimensional sub-algebras which naively could not be there. The investigation of the finite Lie subalgebras of the indefinite Kac-Moody algebras requires new methods beyond the very familiar ones used in the case of finite Lie algebra, which are essentially based on the Dynkin methods. This feature is not completely unrelated with the property, noted in , that the set of infinite dimensional sub-algebras of Lorentzian algebras is quite rich and surprising. We have illustrated this feature in Sec. 5, discussing a class of subalgebras of $`E_{11}`$, but it would be useful to dispose of techniques to build up explicitly or to identify classes of these sub-algebras or to dispose of further examples.
Acknowledgment \- We thank C. Helfgott, A. Kleinschmidt and I. Schnakenburg for pointing us a mistake in the previous version and for useful comments.
## Appendix A Some facts about the lattice II<sup>1,1</sup>
We review some basic facts about the lattice II<sup>1,1</sup>, which is the only Lorentzian even self-dual lattice in dimension 2. The points in this lattice can be described as the vectors:
$$(n,m)$$
(103)
with $`n,m`$ and Gram matrix $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, with eigenvalues $`\pm 1`$. In this way, the scalar product between two vectors $`a=(a_+,a_{})`$ and $`b=(b_+,b_{})`$ can be written as: $`ab=a_+b_{}a_{}b_+`$. We can take $`k_+(1,0)`$ and $`k_{}(0,1)`$ as basis vectors in II<sup>1,1</sup>; with this choice, we have:
$$k_\pm k_\pm =0,k_\pm k_{}=1.$$
(104)
All the vectors in II<sup>1,1</sup> can be written as $`v=pk_++qk_{}`$ (with $`p,q`$); in particular there are only two vectors of squared norm 2, $`\pm (k_++k_{})`$, but infinite vectors of positive ($`4`$) and negative ($`2`$) squared norm, being $`v^2=2pq`$.
## Appendix B Definition of Borcherds Algebras
The best way to think of a Borcherds algebra is to consider it as a generalization of a finite-dimensional simple Lie algebra. The definition is based on the Serre-Chevalley construction of finite-dimensional algebras; we follow .These algebras always have a symmetric matrix and their structure is very similar to that of ordinary Kac-Moody algebras, the only major difference is that generalized Kac-Moody algebras allow the presence of imaginary simple roots. Let $`A=(a_{i,j})`$ a $`n\times n`$ (real) symmetric matrix satisfying the following properties:
* $`a_{ii}=2`$ or $`a_{ii}0`$,
* $`a_{ij}0`$ if $`ij`$,
* $`a_{ij}`$ if $`a_{ii}=2`$.
Then the Borcherds algebra $`๐ข(A)`$ associated with the Cartan matrix $`A`$ is the Lie algebra given by the following generators and relations.
$`3n`$ Generators: $`e_i,f_i`$ and $`h_i`$
Relations:
* $`[h_i,h_j]=0,`$
* $`[e_i,f_j]=\delta _{ij}h_i`$,
* $`[h_i,e_j]=a_{ij}e_j,[h_i,f_j]=a_{ij}e_j`$,
* $`e_{ij}:=(\text{ad }e_i)^{1a_{ij}}e_j=0`$, $`f_{ij}:=(\text{ad }f_i)^{1a_{ij}}f_j=0`$ if $`a_{ii}=2`$ and $`ij`$,
* $`e_{ij}:=[e_i,e_j]=0`$, $`f_{ij}:=[f_i,f_j]=0`$ if $`a_{ii}0`$, $`a_{jj}0`$ and $`a_{ij}=0`$.
The elements $`h_i`$ form a basis for an abelian subalgebra of $`๐ข(A)`$, called Cartan subalgebra $`(A)`$; as it happens for Kac-Moody algebras, $`๐ข(A)`$ has the triangular decomposition:
$$๐ข(A)=๐ฉ_{}(A)๐ฉ_+$$
(105)
and has many of the properties of the usual Kac-Moody algebras (real and imaginary roots, etc.). In particular, in this paper, we have considered Borcherds algebras with just one imaginary simple root (with squared norm 0), which we have added by hand. A rank 2 Borcherds algebra in the lattice II<sup>1,1</sup> can be constructed as follows: the Cartan matrix is given by
$$A=\left(\begin{array}{cc}0& 1\\ 1& 2\end{array}\right)$$
(106)
A possible choice for the simple roots is
$$\alpha _1=k_+,\alpha _2=(k_++k_{}),$$
(107)
with Weyl vector $`\rho =k_+`$ (defined by $`(\rho ,\alpha _i)=1/2(\alpha _i,\alpha _i)`$). |
warning/0506/hep-th0506256.html | ar5iv | text | # 1 Introduction and Conclusion
## 1 Introduction and Conclusion
Recently there has been a considerable interest in BPS solitons in Higgs phase of supersymmetric Yang-Mills theories with eight supercharges. Almost all known BPS objects, like magnetic flux vortices, magnetic monopoles, domain walls, and instantons, have appeared here, sometimes with a bit of twist. These theories can allow many degenerate vacua which can be interpolated by domain walls. With broken $`U(1)`$ gauge theories, one can have magnetic flux vortex. One of the most interesting features has been that there can be magnetic monopoles which appear as beads on vortex strings.
These BPS objects can be interpreted in a simple manner from D-brane point of view. Especially a simple but rich picture appears with $`N`$ parallel D3 branes and $`N_f`$ D7 branes. In this setting, one can have $`N=2`$ supersymmetric $`U(N)`$ gauge theory with $`N_f`$ matter hypermultiplets in the fundamental representation and a single adjoint hypermultiplet. One can add Fayet-Iliopoulos(FI)-terms and mass term for the matter hypermultiplet without breaking the supersymmetry. There are considerable work done along this line to represent the configurations in brane picture.
In this work, we focus on BPS equations, dyonic 1/4 BPS, and 1/8 BPS solutions. In addition, we explore BPS vortex equations when $`N=2,N_f=1`$ and found the cases where there are no vortex solutions of unit or double vorticity.
By studying the known bosonic BPS equations, we found that there are two parameter family of 1/8 BPS equations in 3+1 dimension modulo spatial rotation and $`SU(2)_R\times U(1)_R`$. The FI term breaks $`SU(2)_R`$ to $`U(1)`$ and the mass terms for matter hypermultiplet breaks $`U(1)_R`$ completely. One would expect more general BPS configurations in this setting.
Dyonic objects mean objects carrying โelectricโ charge. Of course there will be no isolated electric charge in the Higgs phase due to screening. Electrically charged solitons could be interpreted as composites of soliton with fundamental strings whose ends carry electric charge. In Higgs phase the electric charge is neutralized by electric charge carried by the Higgs field. As the Higgs fields carry global flavor charge, the conserved flavor charge instead of the total electric flux would appear in the BPS energy formula. Besides dyonic monopoles, we show that dyonic domain walls as well as dyonic composites of domain wall-monopole-vortex are also possible. When parallel D7 branes are not lying on a single line in their transverse space, dyonic BPS configurations which make web-like structures are also possible. These dyonic solutions could be interpreted as the excitations in phase moduli of BPS objects and they belong to 1/4 BPS states.
We also look for BPS solutions preserving 1/8 of eight supersymmetries. By exploring a small perturbation of a homogeneous 1/4 BPS configuration in 3+1 dimensional theories, we argue that there may be no 1/8 BPS configurations satisfying the BPS equations. However we find easily 1/8 BPS configurations in a theory with product gauge group $`U(1)\times U(1)`$ with bi-fundamental and fundamental matter field. In this analysis, the recently discoverd bound states of monopoles and domain walls play some role.
The key aspect here is that the FI parameters breaks the $`SU(2)_R`$ symmetry of eight supersymmetric 5+1 dimensional theory. For a single $`U(1)`$ gauge group, one can use the broken R-symmetry to choose a single direction in $`SU(2)_R`$ space. However with product gauge groups, the FI-parameters cannot be rotated to a single direction in general. This is what allows the presence of 1/8 BPS configurations possible.
As there are multi BPS vortex string configuration in $`U(1)`$ theory with $`N_f=1`$, we may expect there are BPS vortex string configuration in $`U(2)`$ theory with $`N_f=1`$. While there exist degenerate supersymmtric vacua, we will show that classically there exists no BPS vortex configuration with unit and double magnetic flux. We argue that this may imply that there exists no BPS vortex solitons of finite magnetic flux in the theory.
One interesting direction to explore further is the interaction between domain walls and monopoles. (See also a recent work by Sakai and Tong\***.) In string picture, parallel D1 and D3 branes are attracted to each other. This is not apparent from the energy argument of a BPS monopole-vortex-domain composition. The moduli space of domain wall-monopole separation should be analyzed carefully to resolve the question.
Another direction is to study the moduli space dynamics of magnetic monopoles and domain walls when some of nonabelian gauge symmetry is restored. It would be interesting to see whether there exists a similar restoration of symmetry in the moduli space dynamics..
Finally, all BPS solutions we study here have extended structures with infinite energy. There may be finite action BPS solitons in the theory. Especially it may be possible to have finite energy (dyonic) instantons in $`R^3\times S^1`$ (noncommutative) space, which do not have diverging gauge flux.
The plan of this paper is as follows. In Sec.2, we describe 5+1 dimensional supersymmetric Yang-Mills theories and find supersymmetric Lagrangian and its vacuum structure. In Sec.3, we find two parametered BPS equations, especially 1/8 BPS equations. In Sec.4, we study dyonic solutions. In Sec.5, we study 1/8 BPS configurations and find BPS configurations with product gauge group. In Sec.6, we show that there exists no BPS vortex solitons of unit and double magnetic flux when $`N=2`$ and $`N_f=1`$.
Note added : In the early stage of the draft of our paper, we came to know that the authors of Ref. have worked on the classification of 1/8 BPS equations of the similar model we considered.
## 2 Six Dimensional Case
The vector multiplet of super Yang-Mills theory of $`U(N)`$ gauge group with eight supersymmetries in six dimensions is made of $`A_M,\lambda _i(i=1,2),๐^a`$, which are hermitian $`N\times N`$ matrix valued fields. The gaugino field $`\lambda _i,i=1,2`$ is made of two eight component spinors satisfying both chirality and symplectic Majorana conditions
$$\mathrm{\Gamma }^6\lambda _i=\lambda _i(i=1,2),\lambda _i=(i\sigma ^2)_{ij}B(\lambda _j^{})^T$$
(2.1)
where $`B`$ is a matrix such that $`B\mathrm{\Gamma }^MB^1=(\mathrm{\Gamma }^M)^{}`$. Due to this constraint, there are only four physical degrees of freedom in gaugino spinor. Our choice of six dimensional Gamma matrices are
$`\mathrm{\Gamma }^0=1_2i\sigma ^3\sigma ^1,\mathrm{\Gamma }^a=\sigma ^a\sigma ^1\sigma ^1(a=1,2,3)`$
$`\mathrm{\Gamma }^4=1_2\sigma ^2\sigma ^1,\mathrm{\Gamma }^5=1_21_2\sigma ^2`$ (2.2)
In addition, $`\mathrm{\Gamma }^6=\mathrm{\Gamma }^0\mathrm{\Gamma }^1\mathrm{}\mathrm{\Gamma }^5=1_21_2\sigma ^3`$. With the above choice,
$$B=i\sigma ^21_2\sigma ^3.$$
(2.3)
The Lagrangian for the gauge multiplet is
$$_1=\mathrm{tr}\left(\frac{1}{4}F_{MN}F^{MN}\frac{i}{2}\overline{\lambda }_i\mathrm{\Gamma }^MD_M\lambda _i+\frac{1}{2}(๐^a)^2\right)$$
(2.4)
The supersymmetric transformation becomes
$`\delta A_M=i\overline{\lambda }_i\mathrm{\Gamma }_Mฯต_i`$ (2.5)
$`\delta \lambda _i={\displaystyle \frac{1}{2}}F_{MN}\mathrm{\Gamma }^{MN}ฯต_i+i๐^a\sigma _{ij}^aฯต_j`$ (2.6)
$`\delta ๐^a=\overline{ฯต}_i\sigma _{ij}^a\mathrm{\Gamma }^ID_I\lambda _j`$ (2.7)
where the supersymmetric parameter $`ฯต_i`$ is also a chiral spinor and satisfies the symplectic Majorana condition. The Lagrangian and supersymmetric transformation are compatible with the symplectic Majorana condition. The above Lagrangian is invariant under $`SU(2)_R`$ transformation, under which $`\lambda _i`$ and $`๐^a`$ belong to the fundamental and adjoint representations, respectively.
The Lagrangian for an adjoint hypermultiplet $`y_i(i=1,2),\chi `$ where the matter spinor is anti-chiral $`\mathrm{\Gamma }^6\chi =\chi `$, is
$$_2=\mathrm{tr}\left(\frac{1}{2}D_M\overline{y}_iD^My_i+\frac{1}{2}๐^a\sigma _{ij}^a[\overline{y}_j,y_i]i\overline{\chi }\mathrm{\Gamma }^MD_M\chi +\overline{\lambda }_i[\overline{y}_i,\chi ]\overline{\chi }[y_i,\lambda _i]\right)$$
(2.8)
where $`D_My_i=_My_ii[A_M,y_i]`$. Here $`y_i(i=1,2)`$ is a doublet under $`SU(2)_R`$ and $`\chi `$ is a singlet. The supersymmetric transformation is
$$\delta \overline{y}_i=2i\overline{\chi }ฯต_i,\delta \chi =D_My_i\mathrm{\Gamma }^Mฯต_i$$
(2.9)
The matter hypermultiplets $`q_{fi},\psi _f`$ with flavor index $`f=1,\mathrm{},N_f`$ belong to the fundamental representation $`\overline{N}`$ of the gauge group $`U(N)`$. As in the adjoin hypermultiplet, the matter spinor field is anti-chiral. The Lagrangian for the matter multiplet is
$$_3=\mathrm{tr}\left(\frac{1}{2}D_M\overline{q}_{fi}D^Mq_{fi}+\frac{1}{2}๐^a\sigma _{ij}^a\overline{q}_{fj}q_{fi}i\overline{\psi }_f\mathrm{\Gamma }^MD_M\psi _f+\overline{\lambda }_i\overline{q}_{fi}\psi _f\overline{\psi }_fq_{fi}\lambda _i\right)$$
(2.10)
where $`D_Mq_{fi}=_Mq_{fi}+iq_{fi}A_M`$. The supersymmetric transformation is
$$\delta \overline{q}_{fi}=2i\overline{\psi }_fฯต_i,\delta \psi _f=D_Mq_{fi}\mathrm{\Gamma }^Mฯต_i$$
(2.11)
The above Lagrangians are invariant under the $`SU(2)_R`$ symmetry. For a theory with abelian gauge group, one can add the Fayet- Iliopoulos term
$$_{FI}=\frac{1}{2}\mathrm{tr}(\zeta ^a๐^a).$$
(2.12)
If the gauge group is a product group, there would be FI-terms for each independent $`U(1)`$ theory. The FI parameters $`\zeta ^a`$ breaks the $`SU(2)_R`$ symmetry explicitly and so one can use $`SU(2)_R`$ symmetry to rotate them to be
$$\zeta ^1=0,\zeta ^2=0,\zeta ^3=v^2$$
(2.13)
with $`v0`$. We will use both $`\zeta ^a`$ and parameter $`v`$. The $`๐^a`$ field is not dynamical and its field equation leads to
$$๐^a=\frac{e^2}{2}\left\{\zeta ^a\sigma _{ij}^a\left([\overline{y}_j,y_i]+\overline{q}_{fj}q_{fi}\right)\right\}$$
(2.14)
The dimensional reduction to 3+1 dimension induces additional $`U(1)_R`$ symmetry which is a rotation under two reduced space. The dimensional reduction with Scherk-Schwartz mechanism induces two mass parameters $`m_f,m_f^{}`$ for each flavor matter multiplet along the reduced space. If $`x^4,x^5`$ is reduced, then
$$D_4q_{fi}=iq_{fi}(A_4m_f),D_5q_{fi}=iq_{fi}(A_5m_f^{}).$$
(2.15)
This theory with $`U(N)`$ gauge group has a simple D-brane interpretation. It is a Yang-Mills theory on $`N`$ parallel D3 branes near $`N_f`$ D7 branes whose transverse location at $`x^4,x^5`$ is given by the mass parameter. The location of D3 branes along $`x^4,x^5`$ direcion id given by the vacuum expectation value of adjoint scalars $`A_4,A_5`$. The location of D3 branes along transverse 4 directions in D7 branes would be decided by the expectation value of $`y_i`$. The dimensional reduction to $`4+1`$ dimension is a bit simpler with only one mass parameter and no additional R-symmetry. The D-brane interpretation could be D4-D8 system.
One of the vacuum condition $`๐^a=0`$ is the ADHM condition of $`N`$ instantons on $`U(N_f)`$ gauge theory of noncommutative four space. The scalar fields are denoting the separation and size of instantons. As D3 branes act as instantons on $`D7`$ branes, one can see that the vacuum moduli space modulo gauge transformation is the moduli space of instantons when the mass parameters are turned off. With the mass parameters turned on, every D3 brane should lie on some D7 brane at the ground state. Thus, every eigenvalue pair of expectation value of $`(A_4,A_5)`$, which is diagonal at the vacuum, should coincide with $`(m_f,m_f^{})`$ for some $`f`$.
One of the simplest vacua appears when $`N=N_f`$ and all the eigenvalue pair of $`A_4,A_5`$ are distinct, such that there is only one D3 brane for each D7 brane. It is the so-called color-flavor locking phase, where the matter field will have a Higgs condensation $`q_{f1}_{vacuum}=v`$ and the gauge symmetry plus the flavor symmetry is spontaneously broken down to unbroken $`U(1)^N`$ global symmetry.
When $`N=2,N_f=1`$, the vacuum moduli space would be that of two $`U(1)`$ instantons on noncommutative four space, which is the so-called Eguchi-Hanson space. In this case $`y_i`$ does have intrinsic nonabelian components and the gauge group $`U(2)`$ is spontaneously broken to global $`U(1)`$ symmetry.
## 3 BPS Equations
Classically a BPS field configuration is a bosonic field configuration which leaves some of the supersymmetry invariant. We consider now the supersymmetric transformation to obtain the BPS equations. Inspired by the bosonic BPS equations, we rewrite the supersymmetric transformation of the gaugino field as
$`\delta \lambda _i`$ $`=`$ $`\mathrm{\Gamma }^{12}((F_{12}F_{34}\mathrm{\Gamma }^{1234})ฯต_ii๐^3\mathrm{\Gamma }^{12}\sigma _{ij}^3ฯต_j)+\mathrm{\Gamma }^{23}((F_{23}F_{14}\mathrm{\Gamma }^{1234})ฯต_i.`$ (3.16)
$`.i๐^1\mathrm{\Gamma }^{23}\sigma _{ij}^1ฯต_j)+\mathrm{\Gamma }^{31}((F_{31}F_{24}\mathrm{\Gamma }^{1234})ฯต_ii๐^2\mathrm{\Gamma }^{31}\sigma _{ij}^2ฯต_j)`$
$`+\mathrm{\Gamma }^{\mu 0}(F_{\mu 0}F_{\mu 5}\mathrm{\Gamma }^{05})ฯต_i+F_{05}\mathrm{\Gamma }^{05}ฯต_i`$
As $`\mathrm{\Gamma }^4ฯต_i=\mathrm{\Gamma }^{123}\mathrm{\Gamma }^{05}ฯต_i`$, the adjoint spinor transformation is written as
$$\delta \chi =\mathrm{\Gamma }^{123}\left(D_1y_i\mathrm{\Gamma }^{23}+D_2y_i\mathrm{\Gamma }^{31}+D_3y_i\mathrm{\Gamma }^{12}+D_4y_i\mathrm{\Gamma }^{05}\right)ฯต_i+\mathrm{\Gamma }^0(D_0y_iD_5y_i\mathrm{\Gamma }^{05})ฯต_i$$
(3.17)
The spinor in fundamental hypermultiplet transforms as
$`\delta \psi _f`$ $`=`$ $`\mathrm{\Gamma }^{123}\left(D_1q_{fi}\mathrm{\Gamma }^{23}+D_2q_{fi}\mathrm{\Gamma }^{31}+D_3q_{fi}\mathrm{\Gamma }^{12}+D_4q_{fi}\mathrm{\Gamma }^{05}\right)ฯต_i`$ (3.18)
$`+\mathrm{\Gamma }^0(D_0q_{fi}D_5q_{fi}\mathrm{\Gamma }^{05})ฯต_i.`$
We want find some supersymmetric parameter $`ฯต_i`$ such that $`\delta \lambda _i,\delta \chi ,\delta \psi _f`$ remain zero. On eight independent parameters of spinor $`ฯต_i`$, we impose three independent conditions (In the case of $`N=2`$ NLSM, see .),
$$\mathrm{\Gamma }^{05}ฯต_i=\eta ฯต_i,\mathrm{\Gamma }^{12}\sigma _{ij}^3ฯต_j=i\alpha ฯต_i,\mathrm{\Gamma }^{23}\sigma _{ij}^1ฯต_j=i\beta ฯต_i,$$
(3.19)
with $`\alpha ,\beta ,\eta `$ take $`\pm 1`$ independently. Since $`\mathrm{\Gamma }^0\mathrm{\Gamma }^1\mathrm{}\mathrm{\Gamma }^5=1`$ for chiral $`ฯต_i`$, these conditions imply that
$$\mathrm{\Gamma }^{31}\sigma _{ij}^2ฯต_i=i\alpha \beta ฯต_i,\mathrm{\Gamma }^{1234}ฯต_i=\eta ฯต_i.$$
(3.20)
These are conditions on eight independent Majorana parameters in the spinor $`ฯต_i`$, as they are compatible with the symplectic Majorana condition. If we impose any one of the conditions, the number of independent SUSY parameters would be reduced by one half to four of the original value. If we impose any two of them, the number of independent SUSY parameters are reduced to two or 1/4 of the original one. If we impose all three of them, the number of independent parameters is reduced to one, 1/8 of the original value.
One can obtain different conditions by six dimensional Lorentz transformations and $`SU(2)_R`$ transformations. In reduction to $`3+1`$ dimensions, only nontrivial ones modulo remaining symmetries is the rotation between the remaining coordinates and the reduced coordinates. In the reduction to $`3+1`$ dimensions of coordinate $`x^0,x^1,x^2,x^3`$, the above condition can be generalized to new spinor conditions with two parameters,
$`\mathrm{\Gamma }^0(\mathrm{\Gamma }^5\mathrm{cos}\theta +\mathrm{\Gamma }^3\mathrm{sin}\theta )ฯต_i=\eta ฯต_i,\mathrm{\Gamma }^1(\mathrm{\Gamma }^2\mathrm{cos}\phi +\mathrm{\Gamma }^4\mathrm{sin}\phi )\sigma _{ij}^3ฯต_{ij}=i\alpha ฯต_i,`$
$`(\mathrm{\Gamma }^2\mathrm{cos}\phi +\mathrm{\Gamma }^4\mathrm{sin}\phi )(\mathrm{\Gamma }^3\mathrm{cos}\theta \mathrm{\Gamma }^5\mathrm{sin}\theta )\sigma _{ij}^1ฯต_j=i\beta ฯต_i`$ (3.21)
This implies that
$$(\mathrm{\Gamma }^3\mathrm{cos}\theta \mathrm{\Gamma }^5\mathrm{sin}\theta )\mathrm{\Gamma }^1\sigma _{ij}^2ฯต_j=i\alpha \beta ฯต_i,\mathrm{\Gamma }^{124}(\mathrm{\Gamma }^3\mathrm{cos}\theta +\mathrm{\Gamma }^5\mathrm{sin}\theta )ฯต_i=\eta ฯต_i,$$
(3.22)
Note also $`D_4q_{fi}=iq_{fi}(A_4m_f)`$ and $`D_5q_{fi}=iq_{fi}(A_5m_f^{})`$. In reduction to $`4+1`$, we can put $`\phi =0`$ as it is a part of four dimensional spatial rotation.
We use the generalized spinor condition (3.21) to find the BPS equations satisfied by the bosonic configurations for the minimum amount 1/8 of the original supersymmetries. For any vector with spatial indices, we introduce barred indices so that
$`V_{\overline{1}}=V_1,V_{\overline{2}}=V_2\mathrm{cos}\phi +V_4\mathrm{sin}\phi ,V_{\overline{3}}=V_3\mathrm{cos}\theta V_5\mathrm{sin}\theta ,`$
$`V_{\overline{4}}=V_4\mathrm{cos}\phi V_2\mathrm{sin}\phi ,V_{\overline{5}}=V_5\mathrm{cos}\theta +V_3\mathrm{sin}\theta `$ (3.23)
From $`\delta \lambda _i=0`$, we get the gauge field part of BPS equations,
$`F_{0\overline{5}}=0,F_{\overline{\mu }0}\eta F_{\overline{\mu }\overline{5}}=0(\mu =1,\mathrm{}4),F_{1\overline{2}}\eta F_{\overline{3}\overline{4}}+\alpha ๐^3=0,`$
$`F_{\overline{2}\overline{3}}\eta F_{1\overline{4}}+\beta ๐^1=0,F_{\overline{3}1}\eta F_{\overline{2}\overline{4}}\alpha \beta ๐^2=0,`$ (3.24)
From $`\delta \chi =0`$ and $`\delta \psi _f=0`$, we also obtain
$`\beta D_1y_j\sigma _{ji}^1\alpha \beta D_{\overline{2}}y_j\sigma _{ji}^2+\alpha D_{\overline{3}}y_j\sigma _{ji}^3i\eta D_{\overline{4}}y_i=0,`$
$`D_0y_i\eta D_{\overline{5}}y_i=0,D_0q_{fi}\eta D_{\overline{5}}q_{fi}=0,`$
$`\beta D_1q_{fj}\sigma _{ji}^1\alpha \beta D_{\overline{2}}q_{fj}\sigma _{ji}^2+\alpha D_{\overline{3}}q_{fj}\sigma _{ji}^3i\eta D_{\overline{4}}q_{fi}=0`$ (3.25)
These are the BPS equations for 1/8 BPS configurations. The BPS equations preserving more supersymmetry can be obtained by imposing additional conditions to the above BPS equations. For example, 1/4 BPS configurations satisfy two sets of 1/8 BPS equations with, say, both $`\alpha =1`$ and $`\alpha =1`$. There is also a Gauss law constraint for the BPS configurations,
$$\frac{1}{e^2}\underset{\mu =0}{\overset{5}{}}D_{\overline{\mu }}F_{\overline{\mu }0}\frac{i}{2}([\overline{y}_i,D_0y_i][D_0\overline{y}_i,y_i])\frac{i}{2}(\overline{q}_{fi}D_0q_{fi}D_0\overline{q}_{fi}q_{fi})=0.$$
(3.26)
Using the BPS equation, the central charge for the BPS energy bound can be found to be
$`Z`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^3x\mathrm{tr}({\displaystyle \frac{\eta }{2e^2}}F_{\overline{\mu }\overline{\nu }}\stackrel{~}{F}_{\overline{\mu }\overline{\nu }}\alpha \zeta ^3(F_{1\overline{2}}\eta F_{\overline{3}\overline{4}})\beta \zeta ^1(F_{\overline{2}\overline{3}}\eta F_{1\overline{4}}).`$ (3.27)
$`.+\alpha \beta \zeta ^2(F_{\overline{3}1}\eta F_{\overline{2}\overline{4}}))+\eta m_fQ_f\mathrm{cos}\theta +\eta T_{03}\mathrm{sin}\theta +Z^{}`$
where $`\mu ,\nu =1,2,3,4`$ and $`\stackrel{~}{F}_{\mu \nu }=\frac{1}{2}ฯต_{\mu \nu \rho \sigma }F_{\rho \sigma }`$. After dimensional reduction to (3+1) dimensions, $`D_4y_i=i[A_4,y_i]`$ and $`D_4q_{fi}=iq_{fi}(A_4m_f)`$, and so $`F_{14}=D_1A_4`$ and $`F_{45}=i[A_4,A_5]`$. The charge $`Q_f`$ is the one carried by the $`f`$โth-flavor matter field,
$$Q_f=\frac{i}{2}d^3x\mathrm{tr}(\overline{q}_{fi}D_0q_{fi}D_0\overline{q}_{fi}q_{fi}),$$
(3.28)
and $`T_{03}`$ is the linear momentum along $`x^3`$ direction,
$$T_{03}=\frac{1}{2}d^3x\mathrm{tr}\left(\frac{1}{e^2}\underset{\mu =1,2,4,5}{}F_{\mu 0}F_{\mu 3}+(D_0\overline{y}_iD_3y_i+D_3\overline{y}_iD_0y_i)+(D_0\overline{q}_{fi}D_3q_{fi}+D_3\overline{q}_{fi}D_0q_{fi})\right).$$
(3.29)
The boundary term $`Z^{}`$ is given by
$$Z^{}=d^3x\left(\eta _i\mathrm{tr}(F_{i0}A_5)\mathrm{cos}\theta +\mathrm{}\right),$$
(3.30)
where $`\mathrm{}`$ indicates the terms quadratic in matter fields and expected to have zero boundary contribution in both Coulomb and Higgs phases. The first part would have nontrivial contribution in the Coulomb phase where there would be nontrivial electric field.)
The above BPS equations and the energy bound are complicated functions of two parameters $`\phi `$ and $`\theta `$. For example, a complication arises as
$$F_{\overline{3}\overline{4}}=F_{34}\mathrm{cos}\theta \mathrm{cos}\phi F_{32}\mathrm{cos}\theta \mathrm{sin}\phi F_{54}\mathrm{sin}\theta \mathrm{cos}\phi +F_{52}\mathrm{sin}\theta \mathrm{sin}\phi $$
(3.31)
Using the un-bared coordinate indices, we note that the first term of the above expression can be expressed as
$$F_{\overline{\mu }\overline{\nu }}\stackrel{~}{F}_{\overline{\mu }\overline{\nu }}=F_{\mu \nu }\stackrel{~}{F}_{\mu \nu }\mathrm{cos}\theta +4(F_{12}F_{45}+F_{24}F_{15}+F_{41}F_{25})\mathrm{sin}\theta $$
(3.32)
There are also the boundary terms depending on quark fields, which is supposed to make vanishing contributions almost all cases.
Once we fix $`\zeta ^a=v^2\delta _3^a`$, which is possible for the theories of $`U(N)`$ gauge group but not for those with product gauge group like $`U(1)\times U(1)`$, the BPS energy does not depends on the choice of the parameter $`\beta `$. This means that 1/4 BPS configurations defined by $`\alpha `$ and $`\eta `$ parameters could have 1/8 BPS excitations without generating additional energy, which is strange. Indeed we see that this is impossible in some simple case studied in Sec.5.
We can choose two parameters $`\theta ,\phi `$ to be arbitrary. If we fix $`\zeta ^a`$, we no longer have the freedom of $`SU(2)_R`$ transformation, and the parameters $`\theta ,\phi `$ become physically meaningful. One typical cases of BPS equations would be when $`\theta =\phi =0`$. In this case, the barred spacial indices become the un-barred ones and $`_4=_5=0`$. The other extreme may be when $`\theta =\phi =\pi /2`$. In this case the time dependent part becomes $`F_{03}=0,(D_0\eta D_3)anyfield=0`$, and
$`\eta F_{12}+i[A_4,A_5]+\beta ๐^1=0,D_1A_4+\eta D_2A_5+\alpha ๐^\mathrm{๐}=0,`$
$`D_1A_5\eta D_2A_4\alpha \beta ๐^2=0,`$ (3.33)
$`\left(\beta D_1y_j\sigma _{ji}^1+i\eta D_2y_i+i\alpha \beta [A_4,y_j]\sigma _{ji}^2+i\alpha [A_5,y_j]\sigma _{ji}^3\right)=0,`$
$`\left(\beta D_1q_{fj}\sigma _{ji}^1+i\eta D_2q_{fi}i\alpha \beta q_{fj}(A_4m_f)\sigma _{ji}^2i\alpha q_{fj}(A_5m_f^{})\sigma _{ji}^3\right)=0`$
We know quite a bit of the topological objects of the theories in $`\theta =\phi =0`$. The simplest object is a 1/2 BPS vortex soliton along $`x^3`$ direction in $`U(1)`$ theory with $`N_f=1`$. It satisfies the BPS equation with $`\beta =1`$,
$$2F_{12}=v^2|q_1|^2,(D_1iD_2)q_1=0,$$
(3.34)
where $`y_i=0,q_2=0`$, dropping the flavor index. Especially a unit flux vortex has a vortex tension $`T_v=\pi v^2`$. This could be regarded as a D1 string on a single D3 brane in a single D7 brane. The next simplest object is a 1/2 BPS domain wall parallel to $`(x^1,x^2)`$ plane. With $`N=1`$ and $`N_f=2`$ with two different $`m_f`$ along $`x^4`$ direction, the 1/2 BPS equations with $`\alpha \beta =1`$ becomes
$$2_3A_4=v^2\underset{f=1}{\overset{2}{}}|q_{f1}|^2,_3q_{f1}=q_{f1}(A_4m_f),$$
(3.35)
where $`y_i=0,q_{f2}=0`$, $`m_1<m_2`$, and $`A_5=0`$. The $`A_4`$ interpolates between $`m_1`$ and $`m_2`$. It describes the D3 brane on first D7 brane interpolating to the second D7 brane. The wall tension is $`T_{12}=\pi v^2(m_2m_1)`$.
More complicated object is a 1/4 BPS configuration made of magnetic monopole beads in a vortex flux tube. With $`N=N_f=2`$ and in the color-flavor locking phase with $`m_1<m_2`$ and $`A_5=m_f^{}=0`$, the D1 string on the first D3-D7 branes can interpolate to the second D3-D7 branes. The D1 string connecting two D3 branes appears as a magnetic monopole. In the Higgs phase, the magnetic flux is confined to flux string and so the 1/4 BPS object is made of two vortices emerging opposite to the magnetic monopole, where two $`U(1)`$โs of $`U(2)`$ flux are carried to opposite direction. The composite has the energy of a simple sum of vortex tension and monopole mass.
Most complicated 1/4 BPS object is a composite made of vortex and domain walls, which also allows some magnetic monopoles. With $`\beta =1,\alpha =1`$, from the BPS energy one notices that with positive $`\mathrm{tr}F_{12}`$ and $`\mathrm{tr}F_{34}`$, which means postive vortex flux and domain wall charge where $`A_4`$ is increasing, there is negative instanton energy, or monopole energy. This is the so-called bound energy of vortex-domain wall. If a vortex terminates at the domain wall, the wall shape gets deformed in large distance away from the contact point. The detail has been also studied recently. Of course one can add additional monopole kink to this vortex-domain wall junction, which carries the positive monopole energy. In some cases, the magnetic monopole can pass the domain wall. When a vortex penetrating a domain wall is deformed so that the contact points at the both sides of the domain wall do not coincide to the same point, the monopole could not pass the domain wall due to the energy consideration, which means that there could be repulsive potential at the domain wall. It would be interesting to find whether our conjecture is true.
A typical solution of the BPS equations of $`\theta =\phi =\pi /2`$ would be the 1/4 BPS domain wall junction with $`N=1,N_f=3`$. Suppose that the three complex masses $`m_f+im_f^{}`$ lie on vertices of an equitriangle so that $`m_f+im_f^{}=me^{2\pi if/3}`$ with $`f=1,2,3`$. The BPS equation would be give by (3) with $`\zeta ^3=v^2`$ and the wall junction would lie on on $`x^1,x^2`$ plane with $`x^3`$ translation invariance. The ansatiz is that $`y_i=0,q_{f2}=0,A_1=A_2=A_3=0,_3=0`$ and the BPS equation becomes
$`_1A_4+\eta _2A_5={\displaystyle \frac{\alpha }{2}}(v^2|q_{f1}|^2),_1A_5\eta _2A_4=0`$ (3.36)
$`_1q_{f1}\alpha (A_4m_f)q_{f1}=0,\eta _2q_{f1}\alpha (A_5m_f^{})q_{f1}=0`$ (3.37)
The web of wall solutions of this type in a bit more complicated setting has been also studied recently.
## 4 Lorentz Boosted, or Dyonic Solutions
For the BPS configurations, the time dependent part can be solved with
$$A_0=\eta (A_5\mathrm{cos}\theta +A_3\mathrm{sin}\theta ),_0q_{fi}\eta (_3\mathrm{sin}\theta im_f^{}\mathrm{cos}\theta )q_{fi}=0$$
(4.38)
while $`(_0_3\mathrm{sin}\theta )=0`$ for any field in the adjoint representation. One can see that it is a Lorentz boost along $`x^3`$ axis with velocity $`v=\mathrm{sin}\theta `$ when $`|\theta |<\pi /2`$. However, the $`\theta =\pi /2`$ case is still physically distinct as it cannot be obtained through finite boost. The Gauss law is also equivalently Lorentz boosted version. This matches with the energy being increased with $`T_{03}v=๐ช(v^2)`$ for small $`v`$ as $`T_{03}`$ itself is linear in $`v`$ for small $`v`$. For the domain wall junctions, $`T_{03}=0`$ with $`\theta =\pi /2`$ due to the $`x^3`$ translation invariance of the configuration. Thus one cannot boost them along $`x^3`$, but may be able to put some massless wave along $`x^3`$ without breaking the supersymmetry further.
When $`\theta =0`$, $`A_0=\eta A_5`$ and the all adjoint fields are time-independent and $`_0q_{fi}+i\eta m_f^{}q_{fi}=0`$. The $`f`$-th flavor charge becomes
$`Q_f`$ $`=`$ $`\eta {\displaystyle d^4x\mathrm{tr}\left((m_f^{}A_5)\overline{q}_{fi}q_{fi}\right)}.`$ (4.39)
As the total electric charge vanishes in the Higgs phase, we put the constraint $`_fQ_f=0`$. Here we consider the fundamental string connecting D3 branes with net $`U(1)=\mathrm{tr}(U(N))`$ charge vanishes in the Higgs phase. The energy carried by the flavor charge becomes
$$E_Q=\eta \underset{f}{}m_f^{}Q_f=\underset{f}{}d^3xm_f^{}\mathrm{tr}\left((m_f^{}A_5)\overline{q}_{fi}q_{fi}\right)$$
(4.40)
For most of BPS objects considered here, they have moduli space parameter corresponding to a global phase rotation. The excitation along this direction would lead to the dyonic solutions. The Gauss law would give the equation for $`A_5`$ which is exactly the zero mode equation satisfied by the phase moduli coordinate in the background gauge of solutions without $`A_5`$, $`m_f^{}`$ included. The parameters $`m_f^{}`$ serve as coefficients of the excited phase moduli direction vector of the dyonic solution.
Consider a vortex-monopole composite with $`N=N_f=2`$ in a color-flavor locking phase with $`m_1<m_2`$. One can impose additional BPS condition on electric charge section without breaking any additional supersymmetry. One has to solve the above Gauss law which can be solved in priciple in this monopole-vortex background. The result describes a composite of D1-fundamental strings connecting D3 branes, which means that the monopole carries electric charge. However, the $`A_5`$ would approach exponentially vacuum expectation value away from the monopole region, implying that the electric charge is shielded by the Higgs field. For two flavor case, one can choose $`A_5A_4`$ up to constant shift as $`(m_f,m_f^{})`$ lies along a line. Note that $`E_Q(\mathrm{\Delta }m^{})^2`$ and $`Q_2Q_1\mathrm{\Delta }m^{}`$, and so the relative flavor charge fixes $`m_2^{}m_1^{}`$ as in the dyons in Coulomb phase.
When $`N_f3`$, $`D7`$ branes does not need to lie on a line as three points given by the mass parameters do not lie along a line in general. In this case one could have a web of D1, F1 and $`(p,q)`$ strings. For example consider $`N=N_f=3`$ in the color-flavor locking phase. If the D7 branes are separated from each other and lie on almost straight line, one can imagine a D1 string interpolating two D3 branes at the end. When we introduce the fundamental strings connecting, say first and second D3 branes, the resulting configuration would be a vortex string where there are two fundamental monopoles attacted to each other, but the Coulomb repulsion due to the electric charge in short distance keeps them away from each other. This is quite similar to the corresponding configuration in the Coulomb phase. The key difference would be that in the Higgs phase there may be no upper bound on F1 string numbers as the electric repulsion would be shielded in large separation.
It is straightforward to extend this to situations of multiple domain walls. Consider $`N=2`$, $`N_f=3`$ with two domain walls interpolating $`m_1,m_2`$ by first D3 and $`m_2,m_3`$ by second D3 ($`m_1<m_2<m_3`$). If we turn on $`m_2^{}`$ slightly, these two domain walls are attracted, and it is balanced by giving them electric charges proportional to $`m_2^{}`$ distributed on their world volume. This would be web-like structure of D3 branes and sheet of fundamental strings, attached to D7 branes.
Another dyonic BPS configuration is possible. Start with a 1/2 BPS domain wall of a single D3 brane, interpolating two D7 branes in position. Fundamental strings connecting two D7 branes at the wall generates the electric dipole on D3 brane. Two ends of the dipole are shielded by the Higgs field of different flavor, and so the configuration has the Higgs charge. One needs to solve the Gauss law in the domain wall background. From the domain wall world sheet point of view, the fundamental F1 string appears as a charge of phase or magnetic flux on effective 2+1 dimensional theory. Uniform charge configuration would corresponds to unform magnetic flux configuration on effective 2+1 dimensional theory.
In our BPS equation there is additional parameter $`\phi `$. To see its role in $`N=1,N_f=2`$ with $`\zeta ^a=v^2\delta _{a3}`$, $`A_5=m_f^{}=0`$, let us consider the domain wall solution with $`2,4`$ directions mixed. With only dependence on $`x^1`$ and $`x^3`$ and $`A_1=A_2=A_3=0`$, $`\eta =\alpha =1`$, $`\theta =0`$, BPS equations (3.24) and (3.25) for $`A_4`$ become
$`(_3\mathrm{cos}\phi _1\mathrm{sin}\phi )A_4={\displaystyle \frac{1}{2}}(v^2{\displaystyle \underset{f}{}}|q_{f1}|^2),`$
$`(_3\mathrm{cos}\phi _1\mathrm{sin}\phi )q_{f1}+q_{f1}(A_4m_f)=0`$
$`(_3\mathrm{sin}\phi +_1\mathrm{cos}\phi )A_4=0,(_3\mathrm{sin}\phi +_1\mathrm{cos}\phi )q_{f1}=0.`$ (4.41)
This corresponds to a spatial rotation in $`(x^1,x^3)`$ plane. The origin of this fact can be traced back to the correlation between $`(x^2,x^4)`$ and $`(x^1,x^3)`$ in the spinor projection conditions.
## 5 1/8 BPS Objects in Theories with Product Gauge Groups
While we found 1/8 BPS equations which seems to be general up to six dimensional Lorentz boost and $`SU(2)_R`$ symmetry, it is not clear whether 1/8 BPS configurations are allowed. After the dimensional reduction to 3+1 dimensions with two general angle parameters, one cannot make arbitrary six dimensional rotation, especially $`F_{45}=0`$ in $`U(1)`$ theory. While we are interested in the general characteristics of 1/8 BPS configurations, if any exists, it seems very hard to solve the BPS equations.
Let us start with a theory with a simple gauge group, say, $`U(N)`$. To find out what the characteristics of 1/8 BPS configurations are, let us start with BPS configuration of constant field strength with zero matter expectation value. From BPS equations for the gauge fields (3.24) for the constant field strength, we can make $`SU(2)_R`$ rotation to put the FI parameter to 3-th direction and $`SU(2)`$ spatial rotation in $`x^1,x^{\overline{2}},x^{\overline{3}}`$, which rotates both $`ฯต_i`$ and the gauge field strength $`F_{\overline{\mu }\overline{\nu }}`$ with $`\mu ,\nu =1,2,3,4`$. From this one can see that the constant field configuration is at most 1/4 BPS configuration.
Inhomogeneous BPS field configuration can be obtained by extracting magnetic fluxes from the system. To see whether 1/8 BPS configurations are possible when the field configuration is inhomogeneous in space, we ask whether 1/8 BPS perturbation arises in 1/4 BPS homogeneous background.
Let us start with a $`U(1)`$ gauge theory on $`3+1`$ dimension with single flavor. Let us start with a 1/4 BPS configuration which is homogeneous in space and time with $`A_0=A_5`$ and $`\eta =\alpha =1`$ with $`\theta =\phi =0`$. The FI term becomes $`๐^a=e^2v^2/2\delta _{a3}`$ and and we choose the constant 1/4 BPS field strenghs to be
$$F_{12}=\frac{e^2v^2}{2}a,F_{34}=\frac{e^2v^2}{2}(1a),F_{23}=\frac{e^2v^2}{2}b,F_{14}=\frac{e^2v^2}{2}b$$
(5.42)
with constants $`a,b`$. This is a generalization of many previously known homogenous solutions. The homogeneous BPS configuration in $`U(1)`$ Higgs model with single Higgs field represents the unform distribution of vortices on plane, which has the critical total magnetic flux. In $`SU(2)`$ gauge theory, one could have magnetic monopole sheet or homogenous field configuration with uniform instanton density. The energy density is then
$$=\frac{e^2v^4}{4}\left(1+b^2a(1a)\right)$$
(5.43)
In four dimensions, the contribution from the intersection of $`F_{12}`$ and $`F_{34}`$ can decrease the tension when $`0<a<1`$ and can be regarded as an anti-selfdual instanton part with the negative energy, which can be regarded as a bound energy of two uniform magnetic flux. Note that the minimum energy is positive.
In 3+1 dimensions, it represents the bound energy of a domain wall and infinite number of vortex strings penetrating domain walls. The number of flavors does not play any role. For $`b>0`$ and $`a`$ not in this interval induces self-dual instanton density which contributes positive energy. Note that there are critical total flux $`e^2v^2/2`$ in our unit. From the brane point of view, the above BPS solution induces D3 branes with homogeneous field on its world sheet, tilted with respect to D7 branes.
We want to see whether there is any 1/8 BPS deformation of this homogeneous configuration. The BPS equation implies that there should be nonzero $`q_i,i=1,2`$ for 1/8 BPS configurations, which we regard as a small perturbation. (Here we drop flavor index as there is only one flavor.) We solve the 1/8 BPS equation by the perturbation expansion with $`\beta =1`$. To the first order we first solve the matter BPS equation in the unform background,
$$D_1q_j\sigma _{ji}^1+D_2q_j\sigma _{ji}^2+D_3q_j\sigma _{ji}^3iD_4q_i=0,$$
(5.44)
We choose the gauge
$$A_1=0,A_2=\frac{e^2v^2}{2}(ax^1bx^3),A_3=0,A_4=\frac{e^2v^2}{2}((1a)x^3bx^1)$$
(5.45)
The above equation is satisfied if
$$_1q_i+\frac{e^2v^2}{2}x^1q_j(a\sigma ^3b\sigma ^1)_{ji}=0,_3q_i+\frac{e^2v^2}{2}x^3q_j(b\sigma ^1+(1a)\sigma ^3)_{ji}=0$$
(5.46)
One can convince oneself that only $`q_1`$ becomes normalizable along both $`x^1`$ and $`x^3`$ directions for $`b=0`$ and $`0<a<1`$ while $`q_2`$ is not normalizable at all. For $`1/8`$ BPS deviation, we need both normalizable $`q_1`$ and $`q_2`$ modes to start the perturbative approach and so there is no $`1/8`$ BPS deviation from the 1/4 BPS configuration. The BPS equation for the gauge fields indicate the second order effect of the $`q_1`$ perturbation reducing the total magnetic flux and instanton or monopole number. Thus one can guess that the above homogeneous configuration, while remaining 1/4 BPS, is continuously connected to the two intersecting flux sheet along $`x^1x^2`$ and $`x^3x^4`$ plane with finite magnetic monopole charge and negative bound energy. In the brane picture, the end result would be the intersection of D3 brane domain wall and D1 string.
While the above analysis does not provide clear picture about the existence of 1/8 BPS configurations in 8 supersymmetric $`U(N)`$ gauge theories, it suggests that 1/8 BPS configurations are unlikely.
Now consider a theory with $`U(1)\times U(1)`$ gauge group with fundamental matter fields in each gauge group and also many bi-fundamental matter fields of charge $`(+1,1)`$. Let assume that two FI parameters are not parallel and so, say, $`\zeta ^{(1)a}=\delta _{a3}`$ and $`\zeta ^{(2)a}=\delta _{a1}`$. (Here we put the proportional numbers and electric charges to be one for simplicity.) If there is no bi-fundamental matter fields, two theories are not interacting and so it is obvious that there can be 1/8 BPS configurations. They can be made of 1/4 BPS configurations of each gauge group but they are not aligned and so break the supersymmetry further to 1/8. Even when bi-fundamental fields exist, such 1/8 BPS configurations are possible if bi-fundametnal field has zero expectation value.
To see whether bi-fundamental matter field can develope any nontrivial expectation value, let us start with 1/8 BPS homogeneous configuration in this theory of two product gauge group,
$`F_{12}^{(1)}=a,F_{34}^{(1)}=1a`$ (5.47)
$`F_{23}^{(2)}=b,F_{14}^{(2)}=1b`$ (5.48)
The energy density of the configuration becomes
$$=\frac{1}{4}(2a(1a)b(1b))$$
(5.49)
With the gauge
$$A_2^{(1)}=ax^1,A^{(1)}=(1a)x^3,A_2^{(2)}=bx^3,A_4^{(2)}=(1b)x^1$$
(5.50)
The interesting question is whether there exists a nonzero mode for the bi-fundamental field $`q_i`$, whose BPS equation is satisfied if
$`_1q_i+x^1q_j(a\sigma ^3(1b)\sigma ^1)_{ji}=0`$ (5.51)
$`_3q_i+x^3q_j((1a)\sigma ^3b\sigma ^1)_{ji}`$ (5.52)
The normalizable solution along $`x^1,x^3`$ direction is possible if $`a=b=1/2`$, in which case two matrices are proportional to each other and so can be exponentiated easily. Once we found this normalizable zero mode, we can feed it to the BPS equation for the gauge field, which leads to the second order perturbation, which reduces the sum of the magnetic fluxes. Of course there will be also nontrivial BPS deformation of fundamental matter field for each gauge group. One can imagine the continuous deflation of the total flux would lead to some sort of intersecting U(1) magnetic vortex sheets, while remaining 1/8 BPS. From 1/4 BPS case, one can see that first U(1) vortex line along $`x^3`$ direction meets a first U(1) domain parallel to the $`12`$ plane. The second U(1) vortex line along $`x^1`$ direction meets a second U(1) domain wall parallel to the $`23`$ plane.Together they would remain 1/8 BPS. In addition, there would be nontrivial bi-fundamental matter field in this 1/8 BPS configuration, making two configurations to be connected together.
## 6 Nonexistence of BPS Vortices
Most of the analysis on solitons so far have been done when $`N_fN`$. Especially there would be no supersymmetric vacua if $`N_f<N`$ without adjoint hypermultiplet. When $`N_f<N`$, the adjoint hypermultiplet plays a crucial role for supersymmetric vacua to exist. When $`N=2,N_f=1`$, the explicit vacuum solution modulo local gauge transformations is known. At the vacuum the scalars in vector multiplet $`(A_4,A_5)=(m_1,m_1^{})`$, proportional to the identity matrix. With adjoint hypermultiplet, the vacuum equation $`๐^a=0`$ is the ADHM condition on two instantons in noncommutative $`U(1)`$ theory, and the moduli space metric becomes the Eguchi-Hanson space. It is depending on eight parameters, four of which are the position of the center of mass of two D3 branes in D7 branes, and so flat and does not affect our analysis. There are additional four parameters which indicates the relative distance and phase between two D3 branes in D7 branes. Due to the FI term, there would be Higgs condensation on D3 branes. Explicitly,
$`\overline{y}_1=w_1+{\displaystyle \frac{z_1}{2}}\left(\begin{array}{cc}1& \sqrt{\frac{2b}{a}}\\ 0& 1\end{array}\right),y_2=w_2+{\displaystyle \frac{z_2}{2}}\left(\begin{array}{cc}1& \sqrt{\frac{2b}{a}}\\ 0& 1\end{array}\right)`$ (6.57)
$`\overline{q}_1=v\left(\begin{array}{c}\sqrt{1b}\\ \sqrt{1+b}\end{array}\right),q_2=0`$ (6.60)
where $`a=(|z_1|^2+|z_2|^2)/(2v^2)`$ and $`b=\sqrt{a^2+1}a`$. The vacuum moduli space is characterized by four complex parameteres $`w_i,z_i`$. The parameter $`w_i`$ denotes the location of the center of mass points of two D3 branes on D7 background and the parameter $`z_i`$ denotes the relative position.
We know there are BPS multi-vortex solutions when $`N=N_f=1`$. The question is whether any BPS vortex solitons exist when $`N=2,N_f=1`$. Suppose we put a single $`D1`$ string on one of $`D3`$ branes when two D3 branes are in infinite separation. Clearly it is BPS. As we change vacuum moduli parameters so that two D3 branes are almost on top of each other, we may expect that there would be 1/2 BPS vortex solutions. To see whether it is true, we look at a consistent ansatz.
Rather the surprise appears when two D3 branes are on top of each other, or when the vacuum moduli is at minimum two sphere of Eguch-Hanson space. In this case the consistent ansatz becomes
$$\overline{y}_1=\left(\begin{array}{cc}0& Z\\ 0& 0\end{array}\right),y_2=0,\overline{q}_1=\left(\begin{array}{c}0,\\ Q_2\end{array}\right),q_2=0,A_1+iA_2=\mathrm{diag}(A,B)$$
(6.61)
The BPS equation get simplified to be ( $`=\frac{1}{2}(_1i_2)`$)
$`Zi(AB)Z=0,Q_2iBQ_2=0,`$ (6.62)
$`i(\overline{}A\overline{A})=v^2|Z|^2,`$ (6.63)
$`i(\overline{}B\overline{B})=v^2+|Z|^2|Q|^2.`$ (6.64)
Asymptotic value of $`|Z|^2`$ and $`|Q|^2`$ are $`v^2`$ and $`2v^2`$, respectively. The above BPS equations can be combined to
$`_i^2\mathrm{ln}|Z/Q|^2=v^2|Z|^2`$ (6.65)
$`_i^2\mathrm{ln}|Q|^2=v^2+|Z|^2|Q|^2`$ (6.66)
The BPS energy is determined by $`\frac{v^2}{2}(F_A+F_B)=2v^2|Z|^2`$. The vorticity of $`Z`$ and $`Q_2`$ are $`l_1,l_2`$, then the flux $`d^2x(F_AF_B)=2\pi l_1>0`$ and $`d^2xF_B=2\pi l_2>0`$. The $`d^2xF_A=2\pi (l_1+l_2)`$ and the energy is $`\pi v^2(l_1+l_2)`$. From examining the above equations, one can easily draw the fact that there is no solution with $`l_1=0,l_2>0`$ or $`l_1>0,l_2=0`$ or $`l_1=l_2>0`$. The only possibility is $`l_11l_21`$. As we move D3 branes apart, it suggests that there is no BPS configurations possible for vortices with vorticity 1 or 2 even D3 branes are apart. Assuming that the continuity of the BPS configurations here as we do not see any critical separation between D3 branes matter, there seems to be only one logical conclusion, that is, that two D3 branes with any parallel D1 string on them become repulsive. That means there is no BPS configuration with any vorticity and finite separation. This seems to be only consistent result. It would be interesting to verify this conjecture.
Acknowledgement
We thank Muneto Nitta for correspondence about his work, and Norisuke Sakai, Piljin Yi for helpful discussions. K.L. is supported in part by KOSEF R01-2003-000-10319-0 and KOSEF through CQUeST at Sogang University. H.U.Y is supported by grant No. R01-2003-000-10391-0 from the Basic Research Program of the Korea Science & Engineering Foundation. |
warning/0506/hep-lat0506007.html | ar5iv | text | # Towards the infrared limit in ๐บโข๐ผโข(๐) Landau gauge lattice gluodynamics
## I Introduction
Studying the relation of non-perturbative features of QCD, such as confinement and dynamical chiral symmetry breaking, to the properties of propagators, there are two popular approaches at present: lattice gauge theory and Dyson-Schwinger equations (DSE). The latter approach allows to address directly the low-momentum region for a coupled system of quark, gluon and ghost propagators which is of interest for hadron physics Alkofer and von Smekal (2001). In particular their infrared behavior could be related to the mechanism of dynamical symmetry breaking and to confined gluons Alkofer and von Smekal (2001); Fischer and Alkofer (2003).
In fact, the DSE approach has revealed that in the infrared momentum region a diverging ghost propagator $`G`$ is intimately connected with a suppressed gluon propagator $`D_{\mu \nu }`$. In Landau gauge, they can be written as
$`D_{\mu \nu }(q)`$ $`=`$ $`\left(\delta _{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right){\displaystyle \frac{Z_D(q^2)}{q^2}},`$ (1)
$`G(q)`$ $`=`$ $`{\displaystyle \frac{Z_G(q^2)}{q^2}}.`$ (2)
Here $`Z_D(q^2)`$ and $`Z_G(q^2)`$ denote the dressing functions of the corresponding propagators. They describe the deviation from the momentum dependence of the free propagators. Based on the Dyson-Schwinger approach and under mild assumptions these functions are predicted to behave in the limit $`q^20`$ as follows Alkofer and von Smekal (2001):
$$Z_D(q^2)(q^2)^{\kappa _D},Z_G(q^2)(q^2)^{\kappa _G}$$
(3)
with exponents satisfying $`\kappa _D=2\kappa _G`$. In Landau gauge $`\kappa _G0.595`$ Lerche and von Smekal (2002); Zwanziger (2002). Thus the ghost propagator diverges stronger than $`1/q^2`$ and the gluon propagator is infrared suppressed. This is in agreement with the Zwanziger-Gribov horizon condition Zwanziger (2004, 1994); Gribov (1978) as well as with the Kugo-Ojima confinement criterion Kugo and Ojima (1979). Zwanziger Zwanziger (2004) has suggested that in the continuum this behavior of the propagators in Landau gauge results from the restriction of the gauge fields to the Gribov region $`\mathrm{\Omega }`$, where the Faddeev-Popov operator is non-negative.
Using further the ghost-ghost-gluon vertex, the gluon and ghost dressing functions can be used to determine a renormalization group invariant running coupling in a momentum subtraction scheme as von Smekal et al. (1997, 1998); Bloch et al. (2004)
$$\alpha _s(q^2)=\frac{g_0^2}{4\pi }Z_G^2(q^2)Z_D(q^2)$$
(4)
which then enters the quark DSE Alkofer and von Smekal (2001); Fischer and Alkofer (2003). This definition relies on the fact that the ghost-ghost-gluon vertex renormalization function $`Z_1(q^2)`$ is constant, which is true at least to all orders in perturbation theory Taylor (1971). Indeed, a recent numerical investigation of $`Z_1`$ for the $`SU(2)`$ case shows that also nonperturbatively $`Z_1`$ is finite and constant Cucchieri et al. (2004). Applying the behavior given in Eq. (3) the running coupling approaches a finite value $`\alpha _s(0)8.915/N`$ for $`SU(N)`$ at zero momentum in the DSE approach Lerche and von Smekal (2002).
Nevertheless, numerical investigations of those features in lattice simulations are still necessary to check to what extent the truncation of the coupled set of DSEs influences the final result. There are several studies in Landau gauge which confirm the anticipated behavior at least for the $`SU(2)`$ case Bloch et al. (2004); Langfeld et al. (2002); Gattnar et al. (2004); Bloch et al. (2003). Also lattice studies (Oliveira and Silva (2005); Sternbeck et al. (2005a, b), Furui and Nakajima (2004a, b) and references therein) for the $`SU(3)`$ case indicate the correctness of the proposed infrared behavior. However, as recent DSE investigations show Fischer et al. (2005, 2002); Fischer and Alkofer (2002) the infrared behavior of the gluon and ghost dressing functions and of the running coupling is changed on a torus. In particular, the running coupling decreases at low momenta.
This paper presents a lattice study of the gluon and ghost dressing functions and of the running coupling at low momenta in $`SU(3)`$ Landau gauge. We also focus, more carefully than usually, on the problem of the Gribov ambiguity in lattice simulations. In continuum, a gauge orbit has more than one intersection (Gribov copies Gribov (1978)) with the transversality plane (where $`_\mu A_\mu =0`$ holds for the gauge potential $`A_\mu `$) within the Gribov region $`\mathrm{\Omega }`$. Expectation values taken over this region are argued to be equal to those over the fundamental modular region $`\mathrm{\Lambda }`$ which includes only the absolute maximum of the gauge functional Zwanziger (2004).
On a finite lattice, however, this equality cannot be expected Zwanziger (2004). In the literature it is widely taken for granted that the gluon propagator does not depend on the choice of Gribov copy, while an impact on the $`SU(2)`$ ghost propagator has been observed Cucchieri (1997); Bakeev et al. (2004); Nakajima and Furui (2004). However, in a more recent investigation Silva and Oliveira (2004) an influence of the Gribov copies ambiguity on the $`SU(3)`$ gluon propagator has been demonstrated, too. Here we assess the importance of the Gribov ambiguity on a finite lattice for the $`SU(3)`$ ghost propagators.
This paper is structured as follows: In Sec. II we shall define all quantities which are investigated in this study. Then, after specifying the lattice setup used, the dependence of the gluon and ghost propagator on the choice of Gribov copies and lattice discretization as well as finite-volume effects are discussed in Sec. III. We also discuss the problem of exceptional gauge copies in Sec. IV. In Sec. V the behavior of the running coupling is presented. In Appendix A we show how the inversion of the F-P operator can be accelerated by pre-conditioning with a Laplacian operator.
## II Definitions
To study the ghost and gluon propagators using lattice simulations one has to fix the gauge for each thermalized $`SU(3)`$ gauge field configuration $`U=\{U_{x,\mu }\}`$. We adopted the Landau gauge condition which can be implemented by searching for a gauge transformation
$$U_{x,\mu }{}_{}{}^{g}U_{x,\mu }^{}=g_xU_{x,\mu }g_{x+\widehat{\mu }}^{}$$
which maximizes the Landau gauge functional
$$F_U[g]=\frac{1}{4V}\underset{x}{}\underset{\mu =1}{\overset{4}{}}๐ข\text{Tr}{}_{}{}^{g}U_{x,\mu }^{}$$
(5)
while keeping the Monte Carlo configuration $`U`$ fixed. Here $`g_x`$ are elements of $`SU(3)`$.
The functional $`F_U[g]`$ has many different local maxima which can be reached by inequivalent gauge transformations $`g`$, the number of which increases with the lattice size. As the inverse coupling constant $`\beta `$ is decreased, increasingly more of those maxima become accessible by an iterative gauge fixing process starting from a given (random) gauge transformation $`g_x`$. The different gauge copies corresponding to those maxima are called *Gribov copies*, due to their resemblance to the Gribov ambiguity in the continuum Gribov (1978). All Gribov copies $`\{{}_{}{}^{g}U\}`$ belong to the same gauge orbit created by the Monte Carlo configuration $`U`$ and satisfy the differential Landau gauge condition (lattice transversality condition) $`(_\mu {}_{}{}^{g}A_{\mu }^{})(x)=0`$ where
$$(_\mu {}_{}{}^{g}A_{\mu }^{})(x){}_{}{}^{g}A_{\mu }^{}(x+\widehat{\mu }/2){}_{}{}^{g}A_{\mu }^{}(x\widehat{\mu }/2).$$
(6)
Here $`{}_{}{}^{g}A_{\mu }^{}(x+\widehat{\mu }/2)`$ is the non-Abelian (hermitian) lattice gauge potential which may be defined at the midpoint of a link
$`{}_{}{}^{g}A_{\mu }^{}(x+\widehat{\mu }/2)`$ $`=`$ $`{\displaystyle \frac{1}{2iag_0}}\left({}_{}{}^{g}U_{x,\mu }^{}^gU_{x,\mu }^{}\right)`$ (7)
$`{\displaystyle \frac{\mathrm{๐}}{6iag_0}}\text{Tr}\left({}_{}{}^{g}U_{x,\mu }^{}^gU_{x,\mu }^{}\right).`$
In this way it is accurate to $`๐ช(a^2)`$. The bare gauge coupling $`g_0`$ is related to the inverse lattice coupling via $`\beta =6/g_0^2`$ in the case of $`SU(3)`$. In the following, we will drop the label $`g`$ for convenience, i.e. we consider $`U`$ to be already put into the Landau gauge such that $`g=\mathrm{๐}`$ maximizes the functional in Eq. (5) relative to the neighborhood of the identity.
The gluon propagator $`D_{\mu \nu }^{ab}(q^2)`$ is the Fourier transform of the gluon two-point function, i.e. the expectation value
$$D_{\mu \nu }^{ab}(q)=\stackrel{~}{A}_\mu ^a(k)\stackrel{~}{A}_\nu ^b(k)=\delta ^{ab}D_{\mu \nu }(q)$$
(8)
which is required to be color-diagonal. Here $`\stackrel{~}{A}_\mu ^a(k)`$ is the Fourier transform of $`A_\mu ^a(x+\widehat{\mu }/2)`$ and $`q`$ denotes the momentum
$$q_\mu (k_\mu )=\frac{2}{a}\mathrm{sin}\left(\frac{\pi k_\mu }{L_\mu }\right)$$
(9)
which corresponds to a integer-valued lattice momentum $`k`$. Since $`k_\mu (L_\mu /2,L_\mu /2]`$ the lattice equivalent of $`q^2(k)`$ can be realized by different $`k`$. According to Ref. Leinweber et al. (1999), however, a subset of lattice momenta $`k`$ has been considered only for the final analysis of the gluon propagator, although the FFT algorithm provides us with all lattice momenta. Details are given below.
Assuming reality and rotational invariance we envisage for the (continuum) gluon propagator the general tensor structure:
$$D_{\mu \nu }(q)=\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)D(q^2)+\frac{q_\mu q_\nu }{q^2}\frac{F(q^2)}{q^2}$$
(10)
with $`D(q^2)`$ and $`F(q^2)`$ being scalar functions. On the lattice these functions are extracted by projection and are expected to scatter, rather than being smooth functions of $`q^2`$. Using the Landau gauge condition the longitudinal form factor $`F(p^2)`$ vanishes. Recalling the mentioned Gribov ambiguity of the chosen gauge copy there is no a priori reason to assume the estimator of $`D(q^2)`$ is not influenced by the choice.
The ghost propagator is derived from the Faddeev-Popov (F-P) operator, the Hessian of the gauge functional given in Eq. (5). We expect that the properties of the F-P operator differ for the different maxima of the functional (Gribov copies). This should have consequences for the ghost propagator as is shown below.
After some algebra the F-P operator can be written in terms of the (gauge-fixed) link variables $`U_{x,\mu }`$ as
$`M_{xy}^{ab}`$ $`=`$ $`{\displaystyle \underset{\mu }{}}A_{x,\mu }^{ab}\delta _{x,y}B_{x,\mu }^{ab}\delta _{x+\widehat{\mu },y}C_{x,\mu }^{ab}\delta _{x\widehat{\mu },y}`$ (11)
with
$`A_{x,\mu }^{ab}`$ $`=`$ $`๐ข\text{Tr}\left[\{T^a,T^b\}(U_{x,\mu }+U_{x\widehat{\mu },\mu })\right],`$
$`B_{x,\mu }^{ab}`$ $`=`$ $`2๐ข\text{Tr}\left[T^bT^aU_{x,\mu }\right],`$
$`C_{x,\mu }^{ab}`$ $`=`$ $`2๐ข\text{Tr}\left[T^aT^bU_{x\widehat{\mu },\mu }\right].`$
Here $`T^a`$ and $`T^b`$ are the (hermitian) generators of the $`๐ฐ๐ฒ(3)`$ Lie algebra satisfying $`\text{Tr}[T^aT^b]=\delta ^{ab}/2`$.
The ghost propagator is calculated as the following ensemble average
$$G^{ab}(q)=\frac{1}{V}\underset{x,y}{}\mathrm{e}^{2\pi ik(xy)}[M^1]_{xy}^{ab}_U.$$
(12)
It is diagonal in color space: $`G^{ab}(q)=\delta ^{ab}G(q)`$. Following Ref. Suman and Schilling (1996); Cucchieri (1997) we have used the conjugate gradient (CG) algorithm to invert $`M`$ on a plane wave $`\stackrel{}{\psi }_c`$ with color and position components $`\psi _c^a(x)=\delta ^{ac}\mathrm{exp}(2\pi ikx)`$. In fact, we applied the pre-conditioned CG algorithm (PCG) to solve $`M_{xy}^{ab}\varphi ^b(y)=\psi _c^a(x)`$. As pre-conditioning matrix we used the inverse Laplacian operator $`\mathrm{\Delta }^1`$ with diagonal color substructure. This significantly reduces the amount of computing time as it is discussed in more detail in Appendix A.
After solving $`M\stackrel{}{\varphi }=\stackrel{}{\psi }_c`$ the resulting vector $`\stackrel{}{\varphi }`$ is projected back on $`\stackrel{}{\psi }_c`$ such that the average $`G^{cc}(q)`$ over the color index $`c`$ can be taken explicitly. Since the F-P operator $`M`$ is singular if acting on constant modes, only $`k(0,0,0,0)`$ is permitted. Due to high computational requirements to invert the F-P operator for each $`k`$, separately, the estimator on a single, gauge-fixed configuration is evaluated only for a preselected set of momenta $`k`$. In Table 1 a detailed list is given.
## III Results for the ghost and gluon propagators
### III.1 Lattice samples
For the purpose of this study we have analyzed pure $`SU(3)`$ gauge configurations which have been thermalized with the standard Wilson action at three values of the inverse coupling constant $`\beta =5.8`$, $`6.0`$ and $`6.2`$. For thermalization an update cycle of one heatbath and four micro-canonical over-relaxation steps was used. As lattice sizes we used $`16^4`$, $`24^4`$ and $`32^4`$. For tests, exposing the inherent problems of the Gribov problem under the aspect of volume dependence, also smaller lattices ($`8^4`$ and $`12^4`$) have been considered at lower cost.
To each thermalized configuration $`U`$ a random set of $`N_{\mathrm{cp}}`$ local gauge transformation $`\{g_x\}`$ were assigned. Each of those served as starting point for a gauge fixing procedure for which we used standard *over-relaxation* with over-relaxation parameter tuned to $`\omega =1.63`$. Keeping all $`U_{x,\mu }`$ fixed this iterative procedure creates a sequence of local gauge transformations $`g_x`$ at sites $`x`$ with increasing values of the gauge functional (Eq. (5)). Thus, the final Landau gauge is iteratively approximated until the stopping criterion in terms of the transversality (see Eq. (6))
$$\underset{x}{\mathrm{max}}\left[_\mu {}_{}{}^{g}A_{\mu }^{}(x)\right]^2<10^{14}$$
(13)
is fulfilled. Consequently, each random start results in its own local maximum of the gauge functional. Certain extrema of the functional are found multiple times. In fact, this happened frequently on the small lattices, $`8^4`$ and $`12^4`$, but rather seldom on larger lattices. Note that we used the maximum in relation (13) which is very conservative. However, the precision of transversality dictates how symmetric the F-P operator $`M`$ can be considered. This is crucial for its inversion and thus dictates the final precision of the ghost propagator.
To study the dependence on Gribov copies of the propagators, in the course of $`N_{\mathrm{cp}}`$ repetitions for each $`U`$, the gauge copy with the largest functional value is stored under the name *best copy* (bc). The first gauge copy is also stored, labeled as *first copy* (fc). However, it is as good as any other arbitrarily selected gauge copy.
The more gauge copies one gets to inspect, the bigger the likeliness that the copy labeled as bc actually represents the absolute maximum of the functional in Eq. (5). With increasing number $`N_{\mathrm{cp}}`$ of copies the expectation value of gauge variant quantities, evaluated on bc representatives, is converging more or less rapidly as we will discuss next.
### III.2 How severe is the lattice Gribov problem for the propagators?
First we present results of a combined study of the gluon and ghost propagators on the same sets of fc and bc representatives of our thermalized gauge field configurations. This allows us to assess the importance of the Gribov copy problem.
Numerically, it turns out that the dependence of the ghost propagator on the choice of the best copy is most severe for the smallest momentum. In addition, this depends on the lattice size and $`\beta `$. Therefore we studied first the dependence of the ghost and gluon propagators at lowest momentum on the (same) best copies as function of the number of gauge copies $`N_{\mathrm{cp}}`$ under inspection. This was done at $`\beta =6.2`$ where we used $`12^4`$, $`16^4`$ and $`24^4`$ lattices. The number of thermalized configurations used for these three lattice sizes are given in Table 1. To check the dependence on $`\beta `$ also a simulation at $`\beta =5.8`$ on a $`24^4`$ lattice was performed.
The results of this investigation are show in Fig. 1. While there the ghost propagator is shown as an average over the two realizations $`k=(1,0,0,0)`$ and $`k=(0,1,0,0)`$ of the smallest lattice momentum $`q(k)`$, the gluon propagator has been averaged over all four non-equivalent realizations. Note that $`D(k)=D(k)`$. It is clearly visible that the expectation value of the gluon propagator does not change within errors as $`N_{\mathrm{cp}}`$ increases, independent of the lattice size and $`\beta `$. Contrarily, the ghost propagator at $`\beta =5.8`$ on a $`24^4`$ lattice saturates (on average) if calculated on the best among $`N_{\mathrm{cp}}=15`$ gauge copies. At $`\beta =6.2`$ the number of gauge fixings attempts reduces to $`5N_{\mathrm{cp}}10`$ on a $`16^4`$ and $`24^4`$ lattice. On the $`12^4`$ lattice a small impact of Gribov copies is visible, namely $`1<N_{\mathrm{cp}}5`$. The lower panels of Fig. 1 show the relative difference $`\delta F=1F^{\text{cbc}}/F^{\text{bc}}`$ of the corresponding (current best) functional values $`F^{\text{cbc}}`$ to the value $`F^{\text{bc}}`$ of the overall best copy after $`N_{\mathrm{cp}}=20`$, respectively $`N_{\mathrm{cp}}=30`$, attempts. This may serve as an indicator how large $`N_{\mathrm{cp}}`$ has to be on average for the chosen algorithm to have found a maximum of $`F`$ close to the global one.
In order to study further the low-momentum dependence of the gluon and ghost dressing functions, $`Z_D`$ and $`Z_G`$, as given by Eq. (1) and Eq. (2) we have performed similar simulations using lattice sizes $`12^4`$, $`16^4`$, $`24^4`$ and $`32^4`$ at $`\beta =5.8`$, 6.0 and 6.2. Following Ref. Necco and Sommer (2002) these values of $`\beta `$ correspond to $`a^1`$=1.446 GeV, 2.118 GeV and 2.914 GeV using the Sommer scale $`r_0=0.5`$ fm. These values of the lattice spacing $`a`$ associated to $`\beta `$ turn out to be more appropriate as those formerly used by us and others (see Sternbeck et al. (2005a, b); Silva and Oliveira (2004); Leinweber et al. (1999)).
We have fixed a conservative number of $`N_{\mathrm{cp}}=30`$ gauge copies per thermalized configuration on a $`16^4`$ lattice and $`N_{\mathrm{cp}}=40`$ on a $`24^4`$ lattice. Both the gluon and the ghost propagator, respectively their dressing functions, have been measured on the same set of fc and bc copies. Due to the large amount of computing time necessary for the $`32^4`$ lattice we could afford to measure the ghost propagator for the first and best among only $`N_{\mathrm{cp}}=10`$ copies, which is certainly not enough.
The data for the gluon propagator $`D(k)`$ have been determined for all momenta at once. However, we used only a subset of momenta for the final analysis. In fact, inspired by Ref. Leinweber et al. (1999), only data $`D(k)`$ with $`k`$ lying in a cylinder with radius of one momentum unit along one of the diagonals $`\widehat{n}=1/2(\pm 1,\pm 1,\pm 1,\pm 1)`$ have been selected. Since we are using a symmetric lattice structure only data with $`k`$ satisfying $`_\mu k_\mu ^2(_\mu k_\mu \widehat{n}_\mu )^21`$ are surviving this cylindrical cut. In agreement with Leinweber et al. (1999) this recipe has drastically reduced lattice artifacts, in particular for large momenta. Additionally, we try to keep finite volume artifacts at lower momenta under control by removing all data $`D(k)`$ with one or more vanishing momentum components $`k_\mu `$ Leinweber et al. (1999). However, this we have done only for data on a $`12^4`$ and $`16^4`$ lattice. In Sec. III.3 we shall discuss in more detail finite volume effects at various momenta.
In view of this we have chosen appropriate sets of momenta for the ghost propagator, as listed in Table 1 in detail.
The final results of the dressing functions $`Z_D`$ and $`Z_G`$ measured on bc copies are shown in the upper panels of Fig. 2. All momenta $`q(k)`$ have been mapped to physical momenta using the lattice spacings $`a`$ given above. As expected the ghost dressing function diverges with decreasing momenta, while the gluon dressing function decreases after passing a turnover at about $`q^2=0.8`$ GeV<sup>2</sup>. However, for the purpose of the expected infrared behavior given in Eq. (3) the data for momenta $`q^2<0.25`$ GeV<sup>2</sup> are not sufficiently abundant to extract a critical exponent $`\kappa _G>0.5`$ as expected from the Dyson-Schwinger approach. In particular, the fit parameters are not stable under a change of the upper momentum cutoff. The best fits give $`\kappa _G0.25`$.
In the lower panels of Fig. 2 we present the ratio of the dressing functions, $`Z^{\text{fc}}/Z^{\text{bc}}`$, calculated using jackknife from first and best gauge copies as a function of the momentum. There the data from simulations on a $`32^4`$ lattice have been excluded, since only $`N_{\mathrm{cp}}=10`$ gauge copies have been inspected which would result in an underestimate of the ratio $`Z^{\text{fc}}/Z^{\text{bc}}`$. As is clear from these panels we do not observe a systematic dependence on the choice of Gribov copies for the gluon propagator. In contrast the ghost propagator is systematically overestimated for fc (arbitrary) gauge copies. This effect holds up to momenta of about 2 GeV<sup>2</sup>.
Comparing also the ratios for the ghost propagator at $`q<1`$ GeV, the rise at $`\beta =6.0`$ is obviously larger than that at $`\beta =5.8`$. In both cases the data are from simulations on a $`24^4`$ lattice. Thus, it seems that by increasing the physical volume (lower $`\beta `$) the effect of the Gribov ambiguity gets smaller if the same physical momentum is considered. This cannot be due to a too small number $`N_{\mathrm{cp}}`$ of inspected gauge copies since, judged from Fig. 1, $`N_{\mathrm{cp}}=40`$ seems to be on the safe side.
We conclude: the ghost propagator is systematically dependent on the choice of Gribov copies, while the impact on the gluon propagator is not resolvable within our statistics. However, there are indications that the dependence on Gribov copies decreases with increasing physical volume. This is also in agreement with the data listed in the two lattice studies Bakeev et al. (2004); Cucchieri (1997) of the $`SU(2)`$ ghost propagator $`G`$, while it is not explicitly stated there. In fact, in Ref. Bakeev et al. (2004) the ratio $`G^{\text{fc}}/G^{\text{bc}}`$ at $`\beta =2.2`$ on a $`8^4`$ lattice is larger than that on a $`16^4`$ lattice at the same physical momentum.
### III.3 Systematic effects of lattice spacings and volumes
We remind that in Fig. 2 we have dropped all data related to a $`16^4`$ lattice with one or more vanishing momentum components $`k_\mu `$. According to Leinweber et al. (1999) this keeps finite volume effects for the gluon propagator under control. It is quite natural to analyze here the different systematic effects on the gluon and ghost propagators of changing either the lattice spacing $`a`$ or the physical volume $`V`$. However, due to the preselected set of momenta for the ghost propagator and the values chosen for $`\beta `$, we can study this here only in a limited way and in a region of intermediate momenta. For the gluon propagator this has been done in more detail by other authors (see e.g. Bonnet et al. (2001)).
Keeping first the lattice spacing fixed we have found that both the ghost and gluon dressing functions calculated at the same physical momentum $`q^2`$ decrease as the lattice size $`L^4`$ is increased. This is illustrated for various momenta in Fig. 3. There both dressings functions versus the physical momentum are shown for different lattice sizes at either $`\beta =5.8`$ or $`\beta =6.2`$. Note, in this figure we have not dropped data with vanishing momentum components $`k_\mu `$ to emphasize the influence of a finite volume on those (low) momenta. We also show data from simulations on a $`8^4`$ and $`12^4`$ lattice. One clearly sees that the lower the momenta the larger the effect due to a finite volume. In comparison with $`\beta =5.8`$ this is even more drastic at $`\beta =6.2`$. At this $`\beta `$ the lattice spacing is about $`a=0.068`$ fm. Thus the largest volume considered at $`\beta =6.2`$ is about (1.6 fm)<sup>4</sup>, which is even smaller than the physical volume of a $`16^4`$ lattice at $`\beta =5.8`$.
Altogether we can state that for both dressing functions finite volume effects are clearly visible at volumes smaller than (2.2 fm)<sup>4</sup>, which corresponds to a $`16^4`$ lattice at $`\beta =5.8`$. The effect grows with decreasing momenta or decreasing lattice size (see the right panels in Fig. 3). At larger volumes, however, the data for $`q>1`$ GeV coincide within errors for the different lattice sizes (left panels). Even for $`q<1`$ GeV we cannot resolve finite volume effects for both dressing functions based on the data related to a $`24^4`$ and $`32^4`$ lattice.
Based on our chosen values for $`\beta `$ and the lattice sizes we can select equal physical volumes only approximately. Hence also the physical momenta are only approximately the same if the ghost and gluon dressing functions are compared at different $`\beta `$, e.g. at different lattice spacings. Therefore, it is difficult to analyze the systematic effect of changing $`a`$ if for both dressing functions small variations in $`q^2`$ are hidden. Consequently, in Fig. 4 we show the data for the ghost and gluon dressing functions approximately at the same physical volume $`V(2.2fm)^4`$ for different $`a`$, albeit as functions of $`q^2`$. This allows us to disentangle by eyes a change of the data due to varying $`a`$ from the natural dependence of the propagators on $`q^2`$. Inspecting Fig. 4 one concludes that the gluon dressing function at the same physical momentum and volume increases with decreasing lattice spacing. A similar effect (beyond error bars) we cannot report for the ghost dressing function.
## IV The problem of exceptional configurations
We turn now to a peculiarity of the ghost propagator at larger $`\beta `$ which has also been observed by some of us in an earlier $`SU(2)`$ study Bakeev et al. (2004). While inspecting our data we found, though rarely, that there are outliers in the Monte Carlo time histories of the ghost propagator at lowest momentum. Those outliers are not equally distributed around the average value, but are rather significantly larger than this.
In Fig. 5 we present time histories of the ghost propagator $`G(k)`$ measured of fc and bc gauge copies for two smallest momentum realizations $`k=(1,0,0,0)`$ and $`k=(0,1,0,0)`$, separately. From left to right the panels are ordered in increasing order of the lattice sizes $`12^4`$, $`16^4`$ and $`24^4`$ at $`\beta =6.2`$ and $`24^4`$ at $`\beta =5.8`$.
As can be seen from this figure in the majority extreme spikes are reduced (or even not seen) when the ghost propagator could be measured on a better gauge copy (bc) for a particular configuration. Furthermore, it is obvious that the *exceptionality* of a given gauge copy is exhibited not simultaneously for all different realizations of the lowest momentum. Consequently, to reduce the impact of such large values on the average ghost propagator one should better average over all momentum realizations giving rise to the same momentum $`q`$. This has been done for the results shown in Fig. 1 and 2 at least for the lowest momentum at $`\beta =6.2`$. However, compared to the gluon propagator it takes much more computing time to determine the ghost propagator for all its different realizations of momentum $`q(k)`$.
In addition, we have tried to find a correlation of such outliers in the history of the ghost propagator with other quantities measured in our simulations. For example we have checked whether there is a correlation between the values of the ghost propagator $`G(k)`$ as they appear in Fig. 5 with low-lying eigenvalues and eigenvectors of the F-P operator. They are apparent in the contribution of the lowest 10 F-P eigenmodes to the ghost propagator at this particular $`k`$. This we shall present in a subsequent publication Sternbeck et al. (2005c) where we shall discuss the spectral properties of the F-P operator and its relation to the ghost propagator.
## V The running coupling
We shall now focus on the running coupling $`\alpha _s(q^2)`$ as defined in Eq. (4) where $`g_0^2/(4\pi )=3/(2\pi \beta )`$ for $`SU(3)`$. Given the raw data for the gluon and ghost dressing functions on bc gauge copies the average $`Z_G^2(q^2)Z_D(q^2)`$ and its error have been estimated using the *bootstrap* method with drawing 500 random samples. Since the ghost-ghost-gluon-vertex renormalization function $`Z_1`$ has been set to one, there is an overall normalization factor which has been fixed by fitting the data for $`q^2>q_c^2`$ to the well-known perturbative results of the running coupling $`\alpha _{\mathrm{๐ธ}\mathrm{๐๐๐๐}}`$ at 2-loop order (see also Bloch et al. (2004)). Defining $`xq^2/\mathrm{\Lambda }_{\mathrm{๐ธ}\mathrm{๐๐๐๐}}^2`$, the 2-loop running coupling is given by
$$\alpha _{\mathrm{๐ธ}\mathrm{๐๐๐๐}}(x)=\frac{4\pi }{\beta _0\mathrm{ln}x}\left\{1\frac{2\beta _1}{\beta _0^2}\frac{\mathrm{ln}(\mathrm{ln}x)}{\mathrm{ln}x}\right\}.$$
(14)
The $`\beta `$-function coefficients are $`\beta _0=11`$ and $`\beta _1=51`$ for the $`SU(3)`$ gauge group and are independent of the renormalization prescription. The value of $`\mathrm{\Lambda }_{\mathrm{๐ธ}\mathrm{๐๐๐๐}}`$ has been fixed by the same fit. The lower bound $`q_c^2`$ has been chosen such that an optimal value for $`\chi ^2/\mathrm{๐๐๐}`$ has been achieved.
The results are shown in Fig. 6. There also the 1-loop contribution is shown where we used the same lower bound $`q_c^2`$. The best fit of the 2-loop expression to the data gives $`\mathrm{\Lambda }_{\mathrm{๐ธ}\mathrm{๐๐๐๐}}`$=0.88(7) GeV ($`\chi ^2/\mathrm{๐๐๐}=0.96`$), while $`\mathrm{\Lambda }_{\mathrm{๐ท}\mathrm{๐๐๐๐}}`$=0.64(7) is obtained ($`\chi ^2/\mathrm{๐๐๐}=1.05`$) using just the 1-loop part. For $`q_c^2`$ we used $`q_c^2=30`$ GeV<sup>2</sup>. The value for $`\mathrm{\Lambda }_{\mathrm{๐ธ}\mathrm{๐๐๐๐}}`$ is similar within errors to the $`SU(2)`$ result given in Ref. Bloch et al. (2004).
Approaching the infrared limit in Fig. 6 one clearly sees a running coupling $`\alpha _s(q^2)`$ increasing for $`q^2>0.4`$ GeV<sup>2</sup>. However, after passing a maximum at $`q^20.4`$ GeV<sup>2</sup> $`\alpha _s(q^2)`$ decreases again. Such turnover is in agreement with DSE results obtained on a torus Fischer et al. (2002); Fischer and Alkofer (2002); Fischer et al. (2005). Therefore, one can argue that this behavior is a finite lattice effect although we cannot resolve a difference between the different lattice sizes used. A similar infrared behavior for the running coupling has also been observed in different lattice studies Furui and Nakajima (2004b, a). But opposed to Furui and Nakajima (2004a) the existence of a turnover is independent on the choice of Gribov copy, since qualitatively, we have found the same behavior for $`\alpha _s(q^2)`$ calculated on fc gauge copies.
For completeness we mention that running couplings decreasing in the infrared have also been found in lattice studies of the 3-gluon vertex Boucaud et al. (1998, 2003) and the quark-gluon vertex Skullerud and Kizilersu (2002).
Apart from the finite volume argument given above to explain such a behavior, which prevents us from seeing the limit $`\alpha _s(0)0`$ mentioned in the Introduction, one could also put into question whether one can really set $`Z_1(q)=1`$ at lower momenta. A recent investigation dedicated to the ghost-ghost-gluon-vertex renormalization function $`Z_1(q)`$ for the case of $`SU(2)`$ Cucchieri et al. (2004) supports that $`Z_1(q)1`$ at least for $`q>1`$ GeV.
## VI Conclusions
We have reported on a numerical study of the gluon and ghost propagators in Landau gauge using several lattice sizes at $`\beta =5.8`$, $`6.0`$ and $`6.2`$. Studying the dependence on the choice of Gribov copies, it turns out that for the gluon propagator the effect of Gribov copies stays inside numerical uncertainty, while the impact on the ghost propagator increases as the momentum or $`\beta `$ is decreased. However, there are indications that the influence of Gribov copies decreases as the physical volume is increased. This is at least expected in the light of Ref. Zwanziger (2004). There it is argued that in the continuum expectation values of correlation functions $`A(x_1)\mathrm{}A(x_n)`$ over the fundamental modular region $`\mathrm{\Lambda }`$ are equal to those over the Gribov region $`\mathrm{\Omega }`$, since functional integrals are dominated by the common boundary of $`\mathrm{\Lambda }`$ and $`\mathrm{\Omega }`$. Thus Gribov copies inside $`\mathrm{\Omega }`$ should not affect expectation values in the continuum.
While the effect of the Gribov ambiguity on the ghost propagator becomes smaller with increasing $`\beta `$, exceptionally large values appear in the history of the ghost propagator in agreement with what has been observed first in reference Bakeev et al. (2004). These outliers we have not seen simultaneously for all lattice momenta $`k`$ realizing the same lowest momentum $`q(k)`$. However, they are apparent in the contribution of the lowest 10 F-P eigenmodes to the ghost propagator at this particular $`k`$ Sternbeck et al. (2005c). Therefore it is good practice to measure the ghost propagator for more than one $`k`$ with the same momentum $`q(k)`$, in order to reduce the systematic errors coming from such exceptional values.
We have studied the effects of finite volume on the one hand and of finite lattice spacing on the other. The first ones are found to be essential for volumes smaller than $`(2.2\mathrm{fm})^4`$ at the same $`\beta `$ whereas the discretization effects at the same volume are modest. Our available data did not allow us to extend this analysis to physical momenta below 1 GeV, where the Gribov ambiguity shows up and where a similar separation of finite volume and discretization effects would be desirable. However, we could observe from Fig. 2 that enlarging the volume by decreasing $`\beta `$ leads to a reduction of the systematic Gribov effect in the ghost propagator.
The dressing functions, $`Z_D`$ and $`Z_G`$, of the gluon and ghost propagators have allowed us to estimate the behavior of a running coupling $`\alpha _s(q^2)`$ in a momentum subtraction scheme. Going from larger momenta $`q^2`$ to lower ones $`\alpha _s(q^2)`$ is steadily increasing until $`q^20.4`$ GeV<sup>2</sup>. For $`q^2<0.3`$ GeV<sup>2</sup>, however, $`\alpha _s(q^2)`$ is decreasing. A decreasing running coupling at low momenta is in qualitative agreement with recent DSE results obtained on a torus Fischer et al. (2002); Fischer and Alkofer (2002); Fischer et al. (2005). Therefore one might conclude that the decrease is due to *finite* lattice volumes we used. It makes it questionable whether lattice simulations in near future can confirm the predicted infrared behavior of the gluon and ghost dressing functions with related exponents $`\kappa _D=2\kappa _G`$.
## Appendix A Speeding up the inversion of the F-P operator
For the solution of the linear system $`M\stackrel{}{\varphi }=\stackrel{}{\psi }_c`$ with symmetric matrix $`M`$, the conjugate gradient (CG) algorithm is the method of choice. Its convergence rate depends on the condition number, the ratio of largest to lowest eigenvalue of $`M`$. When all $`U_{x,\mu }=\mathrm{๐}`$ obviously the F-P operator is minus the Laplacian $`\mathrm{\Delta }`$ with a diagonal color substructure. Thus instead of solving $`M\stackrel{}{\varphi }=\stackrel{}{\psi }_c`$ one rather solves the transformed system
$$[M\mathrm{\Delta }^1](\mathrm{\Delta }\stackrel{}{\varphi })=\stackrel{}{\psi }_c$$
In this way the condition number is reduced, however, the price to pay is one extra matrix multiplication by $`\mathrm{\Delta }^1`$ per iteration cycle. In terms of CPU time this should be more than compensated by the reduction of iterations.
The pre-conditioned CG algorithm (PCG) can be described as follows:
initialize:
$`\stackrel{}{r}^{(0)}`$ $`=`$ $`\stackrel{}{\psi }M\stackrel{}{\varphi }^{(0)},\stackrel{}{p}^{(0)}=\mathrm{\Delta }^1\stackrel{}{r}^{(0)},`$
$`\gamma ^{(0)}`$ $`=`$ $`(\stackrel{}{p}^{(0)},\stackrel{}{r}^{(0)})`$
start do loop: $`k=0,1,\mathrm{}`$
$`\stackrel{}{z}^{(k)}`$ $`=`$ $`M\stackrel{}{p}^{(k)},\alpha ^{(k)}=\gamma ^{(k)}/(\stackrel{}{z}^{(k)},\stackrel{}{p}^{(k)})`$
$`\stackrel{}{\varphi }^{(k+1)}`$ $`=`$ $`\stackrel{}{\varphi }^{(k)}+\alpha ^{(k)}\stackrel{}{p}^{(k)}`$
$`\stackrel{}{r}^{(k+1)}`$ $`=`$ $`\stackrel{}{r}^{(k)}\alpha ^{(k)}\stackrel{}{z}^{(k)}`$
$`\stackrel{}{z}^{(k+1)}`$ $`=`$ $`\mathrm{\Delta }^1\stackrel{}{r}^{(k+1)}`$
$`\gamma ^{(k+1)}`$ $`=`$ $`(\stackrel{}{z}^{(k+1)},\stackrel{}{r}^{(k+1)})`$
if $`(\gamma ^{(k+1)}<\epsilon )\text{exit do loop}`$
$`\stackrel{}{p}^{(k+1)}`$ $`=`$ $`\stackrel{}{z}^{(k+1)}+{\displaystyle \frac{\gamma ^{(k+1)}}{\gamma ^{(k)}}}\stackrel{}{p}^{(k)}`$
end do loop
Here $`(,)`$ denotes the scalar product.
To perform the additional matrix multiplication with $`\mathrm{\Delta }^1`$ we used two fast Fourier transformations $``$, due to $`(\mathrm{\Delta })^1=^1q^2(k)`$. The performance we achieved is presented in Table 2. We conclude that on larger lattice sizes the reduction of iterations is about 70-75%, while the resulting reduction of CPU time depends on the lattice size. This is because we are using the fast Fourier transformations in a parallel CPU environment. If the ratio of used processors to the lattice size is small (see e.g. the data for $`32^4`$ lattice at this table), almost the same reductions of CPU time as for the number of iterations is achieved.
Further improvement may be achieved by using the multigrid Poisson solver to solve $`\mathrm{\Delta }\stackrel{}{z}^{(k)}=\stackrel{}{r}^{(k)}`$. This method is supposed to perform better on parallel machines. Perhaps a further improvement is possible by using as pre-conditioning matrix $`\stackrel{~}{M}^1=\mathrm{\Delta }^1\mathrm{\Delta }^1M_1\mathrm{\Delta }^1+\mathrm{}`$ which is an approximation of the F-P operator $`M=\mathrm{\Delta }+M_1`$ to a given order Zwanziger (1994) (see also Furui and Nakajima (2004b)). However, the larger the order, the more matrix multiplications per iteration cycle are required. This may reduce the overall performance. We have not checked so far which is the optimal order.
## ACKNOWLEDGMENTS
All simulations have been done on the IBM pSeries 690 at HLRN. We thank R. Alkofer for discussions and H. Stรผben for contributing parts of the program code. We are indebted to C. Fischer for communicating us his DSE results Fischer et al. (2005) prior to publication. This work has been supported by the DFG under contract FOR 465. A. Sternbeck acknowledges support of the DFG-funded graduate school GK 271. |
warning/0506/quant-ph0506087.html | ar5iv | text | # Dissipation and decoherence in a quantum oscillator
## I I. Introduction
In a recent paper in this journal, van KampenNGK has re-examined dissipation and noise in a quantum oscillator, treating it as a sub-system coupled to an environment. In working out his model, he has introduced some nice methods. They are slightly simplified and modified in this paper to revisit the closely related problem of quantum coherence and decoherence.
It is worth emphasizing that the harmonic oscillator is particularly simple, so that the analysis given here is not generalizable to more interesting systems. However, it may make up in explicitness what it lacks in generality. Indeed, the ultimate aim of this research is to further clarify the โpositivity problemโ in time dependent quantum statistical equations.VA ; SP The present work is a first step in that direction, in the hope that deriving largely known results in a simple way will clear the path.
This simplicity may offer amusement if not instruction to my long-time friends and colleagues Jim Langer and Pierre Hohenberg, despite the burdens of their high offices. It is a pleasure and honor to dedicate the paper to them.
## II II. Model and Preliminaries
The model is described by the Hamiltonian
$$H=\frac{1}{2}[P_0^2+\mathrm{\Omega }_0^2Q_0^2]+\frac{1}{2}_k\{P_k^2+\omega _k^2[Q_k+(\alpha _kQ_0/\omega _k^2)]^2\},$$
(1)
where the oscillator labeled $`0`$ will be called โthe sub-systemโ and the others โthe environment.โ The quantum mechanical (and classical) equations of motion obtained from Eq. (1) are
$`\ddot{Q}_0+\left[\mathrm{\Omega }_0^2+{\displaystyle \underset{k}{}}(\alpha _k/\omega _k)^2\right]Q_0`$ $`+`$ $`{\displaystyle \underset{k}{}}\alpha _kQ_k=0`$
$`\ddot{Q}_k+\omega _k^2Q_k+\alpha _kQ_0`$ $`=`$ $`0.`$ (2)
Fourier transforming Eq.(2) and eliminating $`Q_k`$, one obtains
$$g^1(\omega +i0^+)Q_0=0,\mathrm{where}$$
(3)
$$g^1(z)=z^2\mathrm{\Omega }_0^2\underset{k}{}\alpha _k^2\left(\frac{1}{z^2\omega _k^2}+\frac{1}{\omega _k^2}\right),$$
(4)
$`z`$ is a complex variable, and the $`i0^+`$ in Eq.(3) introduces the causal boundary condition. As pointed out in ref.NGK , where it is called $`G`$, $`g^1(z)`$ has zeros on the real $`z`$ axis corresponding to the normal mode frequencies, $`\omega _\nu `$, of the coupled system of oscillators, and $`Q_0(\omega _\nu )`$ is the amplitude of the sub-system oscillator in the $`\nu `$th mode.
Now, ohmic dissipation LC requires
$$\frac{\pi }{2}\underset{k}{}(\alpha _k^2/\omega _k)\delta (\omega \omega _k)J(\omega )=\eta \omega $$
(5)
for $`\omega `$ less than an upper cut-off $`\omega _c`$. Substituting this form into Eq. (4) yields, for $`\omega _c\omega `$,
$`g^1(\omega +i0^+)`$ $`=`$ $`\omega ^2\mathrm{\Omega }_0^2{\displaystyle \frac{2\eta }{\pi }}{\displaystyle _0^{\omega _c}}๐\overline{\omega }{\displaystyle \frac{\omega ^2}{(\omega +i0^+)^2\overline{\omega }^2}}`$ (6)
$`=`$ $`\omega ^2\mathrm{\Omega }_0^2+i\omega \eta ,`$
demonstrating very explicitly that Eqs. (1) and (5) do indeed construct a linearly dissipative environment, with damping constant $`\eta `$, without changing the system frequency $`\mathrm{\Omega }_0`$.
In ref.NGK , it is noted that the orthogonal normal mode transformation matrix $`X`$ defined by
$$Q_0=\underset{\nu }{}X_{0\nu }q_\nu Q_k=\underset{\nu }{}X_{k\nu }q_\nu ,$$
(7)
where the $`q_\nu `$s are the normal co-ordinates, is obtainable from the Green function given in Eq. (4). The normalizations
$`{\displaystyle \underset{\nu }{}}q_\nu ^2`$ $`=`$ $`1\mathrm{and}`$
$`Q_0^2+{\displaystyle \underset{k}{}}Q_k^2=Q_0^2[1`$ $`+`$ $`{\displaystyle \underset{k}{}}\alpha _k^2/(\omega _k^2\omega ^2)^2]=1,`$ (8)
show that the amplitudes corresponding to the mode $`\nu `$, and thus the matrix elements of the transformation, are given by
$`{\displaystyle \frac{1}{X_{0\nu }}}`$ $`=`$ $`\sqrt{1+{\displaystyle \underset{k}{}}\alpha _k^2/(\omega _k^2\omega _\nu ^2)^2}`$
$`=`$ $`\sqrt{{\displaystyle \frac{1}{2\omega _\nu }}{\displaystyle \frac{g^1}{\omega _\nu }}}\mathrm{and}`$
$`X_{k\nu }`$ $`=`$ $`{\displaystyle \frac{\alpha _k}{\omega _\nu ^2\omega _k^2}}X_{0\nu }.`$ (9)
One also learns from ref.NGK that Eq.(9) may be used to deftly perform sums over normal modes. The complex function $`g(z)`$ has poles only on the real axis and thus the spectral representation
$$g(z)=_{\omega _c}^{\omega _c}\frac{d\omega }{2\pi }\frac{s(\omega )}{z\omega }.$$
(10)
Using the explicit form for $`X_{0\nu }`$ given in the second line of Eq.(9) one obtains a formula that will be useful later in this work:
$`{\displaystyle \underset{\nu }{}}X_{o\nu }^2F(\omega _\nu )`$ $`=`$ $`2{\displaystyle \frac{dz}{2\pi i}zg(z)F(z)}`$ (11)
$`=`$ $`2{\displaystyle _0^{\omega _c}}{\displaystyle \frac{d\omega }{2\pi }}\omega s(\omega )F(\omega ).`$
The contour surrounds the real axis, where the function F has been assumed to be regular and zero for $`\omega <0`$; it has been evaluated using Eq.(10).
Note also from Eq.(6) and its complex conjugate that for an ohmic environment
$$s(\omega )=\frac{2\omega \eta }{(\omega ^2\mathrm{\Omega }_0^2)^2+\omega ^2\eta ^2}.$$
(12)
## III III. Time evolution of the reduced density matrix
To carry out the program of this section, initial conditions must be specified. I shall use those of FOC and GvDA , which are a special case of ones originated, to the best of my knowledge, in ref.VHVA . Assume that at $`t=0`$ complete thermal equilibrium is disturbed by a real โaperture functionโ $`\alpha (Q_0)`$. The entire system is then allowed to evolve to time $`t`$, and projected on to position states of the sub-system. The resulting reduced density matrix is given by ($`\mathrm{}=1`$)
$`\rho (Q_{0f}^{},Q_{0f}^{\prime \prime },t)`$
$`Tr\{|Q_{0f}^{\prime \prime }Q_{0f}^{}|\mathrm{e}^{iHt}\alpha (Q_0)\rho _{th}(H)\alpha (Q_0)\mathrm{e}^{+iHt}\}.`$
In this equation, the primed quantities are ordinary numbers, the unprimed ones operators, $`Tr`$ indicates a trace over all states of $`H`$, and $`\rho _{th}\mathrm{exp}\{\beta H\}/Tr\mathrm{exp}\{\beta H\}`$ with $`\beta `$ the reciprocal temperature. Fourier transform the aperture function
$$\alpha (Q_0)=๐a\stackrel{~}{\alpha }(a)\mathrm{e}^{iQ_0a},$$
(14)
and express the projection operator in Eq.(13) as
$`|Q_{0f}^{\prime \prime }Q_{0f}^{}|`$ $`=`$ $`{\displaystyle ๐u๐vf(u,v)\mathrm{e}^{iP_0u}\mathrm{e}^{iQ_0v}}\mathrm{with}`$
$`f(u,v)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{e}^{iQ_{0f}^{}v}\delta (Q_{0f}^{}Q_{of}^{\prime \prime }u),`$ (15)
proved by taking matrix elements of both sides, to see that
$`\rho (Q_{0f}^{},Q_{0f}^{\prime \prime },t)`$ $`=`$ $`{\displaystyle ๐a๐b๐u๐v\stackrel{~}{\alpha }(b)\stackrel{~}{\alpha }(a)f(u,v)๐ฏ},`$
$`\mathrm{with}๐ฏ`$ $``$ $`\mathrm{e}^{iQ_0b}\mathrm{e}^{iP_0(t)u}\mathrm{e}^{iQ_0(t)v}\mathrm{e}^{iQ_0a}.`$ (16)
Here the brackets mean an average with respect to $`\rho _{th}`$, and operators with a time argument are in the Heisenberg picture. A somewhat simpler average is done fairly heroically in ref.NGK using properties of Laguerre polynomials. However, since the sub-system and environment are all harmonic, a low-brow method is available. A single simple harmonic oscillator obeys the well known Debye-Waller identity for thermal averages,
$$\mathrm{e}^{iqc}=\mathrm{e}^{{\scriptscriptstyle \frac{1}{2}}c^2q^2},$$
(17)
where $`q`$ is the position operator and $`c`$ a number. This is reviewed in the Appendix, where it is also shown that a straightforward generalization yields err
$`\mathrm{ln}๐ฏ=`$ $``$ $`\frac{1}{2}[(a+b)^2+v^2]Q_0^2\frac{1}{2}u^2P_0^2buQ_0P_0(t)`$ (18)
$``$ $`bvQ_0Q_0(t)uvP_0Q_0uaP_0(t)Q_0`$
$``$ $`avQ_0(t)Q_0.`$
Since the co-ordinate $`Q_0`$ and the momentum $`P_0`$ are linearly related to the normal mode $`q_\nu `$s and $`p_\nu `$s via the known $`X_{0\nu }`$s the correlators in Eq.(18) are readily calculable, thereby formally completing the task of this section.
## IV IV. An example
To illustrate the usefulness of these methods, consider $`Q_0^2`$, one of the averages occurring in Eq.(18). From Eq.(7)
$`Q_0^2`$ $`=`$ $`{\displaystyle \underset{\nu }{}}{\displaystyle \underset{\nu ^{}}{}}X_{0\nu }X_{0\nu ^{}}q_\nu q_\nu ^{}={\displaystyle \underset{\nu }{}}X_{0\nu }^2q_\nu ^2`$ (19)
$`=`$ $`{\displaystyle \underset{\nu }{}}X_{0\nu }^2\left({\displaystyle \frac{1}{2\omega _\nu }}\mathrm{coth}\frac{1}{2}\beta \omega _\nu \right),`$
because the $`\nu `$s are independent oscillators, and $`q_\nu =0`$. The sum can now be transformed using Eqs.(11) and (12). It is easy to do analytically at zero temperature ($`\beta =\mathrm{}`$). Define the real part of the damped oscillator frequency via $`\mathrm{\Omega }^2=\mathrm{\Omega }_0^2\eta ^2/4`$ and factorize the denominator in Eq. (12) to obtain
$`Q_0^2_{T=0}`$ (20)
$`=`$ $`{\displaystyle \frac{\eta }{2\mathrm{\Omega }}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}\left[{\displaystyle \frac{1}{(\omega \mathrm{\Omega })^2+\eta ^2/4}}{\displaystyle \frac{1}{(\omega +\mathrm{\Omega })^2+\eta ^2/4}}\right]`$
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Omega }}}\left[1{\displaystyle \frac{2}{\pi }}\mathrm{arctan}{\displaystyle \frac{\eta }{2\mathrm{\Omega }}}\right].`$
The last is a known answer for the dissipation-induced squeezing by an ohmic bath. LC ; va This effect is at the root of the effect of damping in reducing the rate of escape from a metastable well.MAR
## V V. Time dependence of the probability
The position-space probability for the sub-system $`P(x,t)`$ is obtained by setting $`Q_{0f}^{}=Q_{0f}^{\prime \prime }=x`$ in Eq.(16), whereupon Eq.(15) requires that the variable $`u0`$ in Eq.(18). The time dependence can then be completely described by
$$[Q_0(t)Q_0(0)]Q_0(0)C(t)+iA(t),$$
(21)
this being the notation used in ref.VHVA in the free-particle limit, with the sign as corrected in footnoteerr . One can now do the (gaussian) integrals in Eq.(16) to obtain, with $`x^{}Q_{0i}^{}Q_{0i}^{\prime \prime },X^{}\frac{1}{2}(Q_{0i}^{}+Q_{0i}^{\prime \prime })`$,
$`P(x,t)`$ $`=`$ $`{\displaystyle ๐Q_{0i}^{}๐Q_{0i}^{\prime \prime }\alpha (Q_{0i}^{})\alpha (Q_{0i}^{\prime \prime })J(x,x^{},X^{},t)},`$
$`J`$ $`=`$ $`{\displaystyle \frac{1}{4\pi A(t)}}{\displaystyle \frac{1}{\sqrt{2\pi Q_0^2}}}\mathrm{exp}[i{\displaystyle \frac{x^{}}{2A(t)}}(xX^{})]`$
$`\times `$ $`\mathrm{exp}`$ $`\left[{\displaystyle \frac{x_{}^{}{}_{}{}^{2}C(t)}{4A^2(t)}}+{\displaystyle \frac{1}{2Q_0^2}}(X^{}i{\displaystyle \frac{x^{}C(t)}{2A(t)}})^2\right].`$ (22)
In the free-particle limit $`\mathrm{\Omega }_00`$, $`Q_0^2\mathrm{}`$ and the normalized gaussian in $`X^{}`$ must be replaced, to preserve dimensions and normalization, by $`L^1`$, where $`L`$ is the size of the system, yielding
$$J=\frac{1}{4\pi A(t)L}\mathrm{exp}\left[i\frac{x^{}}{2A(t)}(xX^{})\frac{x_{}^{}{}_{}{}^{2}C(t)}{4A^2(t)}\right].$$
(23)
Using the method of Section II, known results for a free particle follow from Eq. (21):
$`A(t)`$ $`=`$ $`{\displaystyle \frac{\eta }{\pi }}{\displaystyle _0^{\omega _c}}๐\omega {\displaystyle \frac{\mathrm{sin}\omega t}{\omega (\omega ^2+\eta ^2)}}={\displaystyle \frac{1}{2\eta }}(1\mathrm{e}^{\eta t}),`$
$`C(t)`$ $`=`$ $`{\displaystyle \frac{\eta }{\pi }}{\displaystyle _0^{\omega _c}}๐\omega \mathrm{coth}{\displaystyle \frac{\beta \omega }{2}}{\displaystyle \frac{(1\mathrm{cos}\omega t)}{\omega (\omega ^2+\eta ^2)}}`$
$``$ $`(\beta \eta 1){\displaystyle \frac{1}{\beta \eta ^2}}(1\mathrm{e}^{\eta t}\eta t).`$
## VI VI Decoherence
Eq.(22)โwhich is not restricted to ohmic dissipationโis to my knowledge new, and explicit enough to allow a general study of time development in this dissipative quantum system. Previous attempts VA ; GOA in which I have been involved have for technical reasons been restricted to high temperature, a limitation which it should here be possible to avoid. Work, in progress, in this direction would seem to be justified by experimentalMOOIJ and theoreticalGrifoni interest in quantum information storage.
The free particle limit, Eqs.(23,24), has been derived by many different methodsโperhaps none as straightforward as the one here given. At finite $`\eta `$, these equations (uncontroversially) display decoherence. This is particularly well demonstrated in an example introduced in refFLO . Here the aperture function is taken to be a sum of two Gaussians, each of width $`\sigma `$ and separated by a distance $`d`$. The probability given by Eqs. (23) and (24) can then be written as a sum of the probabilities from each slit alone (sum term) and an interference contribution depending sinusoidally on a time dependent phase (interference term.) There is an unresolved controversy GvDA ; FOC in the published literature about what is meant by decoherence in this completely well defined problem. No one would doubt that the amplitude of the interference term is a measure of coherence. In ref.FLO and other publicationsFC an โattenuation coefficientโ is introduced which is equivalent to the amplitude of the interference term divided by the sum term, evaluated at the mid-point between the slits. This quantity is by construction unity at $`t=0`$. It decreases rapidly with time. At high temperatures, in a sense to be made clear below, it drops to zero, even when the environmental coupling is eliminated. This is interpreted in refs.FLO ; FC as โDecoherence without Dissipation.โ
In a Comment GvDA it is suggested that the measure of decoherence used in FLO does not distinguish between the loss of phase information and the spreading of wave packets on time scales less than the mixing time $`t_{mix}=2\sigma d(m/\mathrm{})`$. Here $`\sigma `$ is the width of each slit, $`d`$ the spacing between them, and I have reintroduced Planckโs constant $`\mathrm{}`$ and the particle mass $`m`$ to make the dimensions transparent.
This has been vigorously rebutted in a ReplyFOC .
Rather than repeat these arguments, I refer the interested reader to them. However, since many readers may be intrigued by the idea of decoherence without dissipation, I close this paper with 4 figures which show that there is no evidence for any such thing in the uncontroversial Eqs.(23,24). Since the disagreement occurs in the limit of no dissipation, consider this case at various temperaturesโgiven by the ratio of the thermal de Broglie wavelength $`l_{th}=\sqrt{\beta /2}(\mathrm{}/\sqrt{m})`$ to $`d`$. The figures, in which $`\sigma /d=0.05`$, show that at low temperatures, $`l_{th}d`$ there is coherence without decoherence. At higher temperatures, there is no coherence at all, even on the short time scale
$$\tau _{FLO}=\sqrt{8\beta m}\frac{\sigma ^2}{d}=2t_{mix}\frac{l_{th}}{d}\frac{\sigma }{d}$$
(25)
introduced in FOC , and thus nothing to decohere. In this limit, coherence is already destroyed by the Hakim-Ambegaokar initial condition.
Added Note. After this paper was completed, I was made aware of FMDG in which normal co-ordinates are used to treat the problem of many non-interacting fermions coupled to a disspative environment in a harmonic oscillator.
## VII Appendix
A wonderfully short proof of the identity for thermal averages
$$\mathrm{exp}[\underset{i}{}d_ia_i+c_ia_i^{}]=\mathrm{exp}[\frac{1}{2}_i(d_ia_i+c_ia_i^{})^2],$$
(26)
where $`a_i`$, $`a_i^{}`$ are boson annihilation and creation operators and the subscript $`i`$ refers to independent harmonic oscillators, is given in a single mannerist, if not rococo, sentence by MerminNDM . Note that, since all the operators in Eq.(16) have c-number commutators, it can be put in the form of the left hand side of Eq.(26) for the subsystem oscillator labeled $`0`$, using the Baker-Haussdorf identity for such operators: $`\mathrm{e}^A\mathrm{e}^B=\mathrm{e}^{[A+B]}\mathrm{e}^{\frac{1}{2}[A,B]}`$. Now, express $`P_0`$ and $`Q_0`$ in terms of the (normal) co-ordinates, of the independent oscillators $`\nu `$. The Debye-Waller identity Eq. (26) then has in the exponent a sum over correlators for each $`\nu `$. Note that as in Eq. (19), a single sum can be replaced by a double sum, because $`q_\nu =p_\nu =0`$, yielding Eq.(18).
## VIII acknowledgements
The figures in this paper were created using Mathematica by Dominique Gobert. I am grateful to him and Jan von Delft for considerable help, and to Frank Wilhelm for background information. This work was started during the summer of 2004 at the Aspen Center for Physics. where I had useful discussions with Vladimir Privman and Dima Mozyrsky. Support from the NSF under grant DMR- 0242120 is acknowledged with thanks. |
warning/0506/cond-mat0506503.html | ar5iv | text | # Mirage phenomena in superconducting quantum corrals
## I Introduction
Understanding impurity effects was one of the early, classical tasks for the theory of superconductivity Abrikosov63 ; Schrieffer64 . The Anderson theorem explained why non-magnetic impurity scatterers do not destroy conventional, isotropic $`s`$-wave superconductivity Anderson59 . For unconventional superconductors, however, both non-magnetic and magnetic impurities are pairbreaking and the sensitivity to impurities and the concomitant suppression of superconductivity may even serve as a fingerprint of an unconventional pairing state. It took three decades after the development of the BCS theory of superconductivity, before experimental studies of impurity effects on microscopic length scales became possible โ mostly due to the advance of scanning tunneling microscopy (STM). The need for experimental probes with a high spatial resolution arose from the discovery of a continuously growing number of unconventional singlet and even triplet superconductors; ruthenates and high-temperature cuprate superconductors belong to the most prominent examples.
During the last decade it has furthermore become possible to manipulate surface structures on the atomic level. In particular the achievements in the research group of D. Eigler led to the design of closed atomic arrangements of adatoms on metallic surfaces Crommie93 ; Heller94 . In these quantum corrals with circular or elliptical shape electrons are confined to well defined geometries and STM techniques were applied to explore the structure of quantum mechanical wavefunctions and their interference patterns Heller94 ; Kliewer00 . Quantum mirage phenomena are observed, when an additional impurity atom is placed at one focus of an elliptic corral Manoharan00 . The terminology of quantum mirages is hereby defined by the observation that mirror images of the local density of states pattern around the impurity at one focus point appear also at the second impurity free focus. In subsequent studies even Kondo resonances could be detected at the impurity-free focus point, if a magnetic impurity adatom was placed in the other focus Aligia05 .
Similar mirage phenomena are expected for elliptic corrals on superconducting surfaces Morr . The local suppression of the superconducting order parameter around an impurity at a selected point in the corral should be transferred to image structures due to standing wave patterns extending over the entire area of the elliptic corral. Superconducting quantum corrals therefore provide unique systems to study the combined effects of impurity induced fingerprint structures of the pairing state and the mirage phenomena of interfering quantum mechanical waves. With the assumption that a superconducting state is established in the bulk and on the surface of the substrate material we discuss in this paper the expected phenomena within the framework of BCS theory and the Bogoliubov-de Gennes equations.
## II The model
We consider the following mean-field Hamiltonian for an isotropic and an anisotropic singlet superconductor, respectively, with magnetic or non-magnetic impurities
$`H={\displaystyle \underset{\sigma }{}}{\displaystyle \mathrm{d}^2๐ซ\psi _\sigma ^{}(๐ซ)H_0\psi _\sigma (๐ซ)}+H_{s/d}+H_{imp}`$ (1)
with $`H_0=(\frac{\mathrm{}^2}{2m}^2\mu )+V(๐ซ)`$. $`\mu `$ is the chemical potential, and $`V(๐ซ)`$ is a hard-wall potential, which confines the electrons to the interior of an elliptic corral. The Schrรถdinger equation for particles in an elliptic geometry was previously solved analytically; the resulting eigenfunctions are a combination of Mathieu functions $`ce_r`$, $`se_r`$ and modified Mathieu functions $`Ce_r`$, $`Se_r`$ Schmid ; Porras :
$`\phi _{r,n_c}^c(\theta ,\eta )`$ $`=`$ $`ce_r(\theta ,k_n^c)Ce_r(\eta ,k_n^c),`$ (2)
$`\phi _{r,n_s}^s(\theta ,\eta )`$ $`=`$ $`se_r(\theta ,k_n^s)Se_r(\eta ,k_n^s),`$ (3)
where $`\theta `$ and $`\eta `$ are elliptical coordinates. $`r`$, $`n_{c(s)}`$ enumerate the quantum numbers for the eigenstates. Their eigenenergies Schmid
$`ฯต={\displaystyle \frac{2\mathrm{}^2}{(ae)^2m}}k_{r,n}^{c/s}`$ (4)
are determined by the zeros of the modified Mathieu functions $`k_{r,n}^{c/s}`$. $`a`$ denotes the length of the semimajor axis and $`e`$ the eccentricity of the ellipse; $`m`$ is the electron mass.
For an isotropic $`s`$-wave superconductor with an attractive contact interaction $`H_s`$ is given by
$`H_s={\displaystyle }\mathrm{d}^2๐ซ[\psi _{}^{}(๐ซ)\psi _{}^{}(๐ซ)\mathrm{\Delta }(๐ซ)+h.c.],`$ (5)
where $`\mathrm{\Delta }(๐ซ)=g\psi _{}(๐ซ)\psi _{}(๐ซ)`$ is the order parameter and $`g>0`$ the pairing interaction strength. Similarly, the mean-field Hamiltonian for a singlet superconductor with a spatially extended pairing interaction is described by
$`H_d={\displaystyle }\mathrm{d}^2๐ซ{\displaystyle }\mathrm{d}^2๐ซ^{}[\psi _{}^{}(๐ซ)\psi _{}^{}(๐ซ^{})\mathrm{\Delta }(๐ซ,๐ซ^{})+h.c.],`$ (6)
where the order parameter $`\mathrm{\Delta }(๐ซ,๐ซ^{})=g(๐ซ,๐ซ^{})\psi _{}(๐ซ^{})\psi _{}(๐ซ)`$ depends on both coordinates $`๐ซ`$ and $`๐ซ^{}`$ of the electrons forming the Cooper pair. Our ansatz for the pairing interaction is $`g(๐ซ,๐ซ^{})=g\delta (|๐ซ๐ซ^{}|R)`$ and assumes an attraction for electrons at a distance $`R`$; this lengthscale $`R`$ should be set by the size of a typical crystal lattice constant. Although $`g(๐ซ,๐ซ^{})`$ is an isotropic real space interaction, it gives rise to an anisotropic local order parameter as we discuss below. The impurity potentials for magnetic ($``$) and non-magnetic ($`+`$) impurities are described by
$`H_{imp}={\displaystyle \mathrm{d}^2๐ซU(๐ซ)\left[\psi _{}^{}(๐ซ)\psi _{}(๐ซ)\pm \psi _{}^{}(๐ซ)\psi _{}(๐ซ)\right]},`$ (7)
where $`U(๐ซ)=_i^nU_0\delta (๐ซ๐ซ_๐ข)`$ with impurity positions $`๐ซ_๐ข`$. In this model for the magnetic scattering center the impurity spins are treated as classical spins, $`S1`$ Shiba ; Schrieffer , which is a reasonable assumption e.g. for Mn and Gd adatom impurities on a niobium surface Yazdani .
The mean-field Hamiltonian Eq. (1) is diagonalized by the transformation
$`\psi _\sigma (๐ซ)`$ $`=`$ $`{\displaystyle \underset{n}{}}\left[u_n(๐ซ)\alpha _{n\sigma }+sgn(\sigma )v_n(๐ซ)\alpha _{n\sigma }^{}\right],`$ (8)
$`\alpha _{n\sigma }`$ $`=`$ $`{\displaystyle \mathrm{d}^2๐ซ\left[u_n(๐ซ)\psi _\sigma (๐ซ)sgn(\sigma )v_n(๐ซ)\psi _\sigma ^{}(๐ซ)\right]},`$ (9)
where $`\alpha _{n\sigma }`$ describes a superconducting quasiparticle with spin $`\sigma `$ and energy $`E_n`$. The ansatz (8) leads to the Bogoliubov-de Gennes equations
$`\left[H_0+U(๐ซ)\right]u_n(๐ซ)+{\displaystyle \mathrm{d}^2๐ซ^{}\mathrm{\Delta }(๐ซ,๐ซ^{})v_n(๐ซ^{})}=E_nu_n(๐ซ),`$ (10)
$`{\displaystyle \mathrm{d}^2๐ซ^{}\mathrm{\Delta }(๐ซ,๐ซ^{})u_n(๐ซ^{})}\left[H_0\pm U(๐ซ)\right]v_n(๐ซ)=E_nv_n(๐ซ),`$ (11)
where the +/- sign in Eq. (11) corresponds to non-magnetic/magnetic impurities, respectively. These equations are simplified in the case of an isotropic superconductor, where the pairing-interaction is pointlike, and therefore $`\mathrm{\Delta }(๐ซ,๐ซ^{})=\delta (๐ซ๐ซ^{})\mathrm{\Delta }(๐ซ)`$. The order parameter $`\mathrm{\Delta }(๐ซ)`$ is determined self-consistently from the condition
$`\mathrm{\Delta }(๐ซ)=g{\displaystyle \underset{n}{}}u_n(๐ซ)v_n(๐ซ)\left[12f(E_n)\right],`$ (12)
where $`f(E)`$ denotes the Fermi function. For an anisotropic superconductor with a finite range pairing interaction we have instead
$`\mathrm{\Delta }(๐ซ,๐ซ^{})=g(๐ซ,๐ซ^{}){\displaystyle \underset{n}{}}\{u_n(๐ซ)v_n(๐ซ^{})[1f(E_n)]`$ (13)
$`u_n(๐ซ^{})v_n(๐ซ)f(E_n)\}.`$
In order to solve the Bogoliubov-de Gennes equations (10) and (11) we expand $`u_n(๐ซ)`$ and $`v_n(๐ซ)`$ in terms of free electron eigenfunctions inside the elliptic corral Schmid
$`u_n(๐ซ)`$ $`=`$ $`{\displaystyle \underset{k}{}}u_{kn}\phi _k(๐ซ),`$ (14)
$`v_n(๐ซ)`$ $`=`$ $`{\displaystyle \underset{k}{}}v_{kn}\phi _k(๐ซ),`$ (15)
where $`k`$ enumerates all the eigenstates given in Eqs. (2) and (3). For the ground-state energies of the isotropic (Eq. (16)) and the anisotropic (Eq. (17)) superconductor we find
$`E_g^s=`$ $`\mathrm{\hspace{0.33em}2}{\displaystyle \underset{n}{}}{\displaystyle \mathrm{d}^2๐ซ\left[v_n(๐ซ)H_0v_n(๐ซ)\mathrm{\Delta }(๐ซ)u_n(๐ซ)v_n(๐ซ)\right]},`$ (16)
$`E_g^d=`$ $`{\displaystyle \underset{n}{}}{\displaystyle }\mathrm{d}^2๐ซ[2v_n(๐ซ)H_0v_n(๐ซ)`$ (17)
$``$ $`{\displaystyle }\mathrm{d}^2๐ซ^{}\mathrm{\Delta }(๐ซ,๐ซ^{})(u_n(๐ซ)v_n(๐ซ^{})+u_n(๐ซ^{})v_n(๐ซ))].`$
The local density of states (LDOS) is obtained from
$`N(๐ซ,\omega )={\displaystyle \underset{n}{}}\left[u_n^2(๐ซ)\delta (\omega E_n)+v_n^2(๐ซ)\delta (\omega +E_n)\right],`$ (18)
and the total density of states (DOS) is $`N(\omega )=\mathrm{d}^2๐ซN(๐ซ,\omega )`$.
## III Results
For the elliptic corral we choose its eccentricity $`e=0.5`$ and the length of the semi-major axis $`a=150`$ ร
; these two parameters determine the single-particle eigenenergies and -states. With an approximate electronic density of one electron per 6ร
$`\times `$6ร
square the total number of electrons in the corral is $`N=1726`$. With this choice the $`863^{rd}`$ eigenstate at the Fermi energy $`ฯต_F`$ has a large probability density at the foci, which is a favorable situation for the occurrence of mirage phenomena. For an isotropic superconductor the mirage phenomenon is observed, whenever there is one of these non-zero probability density eigenstates lying inside the energy gap. For our choice of the pairing interaction strength $`g=1.0\times 10^3ฯต_F`$ the energy gap $`\mathrm{\Delta }_0`$ stretches over about 50 eigenstate energies, where several of these states have the required property. Hence, the existence of mirage phenomena for an isotropic superconducting surface does not depend sensitively on the position of the Fermi energy, in contrast to a metallic, normal conducting surface Schmid and also in contrast to an anisotropic superconductor, as we discuss below.
### III.1 Isotropic $`s`$-wave superconductor
We start our analysis with a corral on the surface of an isotropic $`s`$-wave superconductor. In the presence of non-magnetic impurities the system remains time-reversal invariant. These impurities do not break Cooper pairs, which are built from Kramers-degenerate electronic states Anderson59 . Consequently, no additional states are created in the renormalized energy gap. Magnetic impurities on the contrary, are pair-breaking defects, because the spin-up electron of the Cooper pair is repelled ($`U_0>0`$) and the spin-down electron is attracted by the impurity, leading to a quasiparticle bound state with a spin-down particle peak and a spin-up hole peak in the LDOS shown in Fig. 1 (dashed and dotted curves), where the impurity is located at the right focus of the ellipse. The LDOS at the impurity site vanishes at the gap edge energies, while the LDOS at the mirror point (dotted line) is finite at these energies. The total spectral weight of both peaks, i.e. the spatially and energy (between the gap edge energies $`\mathrm{\Delta }_0`$ and $`\mathrm{\Delta }_0`$) integrated LDOS is unity, characterizing one localized quasiparticle. Hence, the bound state corresponding to the dashed line is a local spin-down quasiparticle state (see also Fig. 2, unfilled triangles), where the spin projection of the particle component is antiparallel and that of the hole component is parallel to the impurity spin. This is the same phenomenon known for the extended $`s`$-wave superconductor without confining boundaries Schrieffer . The dash-dotted and dash-double-dotted curves in Fig. 1 show the LDOS for a stronger impurity potential. The difference now is that the spins of the particle and the hole peak have changed, i.e. we find a spin-up quasiparticle (compare with Fig. 2, filled triangles). This change is due to the fact that both peaks cross the Fermi energy at a critical impurity potential strength $`U_c`$.
In the corral geometry we observe the two-peak structure in the LDOS not only at the impurity site e.g. at one focus point, but also โ with some attenuation โ at the impurity free focus (dotted and dash-double-dotted lines) Morr . The mirage effect is only strong at one of the two peaks at the mirror point LDOS. There the particle peak carries very little spectral weight for $`U_0<U_c`$ (see Fig. 1, dotted line), and so does the hole peak for $`U_0>U_c`$. But with increasing $`U_0`$ the particle and hole peaks at the mirror point are both growing. We emphasize that the LDOS peaks at the mirror point are at the same resonance energies as in the LDOS at the impurity site, because they belong to the same bound state with energy $`\omega _0`$. This is the quantum mirage effect in the LDOS for an elliptic corral built on the surface of an isotropic $`s`$-wave superconductor. As expected, we find this effect to be very robust against a change in the Fermi energy and thus the electron density. Moreover this mirage effect remains almost unaffected by an additional perturbation in form of a second magnetic impurity placed anywhere else inside the corral.
Since the bound-state peaks at energies $`\omega =\pm \omega _0`$, split off from the gap edges, they move with increasing potential strength $`U_0`$ symmetrically towards the Fermi energy (see Fig. 2). $`\omega _0`$ therefore decreases until it reaches the chemical potential for a critical potential $`U_c`$, which depends on the impurity position $`๐ซ_0`$. At $`U_0=U_c`$ the bound state becomes a zero energy state, and the ground state of the superconductor becomes unstable. The existence of this critical point, which signals a ground-state level crossing transition, was first pointed out by Sakurai for $`s`$-wave superconductors with magnetic impurities Sakurai70 . At $`U_c`$ the system undergoes a first order phase transition from a spin zero to a spin $`\frac{1}{2}`$ ground state Schrieffer ; Sakurai70 . While the condensate ground state for $`U_0<U_c`$ consists only of Cooper pairs, the new ground state for $`U_0>U_c`$ contains an additional spin-down electron (or up-spin hole). The impurity induced bound state is now a spin-up quasiparticle (Fig. 1, dash-dotted and dash-double-dotted curves and Fig. 2, filled triangles). A further increase of $`U_0`$ in this regime ($`U_0>U_c`$) is accompanied by a spatial shift of the LDOS at the resonance energies away from the impurity position. This spatial redistribution of the spectral weight leads to decreasing heights of the bound-state peaks in the LDOS Fig. 1 at the impurity position (right focus) and an accompanying increase of the peak heights at the impurity free focus. Note that the total spectral weight of the bound-state peaks always remains unity. Hence we observe a transition from a mirage-effect for small $`U_0`$ to an anti-mirage-effect for large $`U_0`$.
We compare our self-consistent solutions of the Bogoliubov-de Gennes equation for a corral with the results for an open surface using the following formula for the bound-state energies Schrieffer
$`\omega _0=\mathrm{\Delta }_0{\displaystyle \frac{|1(\pi N_F)^2U_0^2|}{1+(\pi N_F)^2U_0^2}},`$ (19)
where $`N_F`$ is the normal state DOS at the Fermi energy. Eq. (19) is obtained from the non-self-consistent T-matrix formalism, which proved to be sufficient and even quantitatively accurate for the description of magnetic impurity induced bound states in $`s`$-wave superconductors without confining walls Schrieffer . In Fig. (2) we use Eq. (19) with $`\pi N_F\overline{N_F}`$ (solid curve), where $`\overline{N_F}`$ desribes the mean density of states averaged over an energy window $`ฯต_F\pm \mathrm{\Delta }_0`$ in the normal state. Fig. 2 shows that the bound-state energies in the elliptic corral are well described by Eq. (19). Comparing the self-consistent and non-self-consistent results for our corral, we find indeed only small differences: neither the quasiparticle energies are shifted notably nor the heights of the quasiparticle peaks are modified significantly.
Another interesting quantity is the renormalized local order parameter $`\mathrm{\Delta }(๐ซ)`$ given by Eq. (12). We observe a suppression of $`\mathrm{\Delta }(๐ซ)`$ at the impurity site, where the lengthscale of the suppression and the wavelength of the order parameter oscillations is given approximately by the Fermi wavelength $`\lambda _F=2\pi \mathrm{}/\sqrt{2mฯต_F}15`$ ร
, but seems to be independent of the size of the energy gap $`\mathrm{\Delta }_0`$. Remarkably, the local order parameter at the impurity position changes discontinuously from a positive to a negative value, when $`U_0`$ crosses the critical potential strength $`U_c`$ Schrieffer ; Flatte97 . For a non-magnetic impurity the reduction of $`\mathrm{\Delta }(๐ซ)`$ in the vicinity of the impurity is due to the repulsive potential for both spin directions resulting in a reduced electronic density at and near the impurity site and therefore to a reduced local pairing amplitude. In Fig. 3 a magnetic impurity is placed at the right focus point. All one-particle eigenstates with a large probability density at the right focus have a large probability density at the left focus as well. The difference plot of the order parameter Fig. 3 strongly resembles the $`863^{rd}`$ eigenstate, implying that $`\mathrm{\Delta }(๐ซ)`$ is suppressed not only at the impurity site (right focus), but almost symmetrically at the impurity free (left) focus, which we therefore call an order parameter mirage effect. We note that the order parameter is negative at the impurity site in Fig. 3.
If there are magnetic impurity spins in both foci, their antiparallel alignment leads to a lower ground-state energy Eq. (16). Moreover we observe that the order parameter suppression is about six times stronger for a parallel than for an antiparallel impurity spin alignment. A parallel alignment leads to a four-peak structure in the LDOS, because there are two similar states localized at each impurity, which are hybridizing and therefore splitting into bonding and antibonding states Flatte2000 ; Morr . If the impurity spins are aligned antiparallel, we find only two peaks inside the energy gap. The intensity of the hybridization depends in general on the distance between the two impurities; for a corral geometry this dependence is quite complex.
### III.2 Anisotropic superconductor
Now we assume the pairing interaction $`g(๐ซ,๐ซ^{})`$ to be attractive in a distance of a typical crystal lattice constant $`R0.04a=6`$ร
. In this case the order parameter $`\mathrm{\Delta }(๐ซ,๐ซ^{})`$ depends on the position inside the corral as well as on the relative coordinate. Although the interaction is still isotropic, the order parameter is now anisotropic. In Fig. 4 (left panel) we show the order parameter in the center of the ellipse, where we identify a $`d`$-wave like structure with sign changes and nodes at angles near $`45^{}`$ with respect to the semimajor axis. At the center the order parameter is least affected by the confining corral, but except for this point of highest symmetry, the structure of the local order parameter is distorted in a complex way. An example for a remarkably complicated symmetry and local structure of $`\mathrm{\Delta }(๐ซ,๐ซ^{})`$ is shown in the right panel of Fig. 4. The spatially extended pairing interaction leads to a V-shape-resembling DOS (see Fig. 5, solid curve), with no sharp gap edges. Furthermore, the DOS has no particle-hole symmetry, because of the asymmetric distribution of eigenenergies for the corral eigenfunctions.
While non-magnetic impurities do not give rise to localized states in a $`s`$-wave superconductor, they do have a pair-breaking effect in an anisotropic superconductor Gorkov . For a non-magnetic impurity placed at different positions inside the corral we calculate the local order parameter, the DOS, and the LDOS. In the DOS we always find bound-state peaks, which move towards the Fermi energy as $`U_0`$ increases, but never reach $`ฯต_F`$. Hence, we do not observe a zero-energy-peak in agreement with the results for an open surface in a particle-hole asymmetric situation Zhu ; Schrieffer . The DOS is modified over the entire energy range as shown in Fig. 5. For weak impurity potentials we observe a two-peak structure in the DOS near the Fermi energy, but after exceeding some critical potential $`\stackrel{~}{U_c}`$, which depends again on the impurity position $`๐ซ_\mathrm{๐}`$, each of the two peaks starts to split into two, so that a four-peak structure emerges near $`\omega =ฯต_F`$ (see Fig. 5). At the resonance energies for small $`U_0`$ (two-peak regime) we find, that the bound state has a large LDOS in the vicinity of the impurity. By increasing $`U_0`$ the original bound state acquires admixtures of the $`862^{nd}`$ and $`866^{th}`$ state Comment . Above the critical potential $`\stackrel{~}{U}_c`$, we identify the $`866^{th}`$ state either at both outer or both inner peaks of the four-peak low-energy DOS, while the LDOS at the two other peaks consists of the small $`U_0`$ bound state mixed with a contribution of the $`862^{nd}`$ state. Hence the splitting of the bound state energy corresponds to a โdemixingโ of the contributing states, which form the bound state. We observed this behavior at different impurity positions on the semimajor axis, but the occurrence of a four-peak structure nevertheless depends sensitively on the Fermi energy. Slight variations of the electronic densities can lead to two-peak structures in the DOS, irrespective of the impurity potential strength.
In order to show a few of a large variety of interesting and astonishing effects we calculate the LDOS in the presence of a non-magnetic impurity in particular at the right focus $`(0.5a,0)`$, at $`(0.44a,0)`$, and at the center (0,0). Interestingly, the important states ($`862^{nd}`$ and $`866^{th}`$) mentioned above are the only contributions to the bound state for an impurity in a focus point. These states have a zero LDOS on the entire semimajor axis, but several maxima about a lattice constant away from the focus. This bound state shows an almost perfect reflection symmetry with respect to the semiminor axis. Shifting the impurity a distance 1.5$`R`$ away from the right focus to $`(0.44a,0)`$, we observe that the LDOS at the lower resonance energies $`\omega _L^\pm =\pm 2.8\times 10^3ฯต_F`$ has a strong contribution of the $`863^{rd}`$ state and, as discussed before, a weaker contribution of the $`862^{nd}`$ state. Fig. 6 shows the LDOS at the positive resonance energy $`\omega _L^+`$. While the LDOS at the impurity site itself is zero, we observe a maximum in the LDOS at the mirror point. The impurity induces asymmetrically spectral weight around both foci, whereas the height is larger at the left than at the right focus, which is close to the impurity. Thus, we observe a surprising, strong anti-mirage effect in the LDOS at all resonance energies with a non-magnetic impurity at (0.44a, 0) independent of the potential strength $`U_0`$.
In Fig. 8 we show the LDOS with an impurity located at the center of the ellipse at the energy of the higher positive resonance energy $`\omega _H^+=7.6\times 10^3ฯต_F`$. We identify the $`862^{nd}`$ eigenstate again, which determines the region near the boundary of the ellipse and consists of rings of peaks aligned along four ellipses. The structure of the LDOS in the vicinity of the impurity is in fact more influenced by the symmetry of the local order parameter than the geometry of the problem. E. g. in a range of a few lattice constants around the center there is no LDOS along the diagonal directions, and locally there is almost a fourfold symmetry as for the absolute value of the order parameter in Fig. 4 (left panel) Haas .
In the vicinity of the impurity site the order parameter is reduced strongly by a non-magnetic impurity. This becomes obvious by comparing the right panels of Figs. 4 and 8, where the impurity is located at $`(0.44a,0)`$. Moreover we observe here and quite generally a modificaton of the spatial structure and the symmetry of the local order parameter, which depends sensitively on the potential strength $`U_0`$. Interestingly the new, impurity induced spatial structure resembles more a $`d`$-wave-like symmetry than the original structure. Fig. 8 clearly indicates, that this $`d`$-wave-like order parameter structure is projected to the impurity free mirror point. Instead of a reduction, the order parameter is actually enhanced at the mirror point. As in the LDOS we therefore also observe an anti-mirage effect for the order parameter. The order parameter at $`(0.44a,0)`$ is changed to a structure similar as in Fig. 8 (right panel), if an impurity is located at the center (!) of the ellipse. In general we find, that the magnitude of the order parameter is not only reduced over the entire area of the corral, but also modified with respect to its local spatial symmetry.
Regarding the connection between the spatial structure of the LDOS in the vicinity of the impurity and the order parameter, we encounter three different situations. (i) the LDOS resembles the order parameter in the absence of an impurity, (ii) the LDOS resembles the order parameter in the presence of an impurity, and (iii) there is no connection at all between the two. The latter is observed for example for a focus point impurity. We find the behavior (i) only for an impurity at the center of the corral, for large impurity potentials, and at positive resonance energies. For smaller impurity potentials we are back to case (ii), where again the LDOS shows the structure of the order parameter with an impurity only at positive resonance energies. This second case seems to reflect the more generic situation. Hence, the LDOS in the presence of a non-magnetic impurity inside the corral can directly reveal the structure and the local symmetry of the order parameter.
## IV Conclusion
We have observed a rich variety of mirage- and anti-mirage-phenomena not only in the LDOS, but also for the local order parameter structure. The induced patterns in both quantities in the vicinity of an impurity lead to mirror images in impurity-free regions, which are characteristic for both, the structure of the electronic wavefunctions in the elliptic geometry and the nature of the superconducting state. With the continuing advances of STM techniques and the already achieved capabilities for the design of atomic corral arrangements on metallic surfaces the realization of quantum corrals on superconducting surfaces indeed appears to be possible. From our model analysis we are led to expect intriguingly rich structures in the LDOS patterns, which by themselves contain information about the nature of the superconducting state.
ACKNOWLEDGEMENTS
One of the authors (A. P. K.) wishes to express his special gratitude for the guidance and support during and ever after his doctoral thesis work supervised by Bernhard Mรผhlschlegel. This work was supported by the Deutsche Forschungsgemeinschaft through SFB 484. |
warning/0506/astro-ph0506323.html | ar5iv | text | # The Ionized Gas and Nuclear Environment in NGC 3783 V. Variability and Modeling of the Intrinsic Ultraviolet AbsorptionBased on observations made with the NASA/ESA Hubble Space Telescope obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS 5-26555, and with the NASA-CNES-CSA Far Ultraviolet Spectroscopic Explorer, which is operated for NASA by the Johns Hopkins University under NASA contract NAS5-32985.
## 1. Introduction
Mass outflow, seen as blueshifted absorption in UV and X-ray spectra, is an important component of active galactic nuclei (see the recent review in Crenshaw et al., 2003). This โintrinsic absorptionโ is ubiquitous in nearby AGNs, appearing in over half of Seyfert 1 galaxies having high-quality UV spectra obtained with the Hubble Space Telescope (HST) (Crenshaw et al., 1999) and the Far Ultraviolet Spectroscopic Explorer (FUSE) (Kriss, 2002). Spectra from the Advanced Satellite for Cosmology and Astrophysics (ASCA) identified X-ray โwarm absorbersโ, modeled as absorption edges, in a similar percentage of objects (Reynolds, 1997; George et al., 1998), and many studies have explored the connection between the absorption observed in the X-ray and UV bandpasses (e.g. Mathur et al., 1994). Large total ejected masses have been inferred for these outflows, exceeding the accretion rate of the central black hole in some cases, indicating mass outflow may play an important role in the overall energetics in AGNs (e.g. Mathur et al., 1995; Reynolds, 1997).
Variability in the intrinsic UV absorption in Seyfert galaxies is common. All objects with high-resolution UV spectra obtained at multiple epochs exhibit substantial variations in their absorption strengths (Crenshaw et al., 2003, and references therein). Absorption variability could result from: (a) a response to changes in the ionizing AGN flux such that the total column density of the absorber remains constant but the ionization structure changes, (b) regions of condensation/evaporation in our line-of-sight to the AGN emission sources due to thermal perturbations, or (c) bulk motion of the absorber transverse to our line-of-sight. In the latter case, the observed equivalent widths could vary due to either a change in the line-of-sight covering factor of the background emission, or a change in the total column density of gas seen due to a shifting of different regions of the outflow across our sightline. In each case, the measured variability characteristics provide important constraints on the dynamics, geometry, and/or physical state of the AGN absorbers. For motion across the background AGN emission, the variability timescales can constrain the transverse component of the kinematics of the absorber and the absorption-emission geometry. For ionization changes, the observed absorption variability can be used to constrain the gas number density and, combined with photoionization models, determine the distance of the absorber from the central source. These parameters are needed to determine the mechanism driving the mass outflow and the source of the absorption gas, and assess its overall role in the energetics of the AGN.
The bright Seyfert 1 galaxy NGC 3783 has a rich UV and X-ray absorption spectrum. High-resolution observations with the Goddard High Resolution Spectrograph (GHRS) and the Space Telescope Imaging Spectrograph (STIS) aboard HST showed the UV absorption is highly variable. Three distinct kinematic components of absorption appeared independently over yearly timescales: components 1 โ 3 having radial velocities $`v_r`$1350, $``$550, and $``$725 km s<sup>-1</sup> and widths $`FWHM`$ 190, 170, and 280 km s<sup>-1</sup> (Kraemer et al., 2001, hereafter KC01).<sup>1</sup><sup>1</sup>1Tentative detection of a weak, fourth component at $`v_r`$1050 km s<sup>-1</sup> was described in Gabel et al. (2003a); we do not treat this component in the following analysis. These long-term changes were found to be inconsistent with ionization changes, and thus interpreted as a signature of transverse motion of the absorbers by KC01. High-resolution X-ray observations with the Chandra X-ray Observatory (CXO) revealed a spectrum with numerous absorption lines from a large range in ionization states (Kaspi et al., 2001).
We have undertaken an intensive, multiwavelength campaign on NGC 3783 with HST/STIS, FUSE, and CXO to monitor the absorption properties. Earlier papers in this series have presented studies of the mean X-ray (Kaspi et al., 2002, Paper I) and UV (Gabel et al., 2003a, Paper II) absorption spectra, analysis of a decrease in radial velocity detected in UV component 1 (Gabel et al., 2003b, Paper III), and variability and detailed modeling of the X-ray absorption (Netzer et al., 2003, Paper IV). In this paper, we present a study of the variability and physical conditions in the UV absorption based on analysis of the FUSE and STIS spectra. In ยง2, we review the observations and present the UV continuum light curve; in ยง3, we present measurements of the absorption parameters and variability. We analyze metastable C III\* $`\lambda `$1175 absorption detected in component 1 in ยง4 to derive the number density in this absorber. In ยง5, we present detailed modeling of the UV absorption, making use of observed variability in the spectrum, and explore time-dependent solutions in response to the continuum variations. We explore global models of the outflow in NGC 3783 in ยง6, based on the combined results from the UV and X-ray analysis.
## 2. Observations and the UV Continuum Light Curve
### 2.1. FUSE and HST/STIS Echelle Spectra
We present results from a total of 18 medium-resolution STIS echelle spectra (S1 โ S18) and 6 FUSE spectra (F1 โ F6) of the nucleus of NGC 3783, obtained between 2000 February and 2002 January. A detailed description of the observations and data reduction is given in Paper II; here we present a brief overview.
Each STIS observation was obtained using the 0$`\stackrel{}{\mathrm{.}}`$$`\times `$ 0$`\stackrel{}{\mathrm{.}}`$2 aperture and the E140M grating, which spans 1150โ1730 ร
, and consisted of two HST orbits for a total exposure time of $``$ 4.9 ks (except S1 which was 5.4 ks). The STIS spectra were processed with IDL software developed at NASAโs Goddard Space Flight Center for the Instrument Definition Team, which includes a procedure to remove background light from each order using a scattered light model devised by Lindler (1999). Our measurements of the residual fluxes in the cores of saturated interstellar Galactic lines show the scattered light was accurately removed (see Paper II). The extracted STIS spectra are sampled in $``$ 0.017 ร
bins, thereby preserving the full resolution of the STIS/E140M grating ($`FWHM`$ 7 km s<sup>-1</sup>).
The FUSE spectra, covering 905 โ 1187 ร
, were obtained through the 30$`\mathrm{}`$ $`\times `$ 30$`\mathrm{}`$ aperture. Each spectrum was processed with the standard calibration pipeline, CALFUSE. For each observation, the eight individual spectra obtained with FUSE from the combination of four mirror/grating channels and two detectors, were coadded for all exposures. We corrected for nonlinear shifts in wavelength scale between the spectra from different detector segments by cross-correlating over small bandpasses. The absolute wavelength scale was determined by matching the velocities of Galactic lines in the FUSE spectrum with those in the STIS bandpass. Mean residual fluxes measured in the cores of saturated Galactic lines are consistent with zero within the noise (i.e., standard deviation of the fluxes) in the troughs of these lines, indicating accurate background removal for all observations except F5. For this observation, we fit the remaining residual background flux in strong interstellar lines to match the other epochs, and subtracted the fit. The spectra were resampled into $``$ 0.02 โ 0.03 ร
bins to increase the signal-to-noise ratio (S/N) while preserving the full resolution of FUSE, which is nominally $`FWHM`$ 20 km s<sup>-1</sup>.
### 2.2. The UV Continuum Light Curve
To put the continuum variations of our observations in perspective, Figure 1a shows the continuum flux light curve at 1470 ร
for all UV spectra of NGC 3783 obtained over the last 22 years. These observations were obtained with the International Ultraviolet Explorer (IUE) and the Faint Object Spectrograph (FOS), GHRS, and STIS (in low dispersion) on HST, which are listed in Table 1. We obtained the most recently processed version of each spectrum from the Multimission Archive at the Space Telescope Science Institute. We measured the continuum fluxes in each spectrum by averaging the elements in a 30 ร
bin centered at 1470 ร
in the observed frame, which is free of contamination from line emission. The 1$`\sigma `$ flux uncertainties were determined from the standard deviations. For the IUE spectra, this technique is known to overestimate the errors (Clavel et al., 1991); therefore, following the procedure described in Kraemer et al. (2002), we scaled these uncertainties by a factor of 0.5 to ensure that observations taken on the same day agreed to within the errors on average. Due to the small wavelength coverage of the GHRS spectra, no continuum regions free of broad line emission were observed. To estimate the GHRS continuum fluxes, we used separate fits to the continua and broad emission lines from the STIS spectra and tested different linear combinations of these fits until we obtained an accurate match to the observed GHRS profiles. Figure 1a shows the STIS echelle observations (JD $``$ 2,452,000) sampled the UV continuum in a range of moderately-low to moderately-high flux states compared to the long-term light curve.
In Figure 1b, the light curve for the monitoring campaign observations is shown in more detail. We have included estimated fluxes from the four FUSE observations that were not simultaneous with STIS observations (F1, F3, F5, F6) by extrapolating from the far-UV bandpass to 1470 ร
. The first four observations (S1 โ S4) were obtained at intervals of several months, the intensive monitoring phase (S5 โ S17) sampled the continuum at 3 โ 8 day intervals, and the final observation (S18) was obtained about 9 months later. The extrema in continuum flux in the STIS observations occurred during the long-term sampling; S2 and S3 observed the continuum in the highest flux state and the final observation found it in the lowest state, with a peak amplitude variation of a factor of $``$2.5. We identify two general phases during the intensive monitoring: the first four observations (S5 โ S8) sampled the continuum in a relatively low-state, after which it increased by up to a factor of 1.7 (S10 and S12) over a period of about 2 weeks. Three FUSE observations (F2 โ F4) were obtained during the low-state STIS observations, while observation F5 observed the continuum in the highest state (four days after S12). In the subsequent variability analysis, we compare observations between mean low and high states derived by averaging representative spectra in each state.
## 3. Measurements of Absorption Parameters
A key issue in the analysis of AGN outflows is extracting accurate ionic column densities from the observed absorption lines, since they provide the basis for determining the physical conditions in the absorbers (i.e., the ionization state, total gas column, number density). The UV absorbers typically only partially occult the background AGN emission, as determined by the relative strengths of the individual members of absorption doublets (Wampler et al., 1993; Hamann et al., 1997; Barlow & Sargent, 1997), and the effects of covering factor ($`C`$) and optical depth ($`\tau `$) must be deconvolved to determine the column densities. In some cases, the $`C\tau `$ solution for observed features can be complicated due to blending of physically distinct absorption components or complex coverage of the background emission. Thus, it is useful to first explore absorption variability by comparing equivalent widths.
### 3.1. Variability in Absorption Equivalent Widths
Figures 2a โ c show the equivalent widths for key lines in the STIS spectra of components 1 โ 3, respectively, plotted as a function of the UV continuum flux. The error bars represent our estimated measurement uncertainties, which are a combination of uncertainties due to spectral noise and fitting the intrinsic (i.e., unabsorbed) fluxes, via propagation of errors. We estimated uncertainties in the intrinsic fluxes by testing different empirical fits over the absorption features and selecting the range of what we deemed to be reasonable line profile shapes. We note that these derived uncertainties may in some cases be too conservative, as indicated by comparing the scatter in measured values with the error bars plotted in Figure 2 (e.g., N V in components 1 and 3). This may be due to systematic trends in how the intrinsic flux was fit in each observation and/or overestimates of the uncertainties associated with the fitting.
In each kinematic component, the lowest-ionization species with detectable absorption other than H I (Si IV in component 1; C IV and N V in component 2; C IV in component 3) show clear variations that are inversely proportional to the continuum flux. This is just as expected for unsaturated lines from relatively low-ionization species in a photoionized gas, as their column densities decrease in higher flux states due to the increased ionization. This gives direct evidence that the ionization structure in these absorbers is dominated by photoionization from the central source. It also indicates these lines are not highly saturated, at least in the epochs with weaker observed absorption. To test the dependence of absorption variability on continuum flux more rigorously, we did a linear regression fit for each line. Figure 2 shows the results of the fits and gives the ratio of the computed slope to the 1$`\sigma `$ uncertainty in the slope, $`m/\mathrm{\Delta }m`$, to characterize the correlation. Fits to all of the lines listed above have a non-zero slope at greater than a 3$`\sigma `$ level. Since the absorption strengths will not necessarily vary linearly with continuum flux, we also tested the Spearman rank correlation. All of the above lines show correlation with the flux at high significance, with probabilities of no correlation computed to be between 0.01 โ 0.3 % from this test (values are listed in Figure 2).
Conversely, the absorption strengths of C IV in component 1 and N V in components 1 and 3 did not vary strongly during the monitoring. This could indicate either these lines were saturated, with partial covering factor, or their ionic column densities were not varying. In the latter case, the lack of variability could be due to either a relatively low electron density such that the ionic populations are unable to respond to continuum changes, or to the ionic structure of the gas, with these ions near the peak ionization states of their parent elements. These possibilities are explored below. We note that since each component has some lines that exhibit no variability (see Figure 2 and ยง3.2), the lines that do vary must result from changes in ionic column densities rather than covering factors.
### 3.2. Effective Line-of-Sight Covering Factors and Column Densities
The doublet method for measuring the absorption parameters provides a solution to a single $`C`$ and $`\tau `$ (at each resolution element in an absorption profile). However, since the background AGN emission is comprised of multiple, physically distinct components with different sizes and geometries (i.e., a featureless continuum source and multiple kinematic components of line emission), additional constraints may be needed to determine the required effective covering factors, which are weighted combinations of covering factors of the different emission sources (Ganguly et al. 1999). Implicit in the doublet solution is that all emission sources have the same covering factor. In Paper II, we used the Lyman series lines to separate the continuum and emission line covering factors for the H I absorption. Here, we use variability in the background AGN emission as an additional constraint to explore the effect of the NLR on the derived covering factors and obtain a consistent model to measure absorption column densities for all lines.
#### 3.2.1 Isolating the Narrow-Line Region Emission Profile
We first compare spectra in different flux states to isolate distinct emission-line components based on variability in the overall profiles. Figure 3a shows the C IV line profiles in mean high-state (S2 and S3) and low-state (S5 โ S8) spectra; these observations were selected for comparison because they show the largest difference in emission-line flux. The continuum flux has been subtracted from each spectrum and the low-state profile has been scaled by a factor of 1.4 to match the flux in the high-velocity wings of the high-state spectrum. The profiles match very well at radial velocities $`|v_r|>`$1500 km s<sup>-1</sup> but diverge at lower velocities, with increasing discrepancy toward line center. This effect is consistent with the superposition of a varying broad component and a non-varying narrower emission-line component, hereafter the BLR and NLR respectively.
Assuming no change in the NLR flux between observations, we isolate its profile by solving the following expression for $`F_{NLR}`$ at each radial velocity:
$$F_{low}(v)\times f_{sc}F_{high}(v)=F_{NLR}(v)\times (f_{sc}1),$$
(1)
where $`F_{low}`$ and $`F_{high}`$ are the total observed emission-line fluxes in each state (i.e., the NLR $`+`$ BLR fluxes) and $`f_{sc}`$ is the scale factor equating the low-state BLR flux with the high-state. The resulting NLR profile is plotted as a dotted line in Figure 3a. This analysis assumes the BLR scales by a uniform factor at all radial velocities between states. Equation 1 is not valid in spectral regions affected by absorption. We have fit the NLR profile using Gaussians for each of the C IV doublet lines constrained to have their intrinsic 2:1 flux ratio and velocity separation, giving a width of $`\sigma =`$ 500 $`\pm `$120 km s<sup>-1</sup> ($`FWHM=`$ 1180 $`\pm `$280 km s<sup>-1</sup>) at a radial velocity of $`+`$50 $`\pm `$60 km s<sup>-1</sup> with respect to the systemic velocity. The combined fit, which appears symmetrical due to the broad line width relative to the separation of the doublet members, is plotted over the residual NLR flux in Figure 3a.
We applied the same analysis to derive the NLR profiles of other emission lines using equation 1. Results for Ly$`\alpha `$, N V, and Si IV ($`f_{sc}=`$1.4, 1.8, and 1.8, respectively) are shown in Figure 3b. A scaled template of the Gaussian fit to the C IV NLR profile is overlaid on each residual NLR profile, including both lines for the N V and Si IV doublets. The discrepancy in the residual fluxes in the blue wing of N V ($`v_r<`$2000 km s<sup>-1</sup>) and the red wing of Ly$`\alpha `$ ($`v_r>`$ 3000 km s<sup>-1</sup>) is due to the different BLR flux scale factors for the two lines and the overlap in their BLR emission. The excess emission in the Si IV profile between $``$ 1000 โ 2500 km s<sup>-1</sup> is due to the O IV\] emission-line multiplet. The Ly$`\alpha `$ NLR is not well matched by the C IV template, appearing narrower in the unabsorbed red wing of the profile. Figure 3 shows the derived NLR fluxes contribute substantially at the wavelengths of many of the absorption features and thus its covering factor could affect the measurement of column densities, which we explore below.
There are some caveats regarding the derivation of the NLR profiles. First, the assumption that the BLR scale factor between the low and high states is independent of radial velocity may introduce an error into the solution. Intensive UV โ optical monitoring of continuum and BLR variability in the Seyfert 1 galaxy NGC 5548 has revealed its BLR variability is consistent with a constant virial product (Peterson & Wandel, 1999; Peterson et al., 2004). A consequence of this is that the integrated BLR profile becomes narrower in higher flux states; if this is the case for the UV lines in NGC 3783, then a velocity-dependent scale factor in equation 1 would be required to fully separate the NLR and BLR. Additionally, the profiles derived from this analysis are broader than are typical for UV NLR lines in Seyfert 1 galaxies. For example, HST observations of NGC 5548 and NGC 4151 while in low flux states with little contamination from BLR emission revealed $`FWHM`$ 500 and 300 km s<sup>-1</sup>, respectively, for the C IV NLR. However, the width derived for NGC 3783 is comparable to the rather broad NLR features in the Seyfert 2 galaxy NGC 1068 (Dietrich & Wagner, 1998; Kraemer et al., 1998a). To assess the NLR fit, we compare their flux ratios relative to the measured \[O III\] $`\lambda `$5007 line with those from other AGNs. From optical spectra of NGC 3783 obtained with a 5$`\mathrm{}`$$`\times `$15$`\mathrm{}`$ aperture by Evans (1988), the ratio of C IV from our fit to the total \[O III\] flux is $``$0.5. In comparison, measurements of NGC 5548 yield a C IV : \[O III\] ratio of $``$ 1.2, obtained with 1$`\mathrm{}`$ (C IV) and 4$`\mathrm{}`$$`\times `$10$`\mathrm{}`$ (\[O III\]) apertures (Kraemer et al., 1998b). Based on this, the C IV flux from our derived profile is reasonable. In contrast, in a sample of more luminous radio-loud AGNs, Wills et al. (1993) detected no UV NLR lines, using the observed \[O III\] line as a template, with upper limits on the C IV : \[O III\] flux ratios of $``$ 0.1 โ 0.5.
#### 3.2.2 Covering Factor Model
For absorption features imprinted on multiple discrete background emission sources, the normalized flux for the j<sup>th</sup> line can be written:
$$I_j=\mathrm{\Sigma }_i[R_j^i(C_j^ie^{\tau _j}+1C_j^i)],$$
(2)
where the i<sup>th</sup> individual emission source has fractional contribution to the total intrinsic (i.e., unabsorbed) flux, $`R_j^i=F_j^i/\mathrm{\Sigma }_i[F_j^i]`$, and line-of-sight covering factor, $`C_j^i`$ (Gabel et al., 2005). The effective covering factor associated with each line is the weighted combination of the individual covering factors:
$$C_j=\mathrm{\Sigma }_i[R_j^iC_j^i].$$
(3)
These equations are extensions of the expressions given in Ganguly et al. (1999) for the continuum and BLR to include an arbitrary number of background emission sources. Combining equations 2 and 3 and solving for $`\tau `$ gives the familiar expression for optical depth (Hamann et al., 1997),
$$\tau _j=\mathrm{ln}(\frac{C_j}{I1+C_j}).$$
(4)
The total column densities for each line are then obtained by integrating
$$N_j(v)=\frac{m_ec}{\pi e^2f\lambda }\tau _j(v),$$
(5)
over radial velocity (Savage & Sembach, 1991), where the optical depths are derived in each velocity bin from equations 3 and 4.
Analysis of the Lyman lines in Paper II revealed the emission-lines are only partially occulted by the UV absorbers. Since the NLR is generally much more extended than the BLR and continuum source in AGNs, we assume here a 3-component geometrical model in which the NLR is entirely unocculted by all components of UV absorption. Later, we explore this assumption based on geometrical constraints from our analysis. In Paper II, the emission-line and continuum covering factors were separated assuming a single emission-line region. Here, we incorporate the NLR โ BLR separation for our 3-component coverage model (hereafter 3-$`C`$) by solving for the continuum and BLR covering factor profiles, $`C^c`$ and $`C^{BLR}`$, using the Lyman lines as in Paper II after first subtracting a model of the Ly$`\alpha `$ NLR. Given the mismatch between the Ly$`\alpha `$ NLR residual and C IV template (see discussion above and Figure 3b), we fit the Ly$`\alpha `$ NLR with a narrower Gaussian that matched both the uncontaminated red wing and core, and was constrained to be below the residual flux in the component 2 absorption feature ($`\sigma =`$350 km s<sup>-1</sup>). Given the uncertainty in the NLR profile, we cannot rule out that the deep component 2 feature partially occults the NLR of Ly$`\alpha `$.
Using equation 3, we calculated effective covering factors for each absorption line, using the derived $`C^c`$ and $`C^{BLR}`$ profiles, and $`C^{NLR}=`$0 at all radial velocities. Figure 4 shows the resulting normalized unocculted flux levels (1 $`C`$) for this 3-$`C`$ model compared to the observed normalized absorption profiles. Regions where the residual fluxes in the absorption troughs are close to the unocculted flux levels ($`I1C`$) indicate the lines are near saturation for this model. For each line, results for both low and high-state spectra are shown to demonstrate the variability and implications of the covering factor on observed line strengths. To increase the S/N in each spectrum, we averaged multiple observations for each state, selecting observations that exhibited the largest variations in absorption strengths and continuum flux during the intensive monitoring. For the STIS lines, the low state shown is the mean of S5 โ S8, and the high state S10 and S12, except for Si IV, which had the weakest absorption in observations S15, S16. For lines in the FUSE spectrum, the low-state is the mean of F2 โ F4, and the high-state F5 โ F6. Column densities for the low and high-states in each component computed with equations 4 and 5 are given in Table 2. Below, we describe the absorption lines for each kinematic component.
For the N V absorption, both doublet members are free of contamination over most of the absorption profiles, providing a comparison between the doublet solution to the covering factor ($`C_d`$) and our 3-$`C`$ model to test the effect of an unocculted NLR. The unocculted flux levels for the doublet solution (1$`C_d`$) are shown in Figure 4a with horizontal tick marks for N V. These values were derived in the cores of each kinematic component in the merged spectrum (Paper II). Close comparison reveals important differences in the implied column densities between the two covering factor models. In components 1 and 3, the doublet solution gives unsaturated absorption. This is seen in Figure 4a, where $`I>1C_d`$ in the weaker doublet member (N V $`\lambda `$1242). The resulting N V column densities are $``$ 8 $`\times `$ 10<sup>14</sup> cm<sup>-2</sup> (component 1) and 1.5 $`\times `$ 10<sup>15</sup> cm<sup>-2</sup> (component 3), with equal values in the low and high-state spectra within measurement uncertainties. Thus, the doublet solution implies the N V absorption in components 1 and 3 is unsaturated, with the above column densities, and without any detectable variability between flux states. A similar result is found consistently for all STIS observations. In contrast, the 3-$`C`$ coverage model implies N V is near saturation in these components, in both flux states, since $`I1C`$ consistently for both doublet lines. The primary difference is due to the NLR flux underlying the red doublet members of these components (see Figure 3b). The lack of substantial variability in the equivalent widths (Figure 2a and 2c) provides independent evidence for saturation of N V in these components. N V in component 2 is unsaturated for both covering factor models, consistent with the variability detected in its equivalent width.
#### 3.2.3 Component 1
Component 1, at $`v_r=`$1350 km s<sup>-1</sup>, exhibits a rich absorption spectrum. In addition to the common C IV, N V, and O VI doublets, it includes lines from the relatively low-ionization species Si IV and C II, P V $`\lambda `$1118 (the P V $`\lambda `$1128 doublet line is contaminated with Galactic absorption and unmeasurable) and the metastable C III\* $`\lambda `$1175 complex (Paper II). None of these lines is detectable in the other kinematic components in NGC 3783, and they are only seen in a small fraction of intrinsic absorbers in AGNs (NGC 4151 is the only other Seyfert with reported C III\* and P V absorption, Bromage et al. 1985, Espey et al. 1998; Si IV appears in $``$ 40% of Seyfert absorbers, Crenshaw et al. 1999). The C III\* $`\lambda `$1175 absorption gives a tight constraint on the number density of the absorber, as shown in ยง4. H I is detected in the Lyman series up to Ly$`\theta `$, while contamination from Galactic absorption prevents detection of the higher order lines; thus the H I column density listed in Table 2 should be considered a lower limit.<sup>2</sup><sup>2</sup>2In Paper II, we missed detection of Ly$`\theta `$ (as well as Ly$`\zeta `$) in component 1; thus the value quoted there, which was measured from the Ly$`ฯต`$ line, is smaller than the present study.
There is evidence that multiple, physically distinct absorption regions are overlapping in kinematic component 1 (KC01; Paper II). In Paper II, it was shown the Si IV covering factor from the doublet solution is lower than that derived from the Lyman lines and the N V doublet. This is confirmed for the new 3-$`C`$ model: the column density measured for the red Si IV doublet line with this covering factor is three times greater than the blue line, indicating the actual Si IV covering factor is significantly lower. In contrast, the 3-$`C`$ model gives consistent results for the two N V doublet members: saturation in the blue-wing and core in both flux states (see above discussion), and similar column densities in the red-wing of the profile where the absorption depths diverge from the unocculted flux levels (i.e., the absorption is unsaturated). Similarly, C IV $`\lambda `$1548 and O VI $`\lambda `$1038 are consistent with being saturated in their blue wings and cores with the 3-$`C`$ model (the other doublet members for these ions are unmeasurable due to contamination with other absorption), and C IV diverges substantially in the red wing (O VI $`\lambda `$1038 is contaminated with a detector artifact at these velocities and cannot be tested). Additionally, there are ion-dependent structural differences in the absorption profiles. As described in Paper II, the red wings of some lines, particularly N V and Ly$`\alpha `$, extend to lower velocities than Si IV and P V. Interestingly, it is at these velocities that the N V and C IV profiles diverge from the unocculted flux levels from the covering factor model.
Thus, we treat the absorption from this kinematic region as coming from two physical components: component 1a has relatively low covering factor and low-ionization and gives Si IV, C II, C III\*, and P V; component 1b is more highly ionized and has higher covering factor, contributing strongly to C IV, N V, O VI and the Lyman lines. The latter lines will have contributions from both components, with absorption from component 1a buried in the higher covering factor absorber (e.g., see Kraemer et al., 2003). To measure the Si IV column density, we used the covering factor derived from the doublet pair ($`C=`$0.35; Paper II). For the other lines associated with component 1a, there is no independent measure of the covering factors (it is not possible to separate the emission-line and continuum covering factors for this component from the single Si IV doublet). We assumed $`C=`$1 for these lines because their underlying emission is predominantly continuum flux. Since they are all weak (see Figure 4), the measurements will not be too far off unless $`C^c`$ is very small in this component. For the O VI, N V, and C IV lines, we adopted the 3-$`C`$ model. Since all these lines show saturation in this component over much of the profiles, we derived lower limits on their integrated column densities by adding estimated uncertainties to the normalized fluxes, $`I+\mathrm{\Delta }I`$, in equation 4. Measured column densities and limits for the low and high-states are listed in Table 2. For modeling purposes (ยง5.2), we also measured the C IV and N V absorption in the red wing of the profile, where the lines are unsaturated and there is no contribution from component 1a based on the Si IV profile ($`v_r=`$1270 โ $``$1170 km s<sup>-1</sup>). This gives 3.7$`\times `$10<sup>14</sup> and 1.0$`\times `$10<sup>14</sup> cm<sup>-2</sup> for the N V and C IV column densities, respectively.
As described in ยง3.1, the Si IV absorption varied, with a general inverse correlation with continuum flux. With the relatively low Si IV covering factor, the moderate equivalent width variations translate into peak changes of a factor of $``$ 4 between the mean low and high states during the intensive phase defined above. The column density decreased somewhat less in the high-state S10 and S12 observations relative to the mean low-state (a factor $``$ 2). Figure 4c shows the same trend for C II: weak absorption is detectable in the ground state (C II $`\lambda `$1335) and fine-structure (C II\* $`\lambda `$1336) lines in the low-state spectrum, but was not distinguishable from noise in the high-state. The weak P V $`\lambda `$1118 and C III\* $`\lambda `$1175 absorption appears to decrease as well, with only $``$1.5$`\sigma `$ detections in the high-state, though the limited S/N in the spectra of these lines makes this less certain. No significant variability was detected in Ly$`\theta `$, nor any other Lyman series lines, between flux states.
#### 3.2.4 Component 2
In component 2 ($`v_r=`$550 km s<sup>-1</sup>), the N V and C IV absorption is relatively weak and variable (Figures 2b and 4a). Our adopted covering factor model gives unsaturated absorption in these doublets, consistent with the observed variability. In contrast, Figure 4b shows the O VI $`\lambda `$1038 line is strong and matches the model unocculted flux levels nearly identically across the entire profile in both states, consistent with heavy saturation of this line. Lyman line absorption is measurable only up to Ly$`\gamma `$ in component 2 due to strong contamination of the higher order lines with Galactic absorption. The Ly$`\gamma `$ line does not vary significantly, thus we adopt the H I column density measured from this line as a lower limit. Figure 4b shows both the C III $`\lambda `$977 and N III\* $`\lambda `$991 lines are contaminated with moderate Galactic H<sub>2</sub> absorption at component 2 velocities, with neither line showing strong evidence for the presence of intrinsic absorption. Upper limits on these lines measured after the removal of the H<sub>2</sub> model are given in Table 2.
#### 3.2.5 Component 3
In component 3 ($`v_r=`$725 km s<sup>-1</sup>), the primary constraints for modeling are the C III $`\lambda `$977 line and the C IV doublet. Both lines exhibit variability, decreasing in epochs with higher continuum flux. The C III line is relatively strong in the low-state FUSE spectrum, but is similar to the noise level in the high state (Figure 4b), giving only an upper limit. Figure 4a shows that for the adopted covering factor model, the C IV $`\lambda `$1551 line is near saturation in the low-state, but well above the unocculted flux level in the high-state. Thus, for this model, the equivalent widths in C IV vary only moderately (Figure 2c), but the variation in column density is large between flux states due to the low effective covering factor ($`C`$ 0.35 in the core of component 3). The C IV $`\lambda `$1548 line is blended with the component 1 C IV $`\lambda `$1551 line, thus no comparison of our adopted model with the doublet solution is possible. The O VI $`\lambda `$1038 and N V doublet lines (see above) are consistent with being saturated in both flux states, based on the unocculted flux levels. The highest order Lyman line not contaminated with other absorption in component 3 is Ly$`ฯต`$. Since it does not vary significantly between flux states, we adopt the resulting H I column density from this line as a lower limit. The N III $`\lambda `$989 component 3 line coincides with the strong N III\* $`\lambda `$991 component 1 line, and cannot be measured. N III\* $`\lambda `$991 in component 3 is moderately contaminated with Galactic H<sub>2</sub>. Given the limited S/N in this region, we take the measurement of N III\* after the removal of the H<sub>2</sub> model as an upper limit.
## 4. Constraints on the Density from the Metastable C III Absorption
The C III\* $`\lambda `$1175 multiplet lines have been used as a density diagnostic for AGN absorbers in several studies. However, as pointed out by Behar et al. (2003), the high densities derived in these studies, including our Paper II, were based on calculations of level populations that only treated the <sup>3</sup>P<sub>1</sub> level. The $`J=`$0 and 2 levels have much lower radiative transition probabilities to the ground state ($`A`$0 and 5$`\times `$10<sup>-3</sup> s<sup>-1</sup>, respectively) than the $`J=`$1 level ($`A=`$75 s<sup>-1</sup>), and thus are populated at densities that are lower by several orders of magnitude (see Figure 1 in Bhatia & Kastner, 1993; Kastner & Bhatia, 1992). This has important consequences for the interpretation of the absorbers since the gas density is needed to determine the location and physical depth of the gas. To address this, we computed the relative populations of the <sup>3</sup>P<sub>J</sub> levels for a range of densities and electron temperatures, $`T_e`$, that are expected for the photoionized AGN outflows seen in the UV. Our calculations extend the results of Bhatia & Kastner (1993), who give results applicable to collisionally ionized plasmas with $`T_e`$ 40000 K, to lower temperatures. Collisional excitation and de-excitation and radiative decay between all levels were treated. We treated the 6 lowest terms/levels of the $`C^{+2}`$ ion: the ground state 2s<sup>2</sup> <sup>1</sup>S<sub>0</sub> term, the three 2s2p <sup>3</sup>P levels, 2s2p <sup>1</sup>P<sub>1</sub>, and 2p<sup>2</sup> <sup>3</sup>P. Temperature dependent collision strengths from Berrington et al. (1985, 1989) were used for transitions between levels and terms, respectively, where we have interpolated between their listed temperatures. Radiative transition rates were obtained from Bhatia & Kastner (1993), Morton (1991), and the NIST Atomic Spectra Database website.
Figure 5 (top panel) shows the computed populations over a large range in density for $`T_e=`$ 16000, 20000, and 40000 K. This shows that the C III\* $`\lambda `$1175 absorption complex can serve as a powerful probe of the physical conditions in the absorber. The relative populations of the three levels are very sensitive to the electron density, but insensitive to temperature. This is seen in the bottom panel of Figure 5, where the ratio of the $`J=`$2 : $`J=`$0 level populations are plotted for the temperatures shown in the top panel. Additionally, the absolute populations of the <sup>3</sup>P levels are very sensitive to the gas temperature.
In practice, full utilization of these lines as diagnostics requires that the absorption features are sufficiently narrow so that the individual lines in the complex are not heavily blended, and sufficient resolving power to separate the lines. These conditions are met with the STIS spectra of component 1 in NGC 3783 as seen in Figure 6, which shows only mild blending of the individual lines. The location of the six multiplet lines are marked and identified by the $`J`$ level of the transition. A fit to the C III\* $`\lambda `$1175 complex is shown as a dashed line, using the width and centroid of the Si IV absorption profile; best-fit column densities for each level are given below the spectrum. The column density ratio of the $`J=`$2 : $`J=`$0 levels, $`N_{J=2}/N_{J=0}=`$ 2.9$`\pm `$1.4, gives $`n_e`$=3$`{}_{1.5}{}^{}{}_{}{}^{+5}`$$`\times `$10<sup>4</sup> cm<sup>-3</sup>, independent of temperature. The lack of detection of the $`\lambda `$1174.93 line is consistent with this density since Figure 5 shows the population of the $`J=1`$ level will be approximately four orders of magnitude lower than the other levels. Use of these level populations as a temperature diagnostic requires knowing the total abundance of the C<sup>+2</sup> ion; since the C III $`\lambda `$977 line is unmeasurable due to contamination (and saturated for the implied column density), this requires the results of photoionization modeling, which is presented in ยง5.
Our calculations do not treat the effects of the radiation field on the level/term populations. This could decrease the ground state or metastable level populations due to continuum pumping, or increase the relative populations of the metastable levels due to recombination followed by cascade to the 2s2p levels. We have computed photoionization models with Cloudy to test this and find it has a negligible effect on the metastable level populations in this component.
## 5. Photoionization Models
### 5.1. Input Parameters and Assumptions
We compare the observed ionic column densities with predictions from the photoionization modeling code Cloudy (Ferland et al., 1998) to constrain the physical conditions and location of the intrinsic absorbers. The absorbers are assumed to be uniform plane parallel slabs of constant density, photoionized by the AGN at a distance $`R`$ from the central source. These calculations apply for a gas in ionization and thermal equilibrium with the ionizing flux; this is explored in ยง5.3 based on the observed variability in the absorption. The models are specified by the spectral energy distribution (SED) of the ionizing continuum, the total hydrogen column density ($`N_H`$) and elemental abundances in the absorber, and the ionization parameter ($`U=Q/4\pi R^2n_Hc`$), which gives the ratio of the density of H ionizing photons at the face of the absorber to the gas density, $`n_H`$.
The choice of SED for the unobservable ionizing continuum is a source of modeling uncertainty in AGNs (e.g. Mathur et al., 1994; Kaspi et al., 2001). To facilitate the comparison with the earlier UV results, we adopted an SED based on the one used by KC01. This consists of multiple power-law components ($`F_\nu \nu ^\alpha `$) constrained by the UV and X-ray observations. The UV power-law component is based on continuum-flux measurements in the two epochs with simultaneous STIS and FUSE observations. After correcting the spectrum for Galactic extinction ($`E(BV)=`$ 0.119; Reichert et al., 1994) with the Cardelli et al. (1989) reddening curve, the continuum flux ratio measured at 950 ร
and 1690 ร
gives a spectral index of $`\alpha _\nu ^{UV}`$ 1. As discussed in KC01, extrapolating the UV power-law to X-ray energies far overestimates the observed flux, thereby requiring a spectral break in the unobservable EUV โ soft X-ray flux to connect the absorption in the two bands. Thus, we extrapolated from the observed UV flux at the Lyman limit to the observed flux at 0.6 keV in the Chandra spectrum, giving a spectral index $`\alpha _\nu ^{EUV}=`$ 1.4. For the hard X-ray index, we used the value derived from the merged Chandra spectrum in Paper I, $`\alpha _\nu ^X`$ 0.7 ($`\mathrm{\Gamma }=`$1.7). For this SED and the adopted distance of 39 Mpc for NGC 3783 (using $`z=`$0.00976 from de Vaucouleurs et al., 1991, and $`H_o=`$75 km s<sup>-1</sup>), this gives an H-ionizing photon luminosity $`Q=`$1.8$`\times `$10<sup>54</sup> s<sup>-1</sup>. The SED adopted here is somewhat different from the one used to model the X-ray absorption in Paper IV, which had $`\alpha =`$0.5, 4.7, and 0.77 in the energy ranges 0.002 โ 0.04, 0.04 โ 0.1, and 0.1 โ 50 keV). We recomputed our models derived in the next section with the SED used in Paper IV and find no significant differences in the predicted column densities for the ions measured in the UV spectrum, after scaling the ionization parameter to account for the somewhat larger relative EUV flux used in the X-ray models. The issue of SED is addressed further below in the discussion of the combined UV and X-ray modeling results (ยง6.2). We adopted the roughly solar elemental abundances (Grevesse & Anders, 1989) used in the KC01 models, and assumed no dust is present in the absorber.
The combined UV and X-ray spectrum of NGC 3783 shows that a large range of ionization states is present in all kinematic components and there are multiple zones of ionization overlapping at all absorption velocities (KC01; Paper IV). Thus some lines, particularly from more highly ionized species, may be comprised of blends of different physical components having different covering factors and optical depths; evidence for this was presented in ยง3 for kinematic component 1, for example. Therefore, we take as primary constraints the lines least likely to be affected by blending or saturation. These include lines from the lowest-ionization species detected in each kinematic component, ions with low abundances of the parent element (e.g., P V), and ions in excited states (e.g. C III\*). Additionally, upper limits on column densities from non-detections provide unambiguous constraints.
### 5.2. Model Results for Mean Low and High State Spectra
Figure 7a shows the solutions in log($`U`$) and log($`N_H`$) from a grid of photoionization models that match the measured column densities for component 1. The contours on the left are for the low covering factor absorber, component 1a. The thickness of the contour for each ion spans the range of solutions corresponding to estimated uncertainties in the measured column densities (e.g., see Arav et al., 2001). Figure 7a includes solutions for Si IV, C II, C III, and P V measured in the low-state spectrum, plotted with hatched marks. The C III solutions are for the measured metastable level column densities, using the electron density derived from the C III\* $`\lambda `$1175 feature in ยง4 and an electron temperature of $`T_e`$1.8$`\times `$10<sup>4</sup> K, consistent with the best-fit Cloudy model. The solutions to the high-state column densities (mean of S15 and S16) for Si IV are also included on the plot, shown as the contour with no hatch marks. These were shifted along the x-axis by log($`U`$)$`=`$0.23, corresponding to the peak amplitude of flux variation observed during the intensive monitoring. This gives the high-state solutions in terms of the low-state ionization parameters, allowing a direct comparison of the two flux states on the same grid: for a correct model with both states in ionization equilibrium and an ionization parameter that scales as the observed UV continuum flux, the low-state and shifted high-state solution contours would overlap. High-state solutions occupying regions in $`U`$ \- $`N_H`$ above and to the left of the low-state solutions overestimate the variability observed between states, while solutions to the right and below underestimate the variability.
Figure 7a shows all ions plotted for the low-state spectrum are fit well by an extended, narrow range of models in $`U`$$`N_H`$ parameter space, with lower bounds $`log(U)>`$ $``$1.7 and $`log(N_H)>`$ 20.3. For higher ionization models, these ionic abundances (particularly Si IV and P V) are very sensitive to small changes in $`U`$ and $`N_H`$, as seen in the narrowing of the solutions spanning the measured limits on the column densities. This is due to the He II opacity โ these solutions are in the region in parameter space where the He II edge becomes optically thick. In these models, the relatively low-ionization species in component 1a exist primarily in a small region in the back end of the slab, where the continuum flux is heavily filtered, thus their abundances are greatly affected by the sensitivity of the He II opacity in the Stromgren shell to the model parameters (e.g. Kraemer et al., 2002). For these high-ionization models, the range of solutions becomes linear in $`U`$ \- $`N_H`$, tracing the He II column density contour.
Combining the shifted solution to the high-state Si IV column density, the range of solutions is limited to low-ionization values. The overlap between Si IV high and low-states gives slightly lower $`U`$ than the fit to all low-state lines, but the solutions are close. For models with $`log(U)>`$1.5 in the low-state, Figure 7a shows the Si IV variability is predicted to be significantly greater than observed. This discrepancy becomes pronounced for higher-ionization solutions: e.g., a factor of 1.7 increase in ionizing flux from the log$`U=`$1, log$`N_H=`$21.6 model results in a Si IV column density of $``$ 10<sup>12</sup> cm<sup>-2</sup>, which is 40 times weaker than observed. Thus, based on these models of the mean high and low-state spectra, the lowest ionization absorber in component 1 has best-fit solution log($`U`$$`=`$1.6$`{}_{0.2}{}^{}{}_{}{}^{+0.2}`$, log($`N_H`$$`=`$ 20.6$`{}_{0.3}{}^{}{}_{}{}^{+0.4}`$. The limits on the high-state C II, C III\*, and P V column densities (not shown on Figure 7a) are also compatible with these solutions. We note the predicted column densities for O VI, N V, and C IV in this absorber would produce heavily saturated absorption in their resonance doublet lines, with $`N_{ion}>`$ 5$`\times `$10<sup>15</sup> cm<sup>-2</sup> for each ion. Model predictions for all ions are listed below their measured values in Table 2.
The physical parameters of the absorber with higher covering factor, component 1b, can also be constrained. For optically thin conditions, the C IV and N V column densities each scale linearly with $`N_H`$, but have different dependences on $`U`$; thus, the C IV : N V ratio uniquely determines the ionization parameter, independent of the total column density. As a result, assuming the ionization structure in component 1b is uniform over all radial velocities, the C IV and N V column densities measured in the unsaturated red wing determines $`U`$ for the entire absorber. Our measurements give log($`U`$)$`=`$0.4$`{}_{0.2}{}^{}{}_{}{}^{+0.3}`$. A lower limit on the total column density comes from the lower limits measured on C IV and N V over the full, saturated profile. For the above $`U`$, this gives log($`N_H`$$``$ 20.3. The O VII column density measured in the CXO spectrum over the velocity range coinciding with component 1 (see Figure 10 in Paper I) provides an upper limit on $`N_H`$. Our fit to this high-velocity region of the profile gives $`N_{OVII}=`$ 5$`\times `$10<sup>16</sup> cm<sup>-2</sup>, with an upper limit of $``$ 10<sup>18</sup> cm<sup>-2</sup> at the 90% confidence level. Incorporating this as an upper limit on O VII gives $`N_H`$ 21.4 for component 1b. The region in $`U`$, $`N_H`$ parameter space spanned by these limits is also plotted on Figure 7a.
The modeling constraints for component 2 are C IV and N V in both the low and high states, and the upper and lower limits on C III and H I, respectively. These solutions are shown on the $`UN_H`$ plot in Figure 7b; N V and C IV are shown as hatched contours in the low-state and as unfilled dashed contours in the high-state, which are shifted to account for the flux difference between states as in Figure 7a. In each state individually, the C IV and N V solutions overlap over an extended range in $`UN_H`$; however, the combined high and low-states are simultaneously fit by only a small region of parameter space selecting the lower ionization solutions; log($`U`$) $`=`$0.45$`{}_{0.1}{}^{}{}_{}{}^{+0.2}`$ (low-state values), log($`N_H`$) $`=`$ 20.4$`{}_{0.1}{}^{}{}_{}{}^{+0.6}`$. These solutions are also consistent with the H I and C III limits.
In component 3, C III $`\lambda `$977 is clearly detectable in the low-state FUSE spectrum, representing the lowest-ionization species detectable in this component, but has only an upper limit in the high state (see ยง3.2.5 and Figure 4b). For our adopted effective covering factor, C IV is near saturation in the low-state, but well below the residual flux level in the high-state giving a correspondingly small column density. The N V doublet is consistent with being saturated in both states. The model results are shown in Figure 7c, with high-state solutions shifted as in Figures 7a and 7b. This shows the C IV high-state and C III low-state column densities (hatched contours) are simultaneously matched by a restricted space of correlated $`U`$ and $`N_H`$ values, with log($`U`$)$``$1.2, log($`N_H`$)$``$19. Including lower limits on N V and H I places lower limits of log($`U`$$``$0.7 and $``$0.6, respectively. An upper bound on the solution comes if we require the model X-ray columns do not exceed those measured in the CXO spectrum. The O VII column density integrated over all radial velocities ($`N_{OVII}=`$ 10<sup>18</sup> cm<sup>-2</sup> from Paper I), provides the most rigid constraint and is plotted in Figure 7c. Incorporating this limit and assuming all the lines measured in the UV arise in a single physical component gives log($`U`$)$`=`$0.5$`{}_{0.1}{}^{}{}_{}{}^{+0.1}`$, log($`N_H`$)$`=`$ 21.1$`{}_{0.2}{}^{}{}_{}{}^{+0.3}`$.
These modeling results can be compared to those from KC01, which were based on the first STIS observation. The ionization and total column density for component 1a are substantially greater here than in KC01. This is due to the constraints implied by the detection of C II and P V in the data presented in this study, and the use of a lower covering factor for Si IV, giving a larger Si IV column density. A reliable covering factor for Si IV was made possible by the high S/N in the merged STIS spectrum (Paper II). The solutions to components 2 and 3 presented above have similar total column densities as the KC01 solutions, but have somewhat lower ionization parameter ($``$ 0.3 โ 0.4 dex). For component 2, this lower ionization was imposed on the solution primarily by the magnitude of variability observed in C IV and N V between high and low flux states. For component 3, the detection of C III $`\lambda `$977 in the FUSE spectrum and the limit placed by O VII measured in the CXO spectrum drove the solution to lower ionization.
### 5.3. Time-Dependent Ionization Solutions
Here we explore variability timescales for the UV absorbers. With the density constraint derived for component 1a from C III\*, the time-dependent ionic populations for this absorber can be probed in detail for comparison with the observed variability, thereby providing better constraints on its physical conditions. For the other absorption components, the observed magnitude and timescales for variability in ionic populations in response to the continuum flux variations can be used to place limits on the density.
#### 5.3.1 Detailed Variability Calculations for Component 1
The response time for an ion is a strong function of several factors: its population ($`n_i`$) relative to adjacent ionization stages, the electron density, and the magnitude of change in ionizing flux incident on the absorber(e.g. Hamann et al., 1997). We computed time-dependent ionic populations in response to changes in the ionizing continuum for the gas in component 1 based on the density and ionization solutions determined above. The problem involves solving for ionic abundances, $`n_i`$, from the system of first-order differential equations:
$$\frac{dn_i}{dt}=[p_i+n_e\alpha _{rec,i1}(T)]n_i+n_en_{i+1}\alpha _{rec,i}(T)+p_{i1}n_{i1},$$
(6)
where $`p_i`$ and $`n_e\alpha _{rec,i}`$ are the ionization rates from and recombination rates to ionization stage $`i`$ (e.g., Krolik & Kriss 1995). Output from the low and high-state Cloudy equilibrium models were used for the parameters in equation 6, corresponding to initial and final states associated with a change in ionizing flux. The time-dependent ionic abundances were then solved using the Runge-Kutta-Fehlberg method, which compares fourth and fifth-order Runge-Kutta estimates to adjust the step size. The rates include all ionization and recombination processes treated by Cloudy (see Ferland et al., 1998). We consider simple step-function increases and decreases in the ionizing flux below.
For example, to compute the time-dependent component 1 Si IV column density in response to an increase in ionization from the low to high states, the initial values of all ionic species of silicon ($`n_i(t=0)`$) were set to the values from the low-state equilibrium model (log($`U`$)$`=`$1.6, log($`N_H`$)$`=`$20.6). The ionization rates, $`p_i`$, were set to the values from the corresponding high-state equilibrium model (log($`U`$)$`=`$1.37). This assumes the ionization rates change instantaneously, which is valid if they scale linearly with the ionizing flux, as is the case in the model considered here. The recombination rates differ somewhat between the initial and final state equilibrium models because of their temperature dependence. However, for the time intervals and magnitude of flux variations considered here, the absorber is not in thermal equilibrium, as determined by the thermal timescale:
$$t_{thermal}=\frac{3/2n_ekT}{n_e^2\mathrm{\Lambda }(T)},$$
(7)
where the numerator is the total thermal energy and the denominator the net cooling rate in the gas per unit volume, with $`\mathrm{\Lambda }(T)`$ representing the difference between total heating and cooling rates, which are obtained from the Cloudy calculations. Even for extreme changes in ionizing flux, the thermal timescale (taken here as the e-folding time) is of order 1/2 year for these conditions. It is much longer than this for the relatively minor flux changes observed in NGC 3783, due to the small value of the net cooling rate. Thus, the gas is not in thermal equilibrium and the actual value of the temperature reflects an average of the long-term history of the cooling and heating rates as they respond to variations in the ionizing flux. For the conditions considered here, the temperature can be assumed to be constant. Thus, we adopt the recombination rates from the initial state models in our solutions to equations 6.
Figure 8 illustrates the dependence of time-dependent populations on the amplitude of ionizing flux variations. The time-dependent behavior of the Si IV column density is given for both increases and decreases in ionizing flux by factors 1.4, 1.7, and 3, showing the time needed to achieve a given level of variability is sensitive to the amplitude of flux change. For models with increased ionization, we used the low-state model solution as the initial state, and for decreased ionizing flux, the high-state model. For large decreases in ionizing flux, the recombination timescale (i.e., e-folding time) is seen to approach the expression from Krolik & Kriss (1995), $`t_{recomb}=n_i/(n_{i+1}n_e\alpha _{rec,i})`$; $`t_{recomb}=`$ 5 days for Si IV for the model considered here. Figure 8 also shows the time interval corresponding to a given fractional change in the ionic abundance is similar for cases of ionization and recombination for the moderate flux variations considered here. However, for very large flux changes, our calculations show the ionization timescale becomes much shorter than the recombination timescale.
We calculated the detailed variability for the component 1a absorber based on the observed flux variations in individual STIS observations. In Figure 9, calculations for multiple step-function variations in flux, corresponding to the continuum light curve during our intensive monitoring, are compared with the measured Si IV column densities. The solid line shows results for the model derived in ยง5.2 based on the low and high-state solutions; the UV light curve is shown in the bottom panel for comparison. The rates and initial ionic populations were taken from the low-state Cloudy model, log($`U`$)=$``$1.6, log$`N_H`$=20.6. At each epoch, the ionization rates were then scaled by the change in flux observed in the subsequent observation, and the time-dependent populations computed using equation 6. The overall Si IV variability is seen to be matched well by this model. It reproduces the somewhat damped decrease in Si IV column density observed in high-states S10, S12 following the increased flux in those observations, and the further decrease in later epochs, where the absorption is a minimum at S15.
In ยง5.2, an extended range of solutions, with correlated $`U`$ and $`N_H`$ values, was found to match the low-state column densities in component 1a reasonably well, with the high-state solution selecting the lower $`U`$ models (Figure 7a). If the physical conditions in this absorber are such that the Si IV population is not able to respond to the ionizing continuum sufficiently, the gas could in principle be in a higher ionization state than modeled by assuming full equilibration between states. With the density for component 1a determined independently from C III\*, it is possible to test this directly. Thus, we computed the time-dependent Si IV population for higher $`U`$ models to test its variability. Figure 9 shows results for a model with log($`U`$$`=`$1, log($`N_H`$$`=`$21.5 (dashed line), which was selected from the low-state solutions in Figure 7a. This high $`U`$ model is seen to far overestimate the variability in Si IV, and thus is excluded by our timing analysis.
#### 5.3.2 Constraints on Density in Components 2 and 3 from Observed Variability
Components 2 and 3 do not have direct estimates of the density, and thus are less well constrained than component 1. However, lower limits on their densities can be derived based on the observed variability. This constraint comes from the characteristic time separation between the low and high-state epochs, which provides an upper limit on the timescale for changes in ionic abundances as they respond to the ionizing continuum.
For component 2, we computed the time-dependent C IV and N V abundances in the same manner as above for component 1 to derive these limits. We assumed a simple step-function change in flux of a factor 1.7, corresponding to the continuum variation between the mean low and high-states during the intensive monitoring. We parameterized the observed variability for these models by requiring the ionic populations varied by at least half the amount observed over a characteristic timescale, which we define as the time interval between the mean of the low states (S5 โ S8) and mean of the high states (S10, S12), $``$ 20 days, where maximum change in these column densities occurred. This places a lower limit on $`n_e`$. The e-folding variability times are proportional to the inverse of the density ($`tn_e^1`$). Thus, we computed the model for a fiducial density and then solved for the value of $`n_e`$ that reproduces the required level of variability over the characteristic timescale.
To demonstrate this, Figure 10 shows the time-dependent N V and C IV ionic abundances for component 2 based on the solution in ยง5.2 (log($`U`$)=$``$0.45, log($`N_H`$)=20.4). These were computed from equation 6 for a fiducial density of $`n_e=`$10<sup>4</sup> cm<sup>-3</sup>, thereby giving $`N_{ion}/N_{ion,t=0}`$ as a function of $`t\times (n_e/10^4)`$. Horizontal dashed lines mark the 50% level of variability; dotted lines give the final, equilibrium values. The lower limits on density follow straightforwardly from the intersection of the abundance curves with the variability limits (dashed lines), which is assumed to occur over a 20 day interval: $`n_e=t/20\times 10^4`$ cm<sup>-3</sup>. The N V variability gives $`n_e`$ 1.3$`\times `$10<sup>3</sup> cm<sup>-3</sup> and C IV gives $`n_e`$2.3$`\times `$10<sup>3</sup> cm<sup>-3</sup> for component 2. Component 3 is less well constrained because for each of the line constraints, at least one state provides only a limit on the column density. A similar analysis for this component gives $`n_e`$7.5$`\times `$10<sup>2</sup> cm<sup>-3</sup>.
## 6. Interpretation
### 6.1. Constraints on Physical Conditions and Geometry of the UV Absorbers
We now derive constraints on the physical conditions and geometry of the absorbers based on our above analysis. The best constraints are for the high-velocity outflow region, component 1, due to the density measured directly from the metastable C III\* feature (ยง4). From the expression for the ionization parameter given in ยง5.1, the distance between the absorber and central source can be solved for the low-ionization absorber (component 1a), using the derived values of the number density ($`n_Hn_e`$/1.2), $`Q`$, and $`U`$ from the modeling in ยง5.2. This gives $`R=`$7.7$`{}_{4.6}{}^{+3.1}\times `$10<sup>19</sup> cm (25 pc). If the individual kinematic components are uniform clouds, i.e., having internal volume filling factors of unity, then their radial physical depths can be determined straightforwardly from the ratios of the total column density to the number density, giving $`\mathrm{\Delta }R=`$1.3$`\times `$10<sup>16</sup> cm for component 1a. Results are summarized in Table 3.
In Paper III, we showed the radial velocities for all component 1 lines in the STIS spectra decreased at a rate of $``$ 50 km s<sup>-1</sup> yr<sup>-1</sup>. This includes Si IV, which comes from the low-ionization, low-covering subcomponent, and the N V and C IV lines which are predominately from component 1b. Thus, these absorbers appear to be dynamically linked and, hence, co-located. Using this result, and adopting the ionization parameter for component 1b derived in the unsaturated red-wing of C IV and N V (ยง5.2), the density for this higher ionization component can be solved from: $`n_{H,1b}=n_{H,1a}\times U_{1a}/U_{1b}`$, giving $`n_{H,1b}=`$ 1.9$`\times `$10<sup>3</sup> cm<sup>-3</sup>. A broad range of values for the radial depth for component 1b follow from the loose constraints on $`N_H`$, 8$`\times `$10<sup>16</sup>$`\mathrm{\Delta }R`$1$`\times `$10<sup>18</sup> cm. This can be compared to the lower limit on the projected transverse size of this absorber, determined by the partial coverage of the BLR from the Lyman line analysis and the size of the BLR based on reverberation mapping (Onken & Peterson 2002), giving $`X_T`$ 1$`\times `$10<sup>16</sup> cm (Paper II). The transverse size of component 1a is not well constrained because the details of the covering factors of the individual sources are not known; the only constraint is that the UV continuum source is at least partially covered by this absorber.
For components 2 and 3, the constraints derived on $`n_e`$ from the observed variability (ยง5.3.2) give upper limits on their distances from the AGN and radial depths. Our derived lower-limit on $`n_e`$ in component 2 based on the calculated response times for C IV and N V gives $`R`$ 7.4$`\times `$10<sup>19</sup> cm (24 pc) and $`\mathrm{\Delta }R`$ 10<sup>17</sup> cm. The density limit for component 3 implies $`R`$ 1.4$`\times `$10<sup>20</sup> cm (45 pc) and $`\mathrm{\Delta }R`$1.3$`\times `$10<sup>18</sup> cm. Lower limits on the projected transverse dimensions of these absorbers follow from the BLR covering factors, giving $`X_T`$ 1 โ 2$`\times `$10<sup>16</sup> cm. Noting that $`\mathrm{\Delta }Rn_H^1`$ and $`Rn_H^{1/2}`$, the absorberโs geometries are seen to depend strongly on their locations when combined with the constraints on the transverse sizes. For example, if the absorbers are uniform clouds having roughly spherical geometries, with dimensions approximately equal to the lower limits given by the BLR covering factors, then distances $``$ 10 and 5 pc are implied for components 2 and 3, respectively. Smaller distances would require flattened geometries. For example, if they are located at 1/2 pc, then they would be extremely thin shell structures, having radial dimensions at least 100 โ 500 times smaller than their transverse sizes.
The results summarized in Table 3 show all UV kinematic components have low global volume filling factors ($`\mathrm{\Delta }R/R`$ 1), and are consistent with being relatively small, discrete clumps. Additionally their derived distances (or limits) are consistent with the location of the inner narrow line region in AGNs, and inside the more extended, diffuse NLR. The geometry implied by these results is consistent with our earlier assumption that the NLR is unocculted by the individual components of UV absorption in deriving the covering factor model (ยง3.2.2).
### 6.2. Connection between the UV and X-ray Absorption
Here, we explore the connection between the observed UV and X-ray absorption. The absorption in the two bandpasses shows similar kinematic structure: in Papers I and II, all X-ray lines having sufficiently high-resolution and S/N were found to span the radial velocities of the three UV kinematic components. Modeling of the CXO spectrum in Paper IV showed the X-ray absorption is highly inhomogeneous, requiring three ionization zones that span a range $`>`$ 50 in $`U`$ to reproduce the full set of lines. A summary of the modeling and geometric constraints from that analysis is given in Table 3, together with our results for the UV absorbers, which imply further inhomogeneities in the outflow in NGC 3783. To convert the X-ray component solutions for a consistent comparison with the UV models, we matched models derived with the SED defined in ยง5.1 with the model predictions in Paper IV, which used an SED with stronger relative emission in the EUV (see ยง5.1). The ionization parameters were scaled by matching the peak ions in each ionization component with the predicted columns in Table 3 of Paper IV. The models computed with the two SEDs were found to give very similar results, matching the dominant ions in each component to better than 90% .
Table 3 shows the three UV components with relatively high-ionization (1b, 2, and 3) share the same ionization parameter as the lowest-ionization region modeled in the X-ray (hereafter XLI), although they have a smaller (integrated) total column density. In Table 4, the predicted column densities from the UV absorbers for key ions with lines in XLI are listed. This shows the component 3 and component 1b absorbers may give significant contribution to some of the XLI lines. Indeed, upper bounds on $`U`$, $`N_H`$ in these components are from the measured O VII column density (ยง5.2). Table 4 shows component 3 may also contribute strongly to Mg IX and about a third of the Si IX and Si X measured in the X-ray spectrum. If $`N_H`$ is near the upper limit for the less well constrained component 1b absorber, it may produce strong Si IX โ Si XI, Mg VIII, and Mg IX. O VI is the only ion with lines detected in both bandpasses. The UV doublet is found to be saturated in all three kinematic components (Figure 4b), giving only a lower limit, and thus consistent with the measurement from the CXO spectrum, $`N_{OVI}`$10<sup>17</sup> cm<sup>-2</sup> (Paper IV).
The XLI model predicts a C IV column density that would produce heavily saturated absorption in the UV doublet in all flux states observed during our monitoring. However, as shown above, the C IV equivalent width varied in components 2 and 3, implying these features are not saturated, at least in the high flux states when their absorption was weaker. Although C IV is likely saturated in component 1, the dominant X-ray absorption coincides kinematically with the lower velocity UV components (Paper I). Our model for component 3, constrained by the unsaturated high-state C IV, fails to reproduce the lowest ionization X-ray species measured in Paper IV, underestimating Si VII, Si VIII, and O V by factors of 10 or more. Component 2 contributes even less due to its relatively small total column density. We have done extensive modeling to test if this discrepancy between the measured UV and X-ray column densities can be reconciled for any values of the model parameters. We find no reasonable choice for an SED that reproduces the measured Si VII and O V column densities while maintaining unsaturated absorption in the C IV UV doublet, due to the similarity in ionization potentials of these ions, nor does any combination of multiple components with different physical conditions. One possible explanation for this discrepancy is that it is due to a geometrical effect. For example, the very large column density of C IV implied by the XLI model may be buried in the variable, unsaturated absorption in components 2 and/or 3. This would require a low covering factor of the UV emission by XLI given the relatively shallow absorption observed in C IV components 2 and 3 (see Figure 4a). Since the X-ray emission region is much more compact than the UV BLR (and likely the UV continuum source as well), it is reasonable that the same gas would have different line-of-sight covering factors in the two bandpasses. Qualitatively, this is consistent with the inhomogeneous wind model described in detail in the next section, in which denser, lower-ionization regions occupy smaller and smaller volumes in the global outflow. The lowest ionization X-ray species could conceivably come from a region that is sufficiently dense and compact that its presence is not seen against a larger absorber giving the unsaturated C IV. Alternatively, it may imply the XLI absorber does not occult the UV absorber at all and that we are seeing physically distinct absorbers in different lines-of-sight in the two bandpasses, thus placing very specific requirements on the absorption โ emission geometry.
Finally, we compare independent constraints on the distances of the absorbers based on the UV and X-ray analysis. In Paper IV, limits on variability in the X-ray absorption were used to give lower limits of $`R`$ 3, 0.6, 0.1 pc for XLI, XMI, and XHI, respectively (see Table 3). Similar lower limits were determined from variability analysis of XMM-Newton observations by Behar et al. (2003). In contrast, a recent study by Krongold et al. (2005) reported variations in the Fe M-shell unresolved transition array in the CXO observations, and they derive an upper limit of $``$ 6 pc for this absorber. Limits can also be derived on the distances for each of the ionization components modeled in Paper IV based on the simple geometrical requirement that $`\mathrm{\Delta }RR`$. If the absorbers are uniform, constant density regions, then for a given model solution with parameters $`U`$ \- $`N_H`$, $`\mathrm{\Delta }RR^2`$ through their dependences on density. This gives upper limits on $`R`$ (lower limit on $`n_H`$); results are listed in Table 3. Filling factors $`<`$1 would imply more stringent limits on $`R`$, since the absorbers would then occupy a larger region than given by $`\mathrm{\Delta }R`$. The limit on XHI, $`R`$ 4 pc, is somewhat less than the distance derived for component 1a based on the C III\* density constraint and UV modeling above; it is similar to the estimates for UV components 2 and 3 based on the assumption of uniform absorbers, with roughly spherical geometries (see discussion in ยง6.1).
### 6.3. Global Model of the Outflow in NGC 3783
#### 6.3.1 Evidence for an Inhomogeneous Wind
With constraints derived on the physical state and geometry of the absorbers, we now investigate implications for the global model of the outflow in NGC 3783 and explore constraints for dynamical models. In Paper IV, the three ionization components modeled for the X-ray absorption were found to be consistent with being in pressure equilibrium ($`PT/U`$), all lying on stable regions of the nearly vertical part of the thermal stability curve (log($`T`$) vs log($`U`$/$`T`$); see Figure 12 in Paper IV). Based on the modeling above, UV kinematic components 1b, 2, and 3 are also at the same pressure. They occupy the low-temperature base of the region of the thermal stability curve where a range of temperatures can co-exist in pressure equilibrium, coinciding with the solution in log($`T`$), log($`U`$/$`T`$) for the low ionization X-ray gas, XLI, seen in Figure 12 in Paper IV. Thus, these results are consistent with the general model presented for NGC 3783 in Paper IV, and described theoretically in Krolik & Kriss (1995, 2001), in which the absorption arises in a multi-phase thermal wind, comprised of embedded regions that are inhomogeneous in temperature and density. In this model, the UV absorbers (and XLI) represent the lowest-ionization, densest material detectable in the spectrum.
These results can be compared further with the recent study by Chelouche & Netzer (2005), which presents a new dynamical model for the X-ray outflow in NGC 3783 that combines some of these general ideas with more specific assumptions and calculations. The main driver of the flow in this model is thermal gas expansion and the main carrier of the flow is the highest ionization, hottest component. The model suggests that cooler, lower-ionization material occupies smaller fractions of the flow where the size distribution resembles what is known from the ISM. An inhomogeneous wind model opens the possibility for having different covering fractions for different ionization components that share the same outflow velocity. It also eases the problem of the large transverse dimension of the X-ray gas in component XLI discussed in Paper IV, since the low-ionization components can have large column densities yet they are composed of small filaments or clouds with relatively small dimensions. In the framework of this model, Chelouche & Netzer (2005) compute the kinetic energy associated with the NGC 3783 outflow to be only a small fraction of the bolometric luminosity, and the mass-loss rate to be comparable to the accretion rate.
The low-ionization, high-velocity UV absorber component 1a does not fit into the picture of inhomogeneities co-existing at pressure equilibrium. It has a gas pressure that is a factor of ten greater than the other components and thus, if embedded in the more diffuse higher-ionization gas without an additional confining mechanism, will eventually evaporate. One possibility is that component 1a is comprised of relatively high density material that has recently been swept up from an external mass source and exposed to the ionizing radiation from the AGN, and is destined to expand to come into pressure equilibrium with the remainder of the flow. This absorber was found to have the smallest covering factor, based on the Si IV doublet absorption, consistent with this being a dense, compact region in the flow. Perhaps it is embedded in, and evaporating into the more diffuse component 1b gas, which is at pressure equilibrium with the rest of the outflow. This would explain the dynamical link between these regions implied by the decrease in radial velocity observed in both absorbers (Paper III). If component 1a is unconfined, the timescale for it to expand, with its density decreasing to that of the 1b region, can be approximated based on the estimated size of the absorber ($`\mathrm{\Delta }R`$ 2$`\times `$10<sup>16</sup> cm) and the thermal velocity ($`v_{th}=`$ 22 km s<sup>-1</sup> at the model temperature of 1.9$`\times `$10<sup>3</sup> K), giving a lifetime of $``$ 150 years. Alternatively, component 1a (and all UV absorbers) could be confined by magnetic pressure, as proposed in some dynamical models (e.g. Emmering et al., 1992; de Kool & Begelman, 1995). Only a moderate field ($`B`$ 10<sup>-3</sup> G) would be required to balance the thermal pressure of this absorber; general calculations by Rees (1987) show fields of this strength could easily be present at the distance derived for the UV absorbers in NGC 3783.
The independent appearance of the UV components on yearly timescales, without correlation with the observed continuum flux (KC01), provides further constraints for physical models of the outflow. In the framework of the thermal wind model described above, it may be a signature of condensations in the outflow that have appeared in our line-of-sight to the AGN. However, we note these timescales for the appearance of absorption components are quite short compared with the evaporation timescale derived for component 1a above. Alternatively, it may be due to motion of the UV absorbers across our sightline, as discussed in KC01 and Paper III. In this scenario, dynamical models must account for the implied transverse component of velocity ($`v`$ 500 km s<sup>-1</sup>, KC01; Paper III), at the $``$ 10 pc distance scale, as well as the radial velocity for the UV absorbers. Additionally, the decreasing radial velocity observed in component 1 must be explained, consistently for both subcomponents 1a and 1b. If this is due to a geometrical effect, in which the absorber is following a curved path across our line-of-sight, it provides constraints on its trajectory and kinematics (see Paper III).
Finally, the combined UV and X-ray results on the outflow geometry and kinematics may provide additional constraints. If all ionization components are roughly co-located, then the gas giving rise to the UV absorbers would occupy a much smaller region than the high-ionization material seen in the X-ray. If the UV absorption components are single uniform clouds, it may be difficult to account for the similarity in outflow velocities for the full range of ionization seen in the two bandpasses; it would require three small, localized regions (the UV absorbers) span the full kinematic range of the much more extended global outflow corresponding to the higher ionization gas. One possibility is that the individual components have low internal filling factors, so that the lower-ionization gas we see is distributed over a larger volume, rather than localized to the small region implied by $`\mathrm{\Delta }R`$ derived above. Another potential explanation is that the gas is not co-located (e.g., see discussion in ยง6.2). This could be explained, for example, by dynamical models that predict the terminal velocities of flow regions will depend on the launch radius of the outflow, such as the MHD wind models presented in Bottorff et al. (2000); in this framework, it is possible that gas with a range of physical conditions and at different locations would have similar terminal velocities.
#### 6.3.2 Connection to the Observed Emission-Line Spectrum
Here we explore a potential connection between the UV absorbers and the optical โ UV line emission. In a recent study, Crenshaw & Kraemer (2005) argued the intrinsic UV absorbers in Seyfert galaxies may be identified with a high-ionization component of gas in the inner narrow line region (NLR), based on their kinematics, ionization, inferred distance, and global covering factor. This is supported by our results for NGC 3783. From STIS slitless spectra of a sample of Seyfert galaxies, Crenshaw & Kraemer (2005) found 9 of 10 objects have bright, central compact \[O III\] emission-line knots with half-width at half-maximum sizes of $``$ 5โ70 pc and half-width at zero intensity radial velocities of 300โ1100 km s<sup>-1</sup>. The derived distances and observed outflow velocities for the intrinsic UV absorbers in NGC 3783 are consistent with this inner NLR component. A further comparison can be made with the line-emission observed in NGC 3783. Comparing to the measurements in Evans (1988), we find the UV absorbers may contribute significantly to some of the high ionization UVโoptical emission-line fluxes. For component 1a, the most constraining line is \[Ne V\] $`\lambda `$3426, with a luminosity of $`L`$10<sup>41</sup> ergs s<sup>-1</sup>, corrected for the NLR reddening determined by Ward & Morris (1984). This can be compared to the prediction from component 1a, giving $`L=`$3$`\times `$10<sup>41</sup> ergs s<sup>-1</sup> $`\times C_g`$, where $`C_g`$ is the global covering factor of the absorber. This limits $`C_g`$ 0.3 for gas with similar physical conditions as component 1a. This component also predicts a strong contribution to the \[O III\] $`\lambda `$5007 emission, and the coronal iron lines (\[Fe VI\] $`\lambda `$5177, \[Fe VII\] $`\lambda `$6085, \[Fe X\] $`\lambda `$6375, and \[Fe XI\] $`\lambda `$7891), reproducing 30 โ 100% of their observed line luminosities for $`C_g=`$1. Absorption component 3, with its relatively high-ionization and large column density, is also predicted to contribute strongly to the higher-ionization coronal iron lines; e.g., it produces $``$3 times the observed \[Fe XI\] $`\lambda `$7891 luminosity for full global coverage, thereby limiting $`C_g`$0.3. It may also contribute to the observed high-ionization UV resonance lines, particularly O VI $`\lambda \lambda `$1032,1037, implying a similar limit on $`C_g`$, with only a small contribution ($``$ 5 โ 10%) coming from resonance scattering. These constraints on $`C_g`$ for the absorbers in NGC 3783 are similar to the general constraint implied by the detection rate of UV absorption in Seyfert 1s ($``$50%, Crenshaw et al., 1999). Finally, we note there is potential evidence for the connection between the UV absorbers and high-ionization line emission in the profiles observed in NGC 3783. Ward & Morris (1984) found a strong asymmetry towards blue wavelengths in the forbidden coronal iron lines, with emission extending to $``$1500 km s<sup>-1</sup>, which is not present in the low-ionization NLR lines such as \[S II\] (see their Figure 6).
## 7. Summary
We have presented analysis of the intrinsic UV absorption in the Seyfert 1 galaxy NGC 3783, based on an intensive monitoring campaign with HST/STIS and FUSE. Our 18 STIS observations included both an intensive phase (3 โ 8 day sampling) and a long-term phase (several months sampling), observing the UV continuum in a range of flux states spanning a factor $``$ 2.5. In the three kinematic components with strong UV absorption (components 1, 2, and 3 at radial velocities $`v_r`$1350, $``$550, and $``$725 km s<sup>-1</sup>, respectively), the equivalent widths of the lowest-ionization lines present were observed to vary inversely with the UV continuum flux. This indicates all components are responding to the ionizing flux, including the short timescales of the intensive observations.
We isolated the NLR emission using variability in the emission-line profiles, in order to explore its effect on the line-of-sight absorption covering factors. We find a coverage model with an unocculted NLR and individual covering factors of the continuum and BLR sources derived from the Lyman lines predicts saturation of several lines that have substantial residual fluxes. These results are consistent with the independent evidence for saturation in these lines from the lack of observed variability and predictions from modeling results. This coverage model is used to obtain ionic column densities in low and high state spectra for photoionization modeling. Based on differences in derived covering factors and kinematic structure, the absorption associated with component 1 appears to be comprised of two physical components (1a and 1b). The decrease in radial velocity observed in lines associated with both regions implies they are linked dynamically and co-located.
We were able to resolve the C III\* $`\lambda `$1175 absorption complex observed in component 1a into its individual multiplet lines and measure column densities for each metastable level. Comparing the ratios of the level populations with our calculations, we determine the electron density for this absorber to be $``$2.5$`\times `$10<sup>4</sup> cm<sup>-3</sup>.
We compare the measured ionic column densities with photoionization models to constrain the physical conditions in each UV component, incorporating variability observed between low and high states. Based on the density from C III\* and the modeling solution, component 1a is found to be located at $``$ 25 pc. It is in a relatively low-ionization state (log($`U`$)$`=`$1.6) compared with the other UV components. Using the density derived for this absorber, we compute time-dependent ionic populations and are able to reproduce the detailed variability observed in this component. Components 1b, 2, and 3 are found to have similar ionization parameters (log($`U`$)$`=`$0.5 โ $``$0.4), with total column densities spanning a range of a factor $``$5. The ionization parameter for these components is the same as that modeled for the lowest-ionization X-ray region from the CXO spectrum, but with smaller total column density. This may indicate different covering factors in the two bandpasses, perhaps due to the different sizes and line-of-sight geometry of the background emission regions. Analysis of the observed variability in components 2 and 3 limits their distances from the central ionizing source to be $``$ 25 and $``$ 50 pc, respectively.
The modeling results for the higher ionization UV components (1b, 2, 3) imply pressures similar to those of the three X-ray ionization components modeled in Paper IV. Thus, the global outflow in NGC 3783 is consistent with an inhomogeneous wind, with relatively small dense regions giving the UV and lower-ionization X-ray absorption, which are embedded in a hotter, more tenuous flow. However, component 1a has a factor of 10 greater pressure than the other absorbers, thus, if embedded in a larger outflow and unconfined it will eventually evaporate, over a timescale of $``$150 years. Alternatively, this absorber could be confined by another mechanism, such as magnetic pressure. The emission-line luminosities predicted from our models of the UV absorbers indicate they may be contributing substantially to the high-ionization lines observed in the UV-optical spectrum. The implied constraints on the global covering factors of the UV absorbers in NGC 3783 are similar to those implied by the detection rate of intrinsic absorption in Seyfert 1 galaxies. Additionally, the distances derived for the UV absorbers are consistent with the location of the inner narrow line region.
Support for proposal 8606 was provided by NASA through a grant from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. Some of the data presented in this paper were obtained from the Multimission Archive at the Space Telescope Science Institute (MAST). Support for MAST for non-HST data is provided by the NASA Office of Space Science via grant NAG5-7584 and by other grants and contracts. We are grateful to Anand Bhatia for insightful discussions and to Gary Ferland for providing us with an updated version of Cloudy. |
warning/0506/cond-mat0506600.html | ar5iv | text | # Comment on โSound Modes broadening in Quasicrystalsโ and its Peer Review
## I The Comment
Recently, de Boissieu et al. deBoissieu proposed a mechanism to explain the broadening of phonons starting from a threshold wavevector in QC, as is systematically observed in experimental data. The mechanism (the coupling of sound waves to a heath wave) is quite general and not new, as it has been known for many years in glasses. There is obviously no doubt that the mechanism itself is sound but we must take issue with the way it is being used through a transcription to the field of quasicrystals. Below, we explain our objections on two levels:
(1) On a general level, phonon broadening is an intrinsic property of quasicrystals, even in a completely harmonic model, such that there is a priori no need for the introduction of an assumption of anharmonicity in the form of a coupling.
(2) On a more detailed level, the authors try to blow new life into a cluster scenario proposed by Janot et al.Janot , by using the idea of localized modes on clusters as a basic ingredient for the microscropic realisation of the mechanism. The localized modes on the clusters are this time no longer directly responsible for the broadening, but they are proposed to be flat modes, that couple to the sound waves. To cite them verbatim: โThe building bricks of all QC structures are atomic clusters. These clusters are not mere geometrical constructions but real physical entities responsible for specific features in the QC vibrational spectrum (e.g., responsible for localized modes)โ.remark
Hence on two levels assumptions are introduced that are not granted or even not needed, while through the presentation the reader might be left with the impression that the experimental data present evidence in support of these assumptions.
(1) There exists an extensive literature on phonons in QC, e.g. on the Fibonacci chain, based on the transfer matrix method.transfer What transpires from such (rigorous) studies is that the phonon eigenmodes are not at all periodic and even not quasiperiodic. (There exist e.g.
so-called recurrent eigenmodes: As a function of its position in space the amplitude of a mode can exhibit humps around the values $`a\times \tau ^n`$, where $`\tau `$ is the golden mean, $`n`$ is the set of integers, and $`a`$ a constant (length). In between the humps, the amplitude of the mode is never zero but it can be very weak). The simple argument of the non-quasiperiodicity of the eigenmodes shows that their Fourier analysis in terms of wave vectors Q will not show dispersion curves of zero line-width as in crystals, but broader features, and that this broadening is not in energy (as de Boissieu et al. think) but in the wave-vector Q. For the recurrent modes, specialists even wonder if it is a mathematically legitimate to take it for granted that they have a Fourier transform in Q-space, and in the case it is not, what kind of information the neutron data contain about these modes.
These features occur even in models that are perfectly harmonic without any coupling between the modes. Of course calculations in the harmonic approximation on successive approximants can never reveal such broadening, because they give rise to zero linewidths by definition. A key intrinsic difficulty of the problem is thus just being missed by such an approach, but it is from the blind spot inherent to such calculations that one might feel the urge to inject additional assumptions of anharmonicity into the problem under the form of localization or couplings.
It would be cheap to push the present objections away, by arguing that real QCs are not one-dimensional, and that surely there will be a loophole of escape from these objections, when we go to higher dimensions. Such vague arguments would (a) reverse the charge of proof, and (b) contain a tacit denial of the horrendous difficulty of eigenvalue problems (with the correct boundary conditions) on quasiperiodic structures. We may add to this that (c) anharmonicity can be experimentally evidenced by the temperature dependence of the Debye-Waller factor.Mossbauer
When one drops full mathematical rigor, approximate eigenmodes that are periodic or quasiperiodic, leading to zero approximate line widths, are only expected to exist in the long-wavelength limit. This correlation between absence of broadening and long wavelengths is confirmed by the experimental observations. But the cluster mechanism proposed by de Boissieu et al. only holds in the limit in which acoustic wavelengths are much larger than a typical cluster size, i.e. the very limit where on the basis of the preceding arguments the broadening of the line widths is rather expected (and observed) to be minor.
(2) In key positions of their paper (the abstract and the final conclusion) the authors stress the important rรดle they propose clusters to play in the microscopic realisation of their mechanism. The authors discuss the relevance of isolated clusters, and present the issue in terms of a unique isolated spherical inclusion in a vibrating medium (e.g. a metallic sphere inside a rubber medium). They dismiss such a scenario with great emphasis, but this does not really clarify the assumptions that underly their paper with respect to the issue if there are isolated clusters in quasicrystals: The scenario chosen to assess this issue is too obviously wrong, and the real issue if there are isolated clusters in quasicrystals is more subtle.
We have to address here a difficult situation of possibly ambiguous terminology, because the authors do indeed introduce a concept of โisolatedโ clusters, different from the one that might be inferred from their presentation. Let us call the type of isolation evoked by the model of a unique spherical inclusion โtype 1โ and the type of isolated clusters used by the authors โtype 2โ. The inadequacy of the type-2 isolated cluster model does not hinge on the abundance of clusters in the structure as one might infer from the type-1 model with its unique cluster. It rather consists in tacitly denying the importance of boundary conditions in the set of coupled differential equations that describe the phonon problem (and which are in general expressed in terms of a dynamical matrix).
The point is easily understood as follows. Take e.g. a one-dimensional crystal that is based on the periodic repetition of the small motif $`LSLLS`$ taken from the Fibonacci chain. The crystal is thus $`\mathrm{}LSLLS.LSLLS.LSLLS\mathrm{}`$. Based on visual clues we could claim that $`LSLLS`$ is a cluster, and that the crystal is a dense packing of clusters. The eigenmodes and eigenvalues of an isolated cluster $`LSLLS`$ are completely different from those of the crystal $`\mathrm{}LSLLS.LSLLS.LSLLS\mathrm{}`$. In fact, the cluster $`LSLLS`$ has six discrete flat modes with a certain dynamical form factor, $`S(๐)`$ that extends throughout reciprocal space. The modes are thus flat due to the finite extension of the cluster in space. The modes of the crystal are completely different: They do not correspond to a few isolated discrete energies, but build a whole dispersion curve and for each eigenmode (i.e. each energy) in the acoustic regime the $`Q`$-dependence is Dirac-like in reciprocal space (if we limit ourselves to one Brillouin zone).
Of course, it would be completely inappropriate to claim on the basis of the visual clue that we can discern clusters $`LSLLS`$ in the crystalline structure that there are flat modes in the crystal, and that there would exist a coupling between sound waves and these flat modes in the crystal. The flaw in such a reasoning is uniquely based on a tacit change of the boundary conditions: It replaces periodic boundary conditions (with a rather smooth variation of the force constants across the โclusterโ boundaries) by an abrupt discontinuity at the โsurfaceโ of the imaginary cluster.
In other words: Type-1 isolation refers to the absence of similar clusters in the surroundings, the cluster is alone. Type-2 isolation refers to a decoupling of the cluster from the surroundings in terms of the force constants. And of course what is relevant for the phonon problem is not if the cluster is alone (type-1 isolation) but how the clusters couple to their surroundings at their supposed boundaries (type-2 isolation). Even in the example of a unique metallic sphere in a rubber medium, it is type-2 isolation that is physically relevant.
QCs are not periodic, and are subject to other boundary conditions than the ones that prevail in a crystal (see below). But this certainly cannot mean that our example would not be appropriate and that the introduction of โclustersโ by the authors would have less of a hidden problem with the prevailing boundary conditions, because the crucial point of our objection lies in the postulated abrupt discontinuity, not in the periodicity. When I am raising an objection against the use of (type 2) isolated clusters, it can thus certainly not imply that I would have missed the passage in the paper where the authors acknowledge that there are no (type 1) isolated clusters in quasicrystals. The statement that there are no (type 1) isolated clusters in QCs (because they are a dense packing of clusters) cannot hide the fact that it is the very use of (type 2) isolated clusters which is the basic ingredient for the microscropic interpretation proposed by the authors.
We could never warn the reader enough against the pitfall that would consist in getting oneโs attention sidetracked towards the issue if there is convincing evidence for the presence of clusters in QCs or otherwise. That would be certainly an interesting topic in its own right, but rather pointless and misleading in the present discussion, as the issue if there are (type-2) isolated clusters in quasicrystals cannot be replaced by an issue if there are clusters in quasicrystals all together, nor by the issue if clusters are physically meaningful in quasicrystals. The verdict on the latter issues will moreover depend on the context of the application: A possible cluster argument in a problem of stability or electronic properties will be different from the one in a phonon problem.
The authors formulate the statement that clusters are not mere geometrical constructions without any proof as though it would be an obvious thruth, and the difficulty that they can overlap is passed under silence. The claim that the origion of the flat modes observed in AlPdMn can be attributed to a localization on clusters is also put forward without any proof.
To introduce the boundary conditions underlying their cluster assumptions, the authors should have given arguments that there is a discontinuity in the force constants at the surface of these clusters. In certain points on the cluster boundaries, the contrary rather seems to be true, viz. when instead of being isolated clusters overlap, which is often the case. In such points it rather looks as though nothing in the whole set of the atomic forces between pairs of atoms in a QC singles out a cluster as an isolated entity, defined by such a discontinuity. The forces between the atoms inside the clusters are not obviously different from those between an atom of the cluster and a neighbouring atom that lies just outside the cluster (but inside the overlapping cluster of the same type). A few phason jumps can create the illusion that a whole cluster has jumped, which also clearly illustrates the relative arbitrariness of assigning an atom to a cluster and of suggesting that a cluster would be an isolated entity whose existence would be obviously defined by a discontinuity in the atomic forces at its surface.
Mathematically spoken, if a cluster is taken large enough it can even be a covering cluster for the whole QC. One can imagine a crystal that could be depictured as a (periodic) arrangement of physically acceptable, overlapping identical clusters of a certain size, and that would not lead to any localization or broadening. Any attempt to escape from this trivial objection must therefore forcedly end up in a discussion of the global, non-periodic arrangement of the clusters and their overlaps.
Discussing QC problems in terms of clusters rather than atoms, is thus just a kind of renormalization procedure, that merely shifts the intrinsic difficulty of non-periodicity to a different length scale, but does not tackle the difficulty itself. It is a blunt denial of the subtlety and difficulty of the eigenvalue problem to overemphasize the rรดle of clusters. We can illustrate this with the Fibonacci chain. It starts with $`LSLLS.LSL.LSLLS.LSLLS.LSL.LSLLS.LSLLS.LLS\mathrm{}.`$, where we have subdivided the sequence in building bricks $`LSLLS`$ and $`LSL`$. Each of the ocurrences $`.LSL.`$ herein is seen to be followed by $`LS`$ as both building bricks $`.LSL.`$ and $`.LSLLS.`$ begin with $`LS`$. Hence the whole sequence can be seen as made from the โcovering clusterโ $`LSLLS`$, whereby we have to allow for overlaps $`LS`$, which appear exactly at the positions where we have separated out $`.LSL.`$. Similarly we could even consider $`LSL`$ as a covering cluster (the overlap would then be $`L`$). Now, the phonon eigenmodes of the isolated sequences $`LSL`$ and $`LSLLS`$ can be calculated (from the corresponding $`4\times 4`$ and $`6\times 6`$ dynamical matrices). What does this handful of eigenmodes tell us about the phonons of the Fibonacci chain? Hardly anything! As we pointed out above, even the phonons of the periodic sequences based on $`LSL`$ or $`LSLLS`$ do not give us the correct picture, despite the fact that in such sequences the clusters are no longer completely isolated (which would be a completely irrealistic boundary condition) and one at least allows for the point that they are embedded in a larger structure (which completely changes the eigenvalue problem).
The idea of clusters $`LSL`$ and $`LSLLS`$ certainly has great eye appeal. One might think at first sight that it must yield great insight in the dynamics of the Fibonacci chain. But as we explained above, all this is mere deception. Already the overlap $`LSLLSLLS`$ of two clusters of the type $`LSLLS`$ will yield competely different solutions for the eigenvalue problem than $`LSLLS`$ itself. The same basic objections about the boundary conditions remain perfectly valid in the three-dimensional case, such that the fact that we work on the one-dimensional case does not present a loophole from these objections. All the use of the clusters $`LSL`$ and $`LSLLS`$ allows us to do is to rewrite the transfer matrix formalism in terms of matrices that correspond to $`LSL`$ and $`LSLLS`$ rather than in terms of the more elementary matrices that correspond to $`S`$ and $`L`$. This illustrates how replacing atoms by clusters is just a renormalization procedure, as we stated. It is an underestimation of the complexity of eigenvalue problems and their boundary conditions (which is global) to suggest that they could be approached locally by focusing oneโs attention to small building bricks. Putting the bricks together just changes everything.
At least in the present context we can thus state that unless a rigorous proof of the contrary is given, it is wise to adopt cautiously the conservative view point that the rigorous application of the idea of clusters, even if they look physically attractive, has remained limited to just a convenient pictorial shorthand to describe parts of the structure, nothing more. We can appreciate from this discussion how both objections (1) and (2) are linked, in the sense that both are based on a tacit modification of highly sensitive details of an eigenvalue problem, that is very hard to spot. The example of how the recurrent modes completely escape the analysis in terms of periodic approximants, shows to what kinds of catastrophies such lack of rigor can lead. Once again, this concern about rigor should not be misrepresented by saying that I would claim that there are no clusters in quasicrystals, or that clusters could not play a role in quasicrystals, etcโฆ
Without any justification, the localized modes invoked are identified with the flat modes that have been reported in AlPdMn, and a coupling mechanism between these localized modes and sound waves is proposed. We have two objections to this:
(a) Such an explanation for the flat modes is just one between several other possibilities. One of the alternatives is documented and can therefore not be ignored: By a scrutiny of the displacement patterns in their numerical simulations Hafner and Krajci Hafner were able to associate the flat modes with a restriction (โconfinementโ) of the vibrations to disclination lines of atoms that are topologically different from average (e.g. the atoms have a 13-fold coordination, rather than a normal 12-fold one). This has nothing to do with the vibration on a cluster.
(b) The issue if the flat modes are due to a localization on clusters is not open-ended within the present state of knowledge. It can be unambiguously settled. It suffices to check if the structure factor of a flat mode is indeed compatible with the dynamical structure factor of a cluster vibration (as Buchenau has done to prove his model for the dynamics of silica). Although the dynamical structure factors of the flat modes have not been published, it must be straightforward to extract this first-rank information from the authorsโ already existing data, and a numerical calculation of the vibrational spectrum of a Bergman or a Mackay cluster with realistic force constants, involving typically 33 to 55 atoms, is certainly not unfeasible.
Hence, before one can formulate any possible approach of the type proposed by the authors, it is a peremptory prerequiste that one first proves on the basis of existing data, that (1) the observed structure factors of the flat modes are compatible with an interpretation of these modes in terms of cluster phonons, and (2) that there are anharmonicities within the system, e.g. on the basis of Debye-Waller factor anomalies of which one has proved beyond any doubt that they cannot possibly be attributed to an onset phason hopping. These are necessary but not sufficient, minimal conditions that have to be met. They stand completely free from any theoretical considerations, and therefore add up completely independently to the two main objections outlined in the present Comment.
## II The Peer Review of the Comment
We would like to give the reader an inkling about the methods that are used to have valid Comments of this type rejected. I would have prefered to quote the referee reports literally rather than paraphrasing them but it has been pointed out to me that it is illegal to reproduce referee reports literally as this constitutes a violation of copyright. I must state that I feel very uncomfortable about this. First of all, it kind of diminishes my credibility and exposes me to cheap and easy criticism that I am distorting the truth because I would not reproduce what has been written literally. Secondly, when one literally reproduces what has been written, one cannot be accused of being responsible for whatever that might be contained in it and look disgracious, while when one has to paraphrase it, one becomes subjectively accredited with this responsibility. Thirdly, it is apparently not enough that anonymous peer review exposes people almost defenseless to the huge prejudices that can be inflicted by sham peer review, especially when it becomes systematical if some group has managed to completely invade the horizon of an editorial board. Victims are this way also obstructed from denouncing what they have undergone and making the community aware of it. I must ask the reader to consider how this obligation to mention what has been stated only indirectly, can only result in a down-sized and filtered account of the adversity and the personal attacks I have been faced with.
In a first reply de Boissieu stated that the verbatim quotation I made at the beginning of my Comment: โThe building bricks of all QC structures are atomic clusters. These clusters are not mere geometrical constructions but real physical entities responsible for specific features in the QC vibrational spectrum (e.g., responsible for localized modes)โ would not be in his paper, and that my Comment would distort the point of view of his paper by making quotations out of context.
He also stated that it was not appropriate to take issue with the fact that his analysis is based on isolated clusters, while on p.5 of his paper he had clearly stated that the picture of one isolated cluster is not adequate for QCs, and that one should rather think of a QC in terms of a dense packing of clusters. The reader should not get confused by this would-be catching me in my own contradictions. The contradictions are entirely from the hand of de Boissieu et al. themselves. Yes, de Boissieu et al. state with great emphasis in the beginning of their paper that they do not use isolated clusters, refering to the example of a metallic sphere in a rubber medium. But, no, this cannot hide that the whole argument is exactly based on an assumed isolation of the clusters. It is just that the isolation at stake is very different from the one suggested by the example of a single metallic sphere in a rubber medium. It is not by sorting two entirely different situations under a same descriptive phrasing that they would become equal. In fact, the problem of isolated clusters is not one of numbers (โoneโ against โa dense packingโ), as one could infer from this reply, but one of mutual overlap. In order to overcome the confusion that could result from this reply, I introduced the definitions and the distinction between type-1 and type-2 isolation in my Comment.
de Boissieu further stated that the scientific content of my comment would be very small, that it contained vague and general views, that were in lack of scientific evidence and that were not supported by recent papers. He stated that I overinterpreted the literature.
Finally, he replaced the true issue of my Comment, viz. that the clusters are not isolated, by another issue, viz. if there are clusters in quasicrystals at all, and then gave a detailed reply on this replacement issue, with several citations. Even if this reply had been entirely correct, it did not address the issues I had raised in my Comment.
With respect to my argument that it is not true that clusters are isolated in the sense of having significantly stronger intra-cluster bonding, de Boissieu stated that he certainly agreed that the existence of clusters is still a matter of debate, but that it has not been proved that the forces are of equal strength throughout the QC and that therefore my argument had no firm ground. He added that the problem of the energetics of QCs is very complicated and one could not expect the final solutions to be given soon.
With respect to my citation of the work of Krajci and Hafner he stated that this was wrong, because the five flat modes reported by these authors have an insignificant participation ratio, and because there has been no analysis of their vibration patterns, such that their true nature is still not clear, even if it is clear that they are associated with disinclination lines.
It was also stated that electron density measurements on a cristalline approximant of AlReSi indicate a larger bonding character within clusters than in between them. (This argument is clearly not general as can immediately be appreciated from the fact that clusters very often overlap, which is the real issue of my Comment).
It was also stated that Gratias approved the cluster approach, while it very clearly transpires from Gratiasโ papers that he rather very cautiously considers them as a convenient shorthand to describe the structure.
It suffices to point out that all this does not address at all the issue I raised, which is that the clusters often overlap. The point is thus not if there are clusters at all in quasicrystals, but that these clusters are not type-2 isolated.
After I had replied to this, the correspondence was sent to two referees. One of them appeared biased to me. But eventually, both of them recommended publication with small modifications. However, I learned later on that one of them wrote a seperate note to the editors wherein he/she stated that he/she did not wish to review my manuscript any further, because the manuscript would not be presented in the proper manner. (The โhe/sheโ reflect a terminology that was used by the APS).
When I had adapted my version to these comments, the editors of PRB requested that I should remove the sentence from my paper that very explicitly stated that one should be careful in not being sidetracked towards the fake issue if there are clusters in QCs, as the true issue is if these clusters are isolated in a very specific sense. I had made a citation towards de Boissieuโs reply to the editors in this respect. It was argued that I would not have the right to cite correspondence from the peer review process. It was also stated that the editorial board was โpositively inclinedโ to accept my Comment for publication. Judging that the way de Boissieu changed the issues was apt to mislead a many reader, I reformulated the sentence such as to keep its contents but to remove the citation. Then PRB put my paper โon holdโ for six months, refusing to give any explanation. They had done that already a first time in the review process. Perhaps I should have understood from this that my Comment was not well considered by the APS itself, and that this gentle use of force was meant to be discouraging enough to make me just give up. Finally after nine months, they sent me all at once a report from a fourth referee. There had very obviously not been the slightest reason to ask advice from a fourth referee.
This referee replaced again the issue if the clusters are isolated by the false issue if there are clusters all together. On this replacement issue he stated that an intense debate was going on within the community, and again developed a whole argumentation about this non-issue, citing the model of closed electronic shells by Janot et al., the fracture experiments by Ebert et al., confirmed by numerical simulations by Roesch et al. and the cluster friction model by Feuerbacher et al. It was also stated that the whole development based on the Fibonacci chain was insignificant, and unsuitable to disprove the assumption that there are clusters in QCs, because real three-dimensional clusters are very different from sections of the Fibonacci chain (e.g. in containing many shells). That the looks of one- and three-dimensional clusters are different may very well be true, but is irrelevant for the real issue, which is not if there are clusters in QCs, but that these clusters are not separated. And the latter issue can equally well be explained on the Fibonacci chain as on a three-dimensional model. It was also stated that as long there was no rigorous proof that there are no clusters in QCs, it would be allowed to assume that clusters are present and to build models on this assumption. Once again, the issue is not if there are clusters all together, but that these clusters are not isolated.
The referee also paraphrased me by stating that I would have claimed that the considerations of de Boissieu are redundant, because phonon broadening is an intrinsic property of QCs. He argued that this might be plausible, but that it would be too simplistic, because the broadening cannot be calculated, while de Boissieu et al. would convincingly explain the experimental data.
I responded more or less as follows (I am forced to paraphrase the exact statements from the report in order not to violate copyright):
I clearly wrote in my Comment that โWe could never warn the reader enough against the pitfall that would consist in getting oneโs attention sidetracked towards the issue if there is convincing evidence for the presence of clusters in QCs or otherwise.โ This sentence summarizes a number of arguments that are clearly developed in my paper: The issue is not if there are clusters or otherwise. The issue is that de Boissieu et al. ignore the consequences of the crucial fact that the clusters often overlap, and that they therefore are not isolated (in the type-2 sense defined in my paper). (And the issue that the clusters are not isolated is a very different one from the one that de Boissieu dismissed after evoking an isolated metallic sphere in a rubber medium, which of the type-1 defined in my paper). The fact that I define two types of isolation already clearly shows that the issue is not if there are clusters all together in QCs. The issue is that these clusters are not isolated in the type-2 sense. I have a whole discussion of this in terms of boundary conditions within the paper. The referee finds it convenient to ignore this all together, and cites work of Janot et al., Ebert et al., Roesch et al. and Feuerbacher et al. as proofs for the existence of clusters while this is totally pointless, as clearly explained in my paper and the sentence I quoted from it above. The referee thus goes resolutely for the pitfall I warned against in my paper, which I formulated because de Boissieu had attempted this elusive move already in his first reply.
The referee builds further on this swap towards a false issue, by stating that the example of the Fibonacci chain, which would take a disproportionate large part of my Comment could not be used to prove the cluster assumption wrong. In my paper, the Fibonacci chain is not being used in order to prove that there would be no clusters in a QC. As I already pointed out above, the issue if there are clusters or otherwise is pointless for our discussion. What matters is that the clusters often overlap (i.e. are not isolated in the type-2 sense) and in order to point this out, the Fibonacci chain is as good as a full-fledged 3D model. And the referee further insists on focusing the attention to this pointless issue of the existence of clusters or otherwise when he writes that three-dimensional clusters are ot comparable to sections of the Fibonacci chain, and that it is legitimate to use the cluster assumption as long as there is no striking argument against it.
The latter contains actually a reversal of de Boissieuโs charge of proof: One could make as many unproved claims as one likes, it is up to others to prove them wrong. de Boissieu and the referee know very well that this is not viable. This shines through clearly enough when it comes down to attacking my work, rather than defending the work of de Boissieu, and one tries to make prevail totally irrelevant objections against it to make us wonder if perhaps I did not fail to meet my charge of proof on some very tiny loopholes.
One of these irrelevant objections is that the Fibonacci chain would not be pertinent in the present discussion. The arguments developed on the Fibonacci chain are not affected by his irrelevant distinctions between the Fibonacci chain and 3D QCs he would like us to believe crucial. The Fibonacci chain is perfectly suited for discussing the type-2 isolation and other issues at stake. Making the same points on a 2D or 3D QC would require to include into the paper elaborate Figures to show the clusters and how they overlap, etcโฆ While with the Fibonacci chain one can describe the whole situation exhaustively and very clearly by referring to the letters L and S. One issue of my paper is that the possible overlap of the clusters (the lack of type-2 isolation) is pointing out a tacit cheat about the values of the interatomic forces: It tacitly implies that the bonds between atoms within the cluster are stronger than the bonds between atoms of a cluster and surrounding atoms. That this is not true is totally obvious when two clusters A and B overlap, as discussed in my paper. And the referee can check it also on a 3D model. It is nothing specific for the Fibonacci chain only. This entails that the clusters are not at all doing what de Boissieu claims they are doing. This point is exactly the same one as voiced by Henley in his paper โClusters, phason elasticity, and entropic stabilization: a theory perspectiveโ (Phil. Mag. 86, 1123 (2006): Ames conference proceedings) in the first sentence of the section โ2. Clustersโ on page 1124. It is this issue, and not the mere presence of clusters or otherwise, that is essential and makes de Boissieuโs position untenable. To make this untenable position prevail nevertheless, the referee carefully eludes discussing this crucial point of overlap. He diverts the attention away from it by hammering incessantly on the irrelevant issue if clusters exist all together.
I also addressed the paraphrasing of my argument that the use of clusters would be redundant. In fact, this is a completely false presentation of the issues, as anyone who reads the paper carefully can see. The referee operates a very subtle shift when he paraphrases my objection by stating that I would have claimed the considerations of de Boissieu et al. are redundant. As far as I can see I wrote: โThere is a priori no need for the introduction of an assumption of anharmonicity in the form of a coupling.โ If one thinks carefdully about it, this does not mean that the considerations are redundant, but that they could be redundant. The snag to this almost subliminal shift is that if I had claimed that the consideratrions are redundant I would be invested with a charge to prove it. While if the considerations could be redundant, it is the charge of proof of de Boissieu that has not been properly met. Hence this is a hidden reversal of the charge of proof, that is very hard to spot. That phonon broadening is an intrinsic property is not merely plausible as he states, it is a mathematical certainty, because the eigenmodes are not quasiperiodic. Hence here the referee unduly questions (again in an inoffensive looking way) an established obvious factual mathematical truth. That nobody is able to solve the horrendously difficult problem of the calculation of the q-dependence of the intrinsic broadening, does not exclude that the broadening observed could be entirely due to this intrinsic broadening. Let one please not jump onto this sentence to put again things in my mouth that I do not say. I do not say that the broadening is entirely intrinsic, I say that cannot be excluded that it could be entirely intrinsic. The mechanism of intrinsic broadening has at least the merrit that it is physically sound, while the cluster scenario is conclusively proved wrong by e.g. the type-2 isolation issue, which the referee carefully eludes to discuss.
It is therefore ridiculous to exploit the difficulty of the problem of calculating the intrinsic broadening to compare this scenario unfavourably with the (illusory) succes of the cluster model by stating that the result of the $`Q`$-dependence of the line width presented by Boissieu et al. is convincing. The model is proved wrong and that the data can be fitted with a polynomial of the fourth degree is hardly informative and a finding that could be derived from scores of other models.
I have never stated that the broadening would be uniquely intrinsic. As I explained it already above I have only evoked this as a possibility. Because, what the argument of the intrinsic broadening was intended to show is that the assumptions de Boissieu makes are in lack of justification, by giving a counter example. E.g. the tacitly implied assumption that there is anharmonicity is gratuitous. In view of the possibility of intrinsic broadening which will occur even in a completely harmonic model, it is a peremptory prerequiste to prove that the forces are anharmonic before one can introduce the assumptions that underly de Boissieuโs model. Asking to develop the intrinsic broadening scenario into a full calculation as the referee does is, again, a reversal of the charge of proof. Moreover it tries to saddle me up with the obligation of a demonstratio diabolica.
The editorial board of PRB refused to consider my reply to this referee report. They even refused to state this refusal. They just moved on towards a statement that this ended the review of my Comment. They eluded answering by addressing non-scientific issues, maintained my Comment rejected and even eluded responding to an appeal of mine. They had artificially made the whole procedure last for more than two years. They even recommended that I try to have it published in another journal, because other journals have other criteria for approval than the APS.
We may finally point out that we already had attempted to write a Comment on the artificial cluster issue, back in 1993 when it was introduced in reference Janot . Janot et al. dropped an off-hand comment on my work towards the end of that paper that my interpretation of the quasielastic data in terms of phason hopping would be wrong. In reality, their data did not warrant such questioning of my work. In fact, reference Janot reported a failure to observe the quasielastic scattering that I had measured and that corresponded to the decrease of the elastic intensity when the temperature was raised. Such elastic data can never be used to challenge the interpretation of the much more detailed and specific quasielastic data. Nevertheless, Janot et al. did this, denigrating my work. I had to discover this as an accomplished fact in the published literature. To undo the damage, I was forced to write a Comment to reference Janot , with reversed rights of reply. Using this reversed situation, Janot answered that my samples were suspicious and that the quasielastic signal was due to preferential segregation of Cu into the grain boundaries. On the editorial board of Physical Review Letters S. Moss stated that he felt much more sympathetic towards my arguments, but that the exchange would be too long to publish it in Physical Review Letters. S. Moss and R. Schuhmann suggested to me that I send it to Physical Review B. But when I did this, I was told that Physical Review B does not handle Comments on papers published in Physical Review Letters. In my Comment I had pointed out that the cluster scenario was analogous to scenarios used in glasses. But in reference deBoissieu it is stated towards the end, that after the work was finished, the authors discovered that a similar approach had been used in glasses.
After the rejection of my Comment on Janot , a proposal of mine for beam time at the ILL to measure phason dynamics was rejected on the basis a statement by Dubois in the scientific evaluation committee that the experiment had already been done by Janot. It was just not true. When I protested, and I expressed my fears that my ideas would be stolen, Janot wrote a letter to me with copy to the director of the ILL, wherein he stated that I would be paranoid, and that they did not intend to measure phason dynamics. A few months later he and de Boissieu made the experiment in my place on IN16, but they melted their sample. They had made their attempt to measure phason dynamics with the same type of sample, on the same type of instrument, in the very same Q-range, with the same energy resolution, and in the same energy and temperature range. Nevertheless, they wrote an ILL report about it wherein they stated that this experiment would be different from mine and wherein they reported that they had figured out in the meantime that the interpretation of the Debye-Waller factor in Janot was wrong. The interpretation of the temperature dependence of the Debye-Waller factor had been the only element of justification on which the whole introduction of the cluster issue had been based. It was wrong. And already at that stage, the obvious error in the reasoning had been that the clusters are not isolated but overlap. Nevertheless, these issues were introduced again in reference deBoissieu . |
warning/0506/math0506201.html | ar5iv | text | # Metric cotype
## 1. Introduction
In 1976 Ribe (see also , , , ) proved that if $`X`$ and $`Y`$ are uniformly homeomorphic Banach spaces then $`X`$ is finitely representable in $`Y`$, and vice versa ($`X`$ is said to be finitely representable in $`Y`$ if there exists a constant $`K>0`$ such that any finite dimensional subspace of $`X`$ is $`K`$-isomorphic to a subspace of $`Y`$). This theorem suggests that โlocal propertiesโ of Banach spaces, i.e. properties whose definition involves statements about finitely many vectors, have a purely metric characterization. Finding explicit manifestations of this phenomenon for specific local properties of Banach spaces (such as type, cotype and super-reflexivity), has long been a major driving force in the bi-Lipschitz theory of metric spaces (see Bourgainโs paper for a discussion of this research program). Indeed, as will become clear below, the search for concrete versions of Ribeโs theorem has fueled some of the fieldโs most important achievements.
The notions of type and cotype of Banach spaces are the basis of a deep and rich theory which encompasses diverse aspects of the local theory of Banach spaces. We refer to , , , , , , , , and the references therein for background on these topics. A Banach space $`X`$ is said to have (Rademacher) type $`p>0`$ if there exists a constant $`T<\mathrm{}`$ such that for every $`n`$ and every $`x_1,\mathrm{},x_nX`$,
(1) $`๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j_X^pT^p{\displaystyle \underset{j=1}{\overset{n}{}}}x_j_X^p.`$
where the expectation $`๐ผ_\epsilon `$ is with respect to a uniform choice of signs $`\epsilon =(\epsilon _1,\mathrm{},\epsilon _n)\{1,1\}^n`$. $`X`$ is said to have (Rademacher) cotype $`q>0`$ if there exists a constant $`C<\mathrm{}`$ such that for every $`n`$ and every $`x_1,\mathrm{},x_nX`$,
(2) $`๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j_X^q{\displaystyle \frac{1}{C^q}}{\displaystyle \underset{j=1}{\overset{n}{}}}x_j_X^q.`$
These notions are clearly linear notions, since their definition involves addition and multiplication by scalars. Ribeโs theorem implies that these notions are preserved under uniform homeomorphisms of Banach spaces, and therefore it would be desirable to reformulate them using only distances between points in the given Banach space. Once this is achieved, one could define the notion of type and cotype of a metric space, and then hopefully transfer some of the deep theory of type and cotype to the context of arbitrary metric spaces. The need for such a theory has recently received renewed impetus due to the discovery of striking applications of metric geometry to theoretical computer science (see , , and the references therein for part of the recent developments in this direction).
Enfloโs pioneering work , , , resulted in the formulation of a nonlinear notion of type, known today as Enflo type. The basic idea is that given a Banach space $`X`$ and $`x_1,\mathrm{},x_nX`$, one can consider the linear function $`f:\{1,1\}^nX`$ given by $`f(\epsilon )=_{j=1}^n\epsilon _jx_j`$. Then (1) becomes
(3)
$$\begin{array}{c}๐ผ_\epsilon f(\epsilon )f(\epsilon )_X^pT^p\underset{j=1}{\overset{n}{}}๐ผ_\epsilon f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n)\hfill \\ \hfill f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n)_X^p.\end{array}$$
One can thus say that a metric space $`(,d_{})`$ has Enflo type $`p`$ if there exists a constant $`T`$ such that for every $`n`$ and every $`f:\{1,1\}^n`$,
(4)
$$\begin{array}{c}๐ผ_\epsilon d_{}(f(\epsilon ),f(\epsilon ))^pT^p\underset{j=1}{\overset{n}{}}๐ผ_\epsilon d_{}(f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n),\hfill \\ \hfill f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n))^p.\end{array}$$
There are two natural concerns about this definition. First of all, while in the category of Banach spaces (4) is clearly a strengthening of (3) (as we are not restricting only to linear functions $`f`$), it isnโt clear whether (4) follows from (3). Indeed, this problem was posed by Enflo in , and in full generality it remains open. Secondly, we do not know if (4) is a useful notion, in the sense that it yields metric variants of certain theorems from the linear theory of type (it should be remarked here that Enflo found striking applications of his notion of type to Hilbertโs fifth problem in infinite dimensions , , , and to the uniform classification of $`L_p`$ spaces ). As we will presently see, in a certain sense both of these issues turned out not to be problematic. Variants of Enflo type were studied by Gromov and Bourgain, Milman and Wolfson . Following we shall say that a metric space $`(,d_{})`$ has BMW type $`p>0`$ if there exists a constant $`K<\mathrm{}`$ such that for every $`n`$ and every $`f:\{1,1\}^n`$,
(5)
$$\begin{array}{c}๐ผ_\epsilon d_{}(f(\epsilon ),f(\epsilon ))^2K^2n^{\frac{2}{p}1}\underset{j=1}{\overset{n}{}}๐ผ_\epsilon d_{}(f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n),\hfill \\ \hfill f(\epsilon _1,\mathrm{},\epsilon _{j1},\epsilon _j,\epsilon _{j+1},\mathrm{},\epsilon _n))^2.\end{array}$$
Bourgain, Milman and Wolfson proved in that if a Banach space has BMW type $`p>0`$ then it also has Rademacher type $`p^{}`$ for all $`0<p^{}<p`$. They also obtained a nonlinear version of the Maurey-Pisier theorem for type , , yielding a characterization of metric spaces which contain bi-Lipschitz copies of the Hamming cube. In Pisier proved that for Banach spaces, Rademacher type $`p`$ implies Enflo type $`p^{}`$ for every $`0<p^{}<p`$. Variants of these problems were studied by Naor and Schechtman in . A stronger notion of nonlinear type, known as Markov type, was introduced by Ball in his study of the Lipschitz extension problem. This important notion has since found applications to various fundamental problems in metric geometry , , , ,
Despite the vast amount of research on nonlinear type, a nonlinear notion of cotype remained elusive. Indeed, the problem of finding a notion of cotype which makes sense for arbitrary metric spaces, and which coincides (or almost coincides) with the notion of Rademacher type when restricted to Banach spaces, became a central open problem in the field.
There are several difficulties involved in defining nonlinear cotype. First of all, one cannot simply reverse inequalities (4) and (5), since the resulting condition fails to hold true even for Hilbert space (with $`p=2`$). Secondly, if Hilbert space satisfies an inequality such as (4), then it must satisfy the same inequality where the distances are raised to any power $`0<r<p`$. This is because Hilbert space, equipped with the metric $`xy^{r/p}`$, is isometric to a subset of Hilbert space (see , ). In the context of nonlinear type, this observation makes perfect sense, since if a Banach space has type $`p`$ then it also has type $`r`$ for every $`0<r<p`$. But, this is no longer true for cotype (in particular, no Banach space has cotype less than $`2`$). One viable definition of cotype of a metric space $`X`$ that was suggested in the early 1980s is the following: Let $``$ be a metric space, and denote by $`\mathrm{Lip}()`$ the Banach space of all real-valued Lipschitz functions on $``$, equipped with the Lipschitz norm. One can then define the nonlinear cotype of $``$ as the (Rademacher) cotype of the (linear) dual $`\mathrm{Lip}()^{}`$. This is a natural definition when $``$ is a Banach space, since we can view $`\mathrm{Lip}()`$ as a nonlinear substitute for the dual space $`^{}`$ (note that in it is shown that there is a norm $`1`$ projection from $`\mathrm{Lip}()`$ onto $`^{}`$). With this point of view, the above definition of cotype is natural due to the principle of local reflexivity , . Unfortunately, Bourgain has shown that under this definition subsets of $`L_1`$ need not have finite nonlinear cotype (while $`L_1`$ has cotype $`2`$). Additionally, the space $`\mathrm{Lip}(M)^{}`$ is very hard to compute, for example it is an intriguing open problem whether even the unit square $`[0,1]^2`$ has nonlinear cotype $`2`$ under the above definition.
In this paper we introduce a notion of cotype of metric spaces, and show that it coincides with Rademacher cotype when restricted to the category of Banach spaces. Namely, we introduce the following concept:
###### Definition 1.1 (Metric cotype).
Let $`(,d_{})`$ be a metric space and$`q>0`$. The space $`(,d_{})`$ is said to have metric cotype $`q`$ with constant $`\mathrm{\Gamma }`$ if for every integer $`n`$, there exists an even integer $`m`$, such that for every $`f:_m^n`$,
(6) $`{\displaystyle \underset{j=1}{\overset{n}{}}}๐ผ_x\left[d_{}(f\left(x+{\displaystyle \frac{m}{2}}e_j\right),f(x))^q\right]\mathrm{\Gamma }^qm^q๐ผ_{\epsilon ,x}\left[d_{}(f(x+\epsilon ),f(x))^q\right],`$
where the expectations above are taken with respect to uniformly chosen $`x_m^n`$ and $`\epsilon \{1,0,1\}^n`$ (here, and in what follows we denote by $`\{e_j\}_{j=1}^n`$ the standard basis of $`^n`$). The smallest constant $`\mathrm{\Gamma }`$ with which inequality (6) holds true is denoted $`\mathrm{\Gamma }_q()`$.
Several remarks on Definition 1.1 are in order. First of all, in the case of Banach spaces, if we apply inequality (6) to linear functions $`f(x)=_{j=1}^nx_jv_j`$, then by homogeneity $`m`$ would cancel, and the resulting inequality will simply become the Rademacher cotype $`q`$ condition (this statement is not precise due to the fact that addition on $`_m^n`$ is performed modulo $`m`$ โ see Section 5.1 for the full argument). Secondly, it is easy to see that in any metric space which contains at least two points, inequality (6) forces the scaling factor $`m`$ to be large (see Lemma 2.3) โ this is an essential difference between Enflo type and metric cotype. Finally, the averaging over $`\epsilon \{1,0,1\}^n`$ is natural here, since this forces the right-hand side of (6) to be a uniform average over all pairs in $`_m^n`$ whose distance is at most $`1`$ in the $`\mathrm{}_{\mathrm{}}`$ metric.
The following theorem is the main result of this paper:
###### Theorem 1.2.
Let $`X`$ be a Banach space, and $`q[2,\mathrm{})`$. Then $`X`$ has metric cotype $`q`$ if and only if $`X`$ has Rademacher cotype $`q`$. Moreover,
$$\frac{1}{2\pi }C_q(X)\mathrm{\Gamma }_q(X)90C_q(X).$$
Apart from settling the nonlinear cotype problem described above, this notion has various applications. Thus, in the remainder of this paper we proceed to study metric cotype and some of its applications, which we describe below. We believe that additional applications of this notion and its variants will be discovered in the future. In particular, it seems worthwhile to study the interaction between metric type and metric cotype (such as in Kwapienโs theorem ), the possible โMarkovโ variants of metric cotype (ร la Ball ) and their relation to the Lipschitz extension problem, and the relation between metric cotype and the nonlinear Dvoretzky theorem (see , for information about the nonlinear Dvoretzky theorem, and for the connection between cotype and Dvoretzkyโs theorem).
### 1.1. Some applications of metric cotype
#### 1) A nonlinear version of the Maurey-Pisier theorem.
Given two metric spaces $`(,d_{})`$ and $`(๐ฉ,d_๐ฉ)`$, and an injective mapping $`f:๐ฉ`$, we denote the distortion of $`f`$ by
$$\mathrm{dist}(f):=f_{\mathrm{Lip}}f^1_{\mathrm{Lip}}=\underset{\begin{array}{c}x,y\\ xy\end{array}}{sup}\frac{d_๐ฉ(f(x),f(y))}{d_{}(x,y)}\underset{\begin{array}{c}x,y\\ xy\end{array}}{sup}\frac{d_{}(x,y)}{d_๐ฉ(f(x),f(y))}.$$
The smallest distortion with which $``$ can be embedded into $`๐ฉ`$ is denoted $`c_๐ฉ()`$; i.e.,
$$c_๐ฉ():=inf\{\mathrm{dist}(f):f:๐ฉ\}.$$
If $`c_๐ฉ()\alpha `$ then we sometimes use the notation $`\stackrel{๐ผ}{}๐ฉ`$. When $`๐ฉ=L_p`$ for some $`p1`$, we write $`c_๐ฉ()=c_p()`$.
For a Banach space $`X`$ write
$$p_X=sup\{p1:T_p(X)<\mathrm{}\}\mathrm{and}q_X=inf\{q2:C_q(X)<\mathrm{}\}.$$
$`X`$ is said to have nontrivial type if $`p_X>1`$, and $`X`$ is said to have nontrivial cotype if $`q_X<\mathrm{}`$.
In Pisier proved that $`X`$ has no nontrivial type if and only if for every $`n`$ and every $`\epsilon >0`$, $`\mathrm{}_1^n\stackrel{1+\epsilon }{}X`$. A nonlinear analog of this result was proved by Bourgain, Milman and Wolfson (see also Pisierโs proof in ). They showed that a metric space $``$ does not have BMW type larger than $`1`$ if and only if for every $`n`$ and every $`\epsilon >0`$, $`(\{0,1\}^n,_1)\stackrel{1+\epsilon }{}`$. In Maurey and Pisier proved that a Banach space $`X`$ has no nontrivial cotype if and only for every $`n`$ and every $`\epsilon >0`$, $`\mathrm{}_{\mathrm{}}^n\stackrel{1+\epsilon }{}X`$. To obtain a nonlinear analog of this theorem we need to introduce a variant of metric cotype (which is analogous to the variant of Enflo type that was used in .
###### Definition 1.3 (Variants of metric cotype ร la Bourgain, Milman andWolfson).
Let $`(,d_{})`$ be a metric space and $`1pq`$. We denote by $`\mathrm{\Gamma }_q^{(p)}()`$ the least constant $`\mathrm{\Gamma }`$ such that for every integer $`n`$ there exists an even integer $`m`$, such that for every $`f:_m^n`$,
(7)
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}๐ผ_x\left[d_{}(f\left(x+\frac{m}{2}e_j\right),f(x))^p\right]\hfill \\ \hfill \mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}๐ผ_{\epsilon ,x}\left[d_{}(f(x+\epsilon ),f(x))^p\right],\end{array}$$
where the expectations above are taken with respect to uniformly chosen $`x_m^n`$ and $`\epsilon \{1,0,1\}^n`$. Note that $`\mathrm{\Gamma }_q^{(q)}()=\mathrm{\Gamma }_q()`$. When $`1p<q`$ we shall refer to (7) as a weak metric cotype $`q`$ inequality with exponent $`p`$ and constant $`\mathrm{\Gamma }`$.
The following theorem is analogous to Theorem 1.2.
###### Theorem 1.4.
Let $`X`$ be a Banach space, and assume that for some $`1p<q`$, $`\mathrm{\Gamma }_q^{(p)}(X)<\mathrm{}`$. Then $`X`$ has cotype $`q^{}`$ for every $`q^{}>q`$. If $`q=2`$ then $`X`$ has cotype $`2`$. On the other hand,
$$\mathrm{\Gamma }_q^{(p)}(X)c_{pq}C_q(X),$$
where $`c_{pq}`$ is a universal constant depending only on $`p`$ and $`q`$.
In what follows, for $`m,n`$ and $`p[1,\mathrm{}]`$ we let $`[m]_p^n`$ denote the set $`\{0,1,\mathrm{},m\}^n`$, equipped with the metric induced by $`\mathrm{}_p^n`$. The following theorem is a metric version of the Maurey-Pisier theorem (for cotype):
###### Theorem 1.5.
Let $``$ be a metric space such that $`\mathrm{\Gamma }_q^{(2)}()=\mathrm{}`$ for all $`q<\mathrm{}`$. Then for every $`m,n`$ and every $`\epsilon >0`$,
$$[m]_{\mathrm{}}^n\stackrel{1+\epsilon }{}.$$
We remark that in Maurey and Pisier prove a stronger result, namely that for a Banach space $`X`$, for every $`n`$ and every $`\epsilon >0`$, $`\mathrm{}_{p_X}^n\stackrel{1+\epsilon }{}X`$ and $`\mathrm{}_{q_X}^n\stackrel{1+\epsilon }{}X`$. Even in the case of nonlinear type, the results of Bourgain, Milman and Wolfson yield an incomplete analog of this result in the case of BMW type greater than $`1`$. The same phenomenon seems to occur when one tries to obtain a nonlinear analog of the full Maurey-Pisier theorem for cotype. We believe that this issue deserves more attention in future research.
#### 2) Solution of a problem posed by Arora, Lovรกsz, Newman, Rabani,Rabinovich and Vempala.
The following question appears in \[3, Conj. 5.1\]:
> Let $``$ be a *baseline* metric class which does not contain all finite metrics with distortion arbitrarily close to $`1`$. Does this imply that there exists $`\alpha >0`$ and arbitrarily large $`n`$-point metric spaces $`_n`$ such that for every $`๐ฉ`$, $`c_๐ฉ(_n)(\mathrm{log}n)^\alpha `$?
We refer to \[3, ยง2\] for the definition of baseline metrics, since we will not use this notion in what follows. We also refer to for background and motivation from combinatorial optimization for this problem, where several partial results in this direction are obtained. An extended abstract of the current paper also contains more information on the connection to Computer Science. Here we apply metric cotype to settle this conjecture positively, without any restriction on the class $``$.
To state our result we first introduce some notation. If $``$ is a family of metric spaces we write
$$c_{}(๐ฉ)=inf\{c_{}(๐ฉ):\}.$$
For an integer $`n1`$ we define
$$๐_n()=sup\{c_{}(๐ฉ):๐ฉ\text{is a metric space},|๐ฉ|n\}.$$
Observe that if, for example, $``$ consists of all the subsets of Hilbert space (or $`L_1`$), then Bourgainโs embedding theorem implies that $`๐_n()=O(\mathrm{log}n)`$.
For $`K>0`$ we define the $`K`$-cotype (with exponent $`2`$) of a family of metric spaces $``$ as
$$q_{}^{(2)}(K)=\underset{}{sup}inf\{q(0,\mathrm{}]:\mathrm{\Gamma }_q^{(2)}()K\}.$$
Finally we let
$$q_{}^{(2)}=\underset{\mathrm{}>K>0}{inf}q_{}^{(2)}(K).$$
The following theorem settles positively the problem stated above:
###### Theorem 1.6.
Let $``$ be a family of metric spaces. Then the following conditions are equivalent:
1. There exists a finite metric space $``$ for which $`c_{}()>1`$.
2. $`q_{}^{(2)}<\mathrm{}`$.
3. There exists $`0<\alpha <\mathrm{}`$ such that $`๐_n()=\mathrm{\Omega }\left((\mathrm{log}n)^\alpha \right)`$.
#### 3) A quantitative version of Matou \ฬhskip-5.75ptseks BD Ramsey theorem.
In Matouลกek proved the following result, which he calls the Bounded Distortion (BD) Ramsey theorem. We refer to for motivation and background on these types of results.
Theorem 1.7 (Matouลกekโs BD Ramsey theorem). Let $`X`$ be a finite metric space and $`\epsilon >0`$, $`\gamma >1`$. Then there exists a metric space $`Y=Y(X,\epsilon ,\gamma )`$, such that for every metric space $`Z`$,
$$c_Z(Y)<\gamma c_Z(X)<1+\epsilon .$$
We obtain a new proof of Theorem 1.7, which is quantitative and concrete:
Theorem 1.8 (Quantitative version of Matouลกekโs BD Ramsey theorem). There exists a universal constant $`C`$ with the following properties. Let $`X`$ be an $`n`$-point metric space and $`\epsilon (0,1)`$, $`\gamma >1`$. Then for every integer $`N(C\gamma )^{2^{5A}}`$, where
$$A=\mathrm{max}\{\frac{4\mathrm{diam}(X)}{\epsilon \mathrm{min}_{xy}d_X(x,y)},n\},$$
if a metric space $`Z`$ satisfies $`c_Z(X)>1+\epsilon `$ then, $`c_Z\left(\left[N^5\right]_{\mathrm{}}^N\right)>\gamma `$.
We note that Matouลกekโs argument in uses Ramsey theory, and is nonconstructive (at best it can yield tower-type bounds on the size of $`Z`$, which are much worse than what the cotype-based approach gives).
#### 4) Uniform embeddings and Smirnovs problem.
Let $`(,d_{})`$ and $`(๐ฉ,d_๐ฉ)`$ be metric spaces. A mapping $`f:๐ฉ`$ is called a uniform embedding if $`f`$ is injective, and both $`f`$ and $`f^1`$ are uniformly continuous. There is a large body of work on the uniform classification of metric spaces โ we refer to the survey article , the book , and the references therein for background on this topic. In spite of this, several fundamental questions remain open. For example, it was not known for which values of $`0<p,q<\mathrm{}`$, $`L_p`$ embeds uniformly into $`L_q`$. As we will presently see, our results yield a complete characterization of these values of $`p,q`$.
In the late 1950โs Smirnov asked whether every separable metric space embeds uniformly into $`L_2`$ (see ). Smirnovโs problem was settled negatively by Enflo in . Following Enflo, we shall say that a metric space $``$ is a universal uniform embedding space if every separable metric space embeds uniformly into $``$. Since every separable metric space is isometric to a subset of $`C[0,1]`$, this is equivalent to asking whether $`C[0,1]`$ is uniformly homeomorphic to a subset of $``$ (the space $`C[0,1]`$ can be replaced here by $`c_0`$ due to Aharoniโs theorem ). Enflo proved that $`c_0`$ does not uniformly embed into Hilbert space. In , Aharoni, Maurey and Mityagin systematically studied metric spaces which are uniformly homeomorphic to a subset of Hilbert space, and obtained an elegant characterization of Banach spaces which are uniformly homeomorphic to a subset of $`L_2`$. In particular, the results of imply that for $`p>2`$, $`L_p`$ is not uniformly homeomorphic to a subset of $`L_2`$.
Here we prove that in the class of Banach spaces with nontrivial type, if $`Y`$ embeds uniformly into $`X`$, then $`Y`$ inherits the cotype of $`X`$. More precisely:
###### Theorem 1.9.
Let $`X`$ be a Banach space with nontrivial type. Assume that $`Y`$ is a Banach space which uniformly embeds into $`X`$. Then $`q_Yq_X`$.
As a corollary, we complete the characterization of the values of $`0<p`$,$`q<\mathrm{}`$ for which $`L_p`$ embeds uniformly into $`L_q`$:
###### Theorem 1.10.
For $`p,q>0`$, $`L_p`$ embeds uniformly into $`L_q`$ if and only if $`pq`$ or $`qp2`$.
We believe that the assumption that $`X`$ has nontrivial type in Theorem 1.9 can be removed โ in Section 8 we present a concrete problem which would imply this fact. If true, this would imply that cotype is preserved under uniform embeddings of Banach spaces. In particular, it would follow that a universal uniform embedding space cannot have nontrivial cotype, and thus by the Maurey-Pisier theorem it must contain $`\mathrm{}_{\mathrm{}}^n`$โs with distortion uniformly bounded in $`n`$.
#### 5) Coarse embeddings.
Let $`(,d_{})`$ and $`(๐ฉ,d_๐ฉ)`$ be metric spaces. A mapping $`f:๐ฉ`$ is called a coarse embedding if there exists two nondecreasing functions $`\alpha ,\beta :[0,\mathrm{})[0,\mathrm{})`$ such that $`lim_t\mathrm{}\alpha (t)=\mathrm{}`$, and for every $`x,y`$,
$$\alpha (d_{}(x,y))d_๐ฉ(f(x),f(y))\beta (d_{}(x,y)).$$
This (seemingly weak) notion of embedding was introduced by Gromov (see ), and has several important geometric applications. In particular, Yu obtained a striking connection between the Novikov and Baum-Connes conjectures and coarse embeddings into Hilbert spaces. In Kasparov and Yu generalized this to coarse embeddings into arbitrary uniformly convex Banach spaces. It was unclear, however, whether this is indeed a strict generalization, i.e. whether or not the existence of a coarse embedding into a uniformly convex Banach space implies the existence of a coarse embedding into a Hilbert space. This was resolved by Johnson and Randrianarivony in , who proved that for $`p>2`$, $`L_p`$ does not coarsely embed into $`L_2`$. In , Randrianarivony proceeded to obtain a characterization of Banach spaces which embed coarsely into $`L_2`$, in the spirit of the result of Aharoni, Maurey and Mityagin . There are very few known methods of proving coarse nonembeddability results. Apart from the papers , quoted above, we refer to , , for results of this type. Here we use metric cotype to prove the following coarse variants of Theorem 1.9 and Theorem 1.10, which generalize, in particular, the theorem of Johnson and Randrianarivony.
###### Theorem 1.11.
Let $`X`$ be a Banach space with nontrivial type. Assume that $`Y`$ is a Banach space which coarsely embeds into $`X`$. Then $`q_Yq_X`$. In particular, for $`p,q>0`$, $`L_p`$ embeds coarsely into $`L_q`$ if and only if $`pq`$ or $`qp2`$.
#### 6) Bi-Lipschitz embeddings of the integer lattice.
Bi-Lipschitz embeddings of the integer lattice $`[m]_p^n`$ were investigated by Bourgain in and by the present authors in where it was shown that if $`2p<\mathrm{}`$ and $`Y`$ is a Banach space which admits an equivalent norm whose modulus of uniform convexity has power type $`2`$, then
(8)
$$c_Y\left([m]_p^n\right)=\mathrm{\Theta }\left(\mathrm{min}\{n^{\frac{1}{2}\frac{1}{p}},m^{1\frac{2}{p}}\}\right).$$
The implied constants in the above asymptotic equivalence depend on $`p`$ and on the $`2`$-convexity constant of $`Y`$. Moreover, it was shown in that
$$c_Y([m]_{\mathrm{}}^n)=\mathrm{\Omega }\left(\mathrm{min}\{\sqrt{\frac{n}{\mathrm{log}n}},\frac{m}{\sqrt{\mathrm{log}m}}\}\right).$$
It was conjectured in that the logarithmic terms above are unnecessary. Using our results on metric cotype we settle this conjecture positively, by proving the following general theorem:
###### Theorem 1.12.
Let $`Y`$ be a Banach space with nontrivial type which has cotype $`q`$. Then
$$c_Y([m]_{\mathrm{}}^n)=\mathrm{\Omega }\left(\mathrm{min}\{n^{1/q},m\}\right).$$
Similarly, our methods imply that (8) holds true for any Banach space $`Y`$ with nontrivial type and cotype $`2`$ (note that these conditions are strictly weaker than being $`2`$-convex, as shown e.g. in ). Moreover, it is possible to generalize the lower bound in (8) to Banach spaces with nontrivial type, and cotype $`2qp`$, in which case the lower bound becomes $`\mathrm{min}\{n^{\frac{1}{q}\frac{1}{p}},m^{1\frac{q}{p}}\}`$.
#### 7) Quadratic inequalities on the cut-cone.
An intriguing aspect of Theorem 1.2 is that $`L_1`$ has metric cotype $`2`$. Thus, we obtain a nontrivial inequality on $`L_1`$ which involves distances squared. To the best of our knowledge, all the known nonembeddability results for $`L_1`$ are based on Poincarรฉ type inequalities in which distances are raised to the power $`1`$. Clearly, any such inequality reduces to an inequality on the real line. Equivalently, by the cut-cone representation of $`L_1`$ metrics (see ) it is enough to prove any such inequality for cut metrics, which are particularly simple. Theorem 1.2 seems to be the first truly โinfinite dimensionalโ metric inequality in $`L_1`$, in the sense that its nonlinearity does not allow a straightforward reduction to the one-dimensional case. We believe that understanding such inequalities on $`L_1`$ deserves further scrutiny, especially as they hint at certain nontrivial (and nonlinear) interactions between cuts.
## 2. Preliminaries and notation
We start by setting notation and conventions. Consider the standard $`\mathrm{}_{\mathrm{}}`$ Cayley graph on $`_m^n`$, namely $`x,y_m^n`$ are joined by an edge if and only if they are distinct and $`xy\{1,0,1\}^n`$. This induces a shortest-path metric on $`_m^n`$ which we denote by $`d_{_m^n}(,)`$. Equivalently, the metric space $`(_m^n,d_{_m^n})`$ is precisely the quotient $`(^n,_{\mathrm{}})/(m)^n`$ (for background on quotient metrics see , ). The ball of radius $`r`$ around $`x_m^n`$ will be denoted $`B_{_m^n}(x,r)`$. We denote by $`\mu `$ the normalized counting measure on $`_m^n`$ (which is clearly the Haar measure on this group). We also denote by $`\sigma `$ the normalized counting measure on $`\{1,0,1\}^n`$. In what follows, whenever we average over uniformly chosen signs $`\epsilon \{1,1\}^n`$ we use the probabilistic notation $`๐ผ_\epsilon `$ (in this sense we break from the notation used in the introduction, for the sake of clarity of the ensuing arguments).
In what follows all Banach spaces are assumed to be over the complex numbers $``$. All of our results hold for real Banach spaces as well, by a straightforward complexification argument.
Given a Banach space $`X`$ and $`p,q[1,\mathrm{})`$ we denote by $`C_q^{(p)}(X)`$ the infimum over all constants $`C>0`$ such that for every integer $`n`$ and every $`x_1,\mathrm{},x_nX`$,
(9) $`\left(๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j_X^p\right)^{1/p}{\displaystyle \frac{1}{C}}({\displaystyle \underset{j=1}{\overset{n}{}}}x_j_X^q)^{1/q}.`$
Thus, by our previous notation, $`C_q^{(q)}(X)=C_q(X)`$. Kahaneโs inequality says that for $`1p,q<\mathrm{}`$ there exists a constant $`1A_{pq}<\mathrm{}`$ such that for every Banach space $`X`$, every integer $`n`$, and every $`x_1,\mathrm{},x_nX`$,
(10) $`\left(๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j_X^p\right)^{1/p}A_{pq}\left(๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j_X^q\right)^{1/q}.`$
Where clearly $`A_{pq}=1`$ if $`pq`$, and for every $`1q<p<\mathrm{}`$, $`A_{pq}=O\left(\sqrt{p}\right)`$ (see ). It follows in particular from (10) that if $`X`$ has cotype $`q`$ then for every $`p[1,\mathrm{})`$, $`C_q^{(p)}(X)=O_{p,q}(C_q(X))`$, where the implied constant may depend on $`p`$ and $`q`$.
Given $`A\{1,\mathrm{},n\}`$, we consider the Walsh functions $`W_A:\{1,1\}^n`$, defined as
$$W_A(\epsilon _1,\mathrm{},\epsilon _m)=\underset{jA}{}\epsilon _j.$$
Every $`f:\{1,1\}^nX`$ can be written as
$$f(\epsilon _1,\mathrm{},\epsilon _n)=\underset{A\{1,\mathrm{},n\}}{}\widehat{f}(A)W_A(\epsilon _1,\mathrm{},\epsilon _n),$$
where $`\widehat{f}(A)X`$ are given by
$$\widehat{f}(A)=๐ผ_\epsilon \left(f(\epsilon )W_A(\epsilon )\right).$$
The Rademacher projection of $`f`$ is defined by
$$\mathrm{๐๐๐}(f)=\underset{j=1}{\overset{n}{}}\widehat{f}(A)W_{\{j\}}.$$
The $`K`$-convexity constant of $`X`$, denoted $`K(X)`$, is the smallest constant $`K`$ such that for every $`n`$ and every $`f:\{1,1\}^nX`$,
$$๐ผ_\epsilon \mathrm{๐๐๐}(f)(\epsilon )_X^2K^2๐ผ_\epsilon f(\epsilon )_X^2.$$
In other words,
$$K(X)=\underset{n}{sup}\mathrm{๐๐๐}_{L_2(\{1,1\}^n,X)L_2(\{1,1\}^n,X)}.$$
$`X`$ is said to be $`K`$-convex if $`K(X)<\mathrm{}`$. More generally, for $`p1`$ we define
$$K_p(X)=\underset{n}{sup}\mathrm{๐๐๐}_{L_p(\{1,1\}^n,X)L_p(\{1,1\}^n,X)}.$$
It is a well known consequence of Kahaneโs inequality and duality that for every $`p>1`$,
$$K_p(X)O\left(\frac{p}{\sqrt{p1}}\right)K(X).$$
The following deep theorem was proved by Pisier in :
###### Theorem 2.1 (Pisierโs $`K`$-convexity theorem ).
Let $`X`$ be a Banach space. Then
$$q_X>1K(X)<\mathrm{}.$$
Next, we recall some facts concerning Fourier analysis on the group $`_m^n`$. Given $`k=(k_1,\mathrm{},k_n)_m^n`$ we consider the Walsh function $`W_k:_m^n`$:
$$W_k(x)=\mathrm{exp}\left(\frac{2\pi i}{m}\underset{j=1}{\overset{m}{}}k_jx_j\right).$$
Then, for any Banach space $`X`$, any $`f:Z_m^nX`$ can be decomposed as follows:
$$f(x)=\underset{k_m^n}{}W_k(x)\widehat{f}(k),$$
where
$$\widehat{f}(k)=_{_m^n}f(y)\overline{W_k(y)}๐\mu (y)X.$$
If $`X`$ is a Hilbert space then Parsevalโs identity becomes:
$$_{Z_m^n}f(x)_X^2๐\mu (x)=\underset{k_m^n}{}\widehat{f}(k)_X^2.$$
### 2.1. Definitions and basic facts related to metric cotype
###### Definition 2.2.
Given $`1pq`$, an integer $`n`$ and an even integer $`m`$, let $`\mathrm{\Gamma }_q^{(p)}(;n,m)`$ be the infimum over all $`\mathrm{\Gamma }>0`$ such that for every $`f:_m^n`$,
(11)
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}d_{}(f\left(x+\frac{m}{2}e_j\right),f(x))^p๐\mu (x)\hfill \\ \hfill \mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}_{\{1,0,1\}^n}_{_m^n}d_{}(f\left(x+\epsilon \right),f(x))^p๐\mu (x)๐\sigma (\epsilon ).\end{array}$$
When $`p=q`$ we write $`\mathrm{\Gamma }_q(;n,m):=\mathrm{\Gamma }_q^{(q)}(;n,m)`$ . With this notation,
$$\mathrm{\Gamma }_q^{(p)}()=\underset{n}{sup}\underset{m2}{inf}\mathrm{\Gamma }_q^{(p)}(;n,m).$$
We also denote by $`m_q^{(p)}(;n,\mathrm{\Gamma })`$ the smallest even integer $`m`$ for which (11) holds. As usual, when $`p=q`$ we write $`m_q(;n,\mathrm{\Gamma }):=m_q^{(q)}(;n,\mathrm{\Gamma })`$.
The following lemma shows that for nontrivial metric spaces $``$,$`m_q(;n,\mathrm{\Gamma })`$ must be large.
###### Lemma 2.3.
Let $`(,d_{})`$ be a metric space which contains at least two points. Then for every integer $`n`$, every $`\mathrm{\Gamma }>0`$, and every $`p,q>0`$,
$$m_q^{(p)}(;n,\mathrm{\Gamma })\frac{n^{1/q}}{\mathrm{\Gamma }}.$$
###### Proof.
Fix $`u,v`$, $`uv`$, and without loss of generality normalize the metric so that $`d_{}(u,v)=1`$. Denote $`m=m_q^{(p)}(;n,\mathrm{\Gamma })`$. Let $`f:_m^n`$ be the random mapping such that for every $`x_m^n`$, $`\mathrm{Pr}[f(x)=u]=\mathrm{Pr}[f(x)=v]=\frac{1}{2}`$, and $`\{f(x)\}_{x_m^n}`$ are independent random variables. Then for every distinct $`x,y_m^n`$, $`๐ผ\left[d_{}(f(x),f(y))^p\right]=\frac{1}{2}`$. Thus, the required result follows by applying (11) to $`f`$ and taking expectation. โ
###### Lemma 2.4.
For every two integers $`n,k`$, and every even integer $`m`$,
$$\mathrm{\Gamma }_q^{(p)}(;n,km)\mathrm{\Gamma }_q^{(p)}(;n,m).$$
###### Proof.
Fix $`f:_{km}^n`$. For every $`y_k^n`$ define $`f_y:_m^n`$ by
$$f_y(x)=f(kx+y).$$
Fix $`\mathrm{\Gamma }>\mathrm{\Gamma }_q^{(p)}(;n,m)`$. Applying the definition of $`\mathrm{\Gamma }_q^{(p)}(;n,m)`$ to $`f_y`$, we get that
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}d_{}(f\left(kx+\frac{km}{2}e_j+y\right),f(kx+y))^p๐\mu _{_m^n}(x)\hfill \\ \hfill \mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}_{\{1,0,1\}^n}_{_m^n}d_{}(f\left(kx+k\epsilon +y\right),f(kx+y))^p๐\mu _{_m^n}(x)๐\sigma (\epsilon ).\end{array}$$
Integrating this inequality with respect to $`y_k^n`$ we see that
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_{km}^n}}d_{}(f\left(z+{\displaystyle \frac{km}{2}}e_j\right),f\left(z\right))^p๐\mu _{_{km}^n}\left(z\right)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_k^n}}{\displaystyle _{_m^n}}d_{}(f\left(kx+{\displaystyle \frac{km}{2}}e_j+y\right),f\left(kx+y\right))^p๐\mu _{_m^n}\left(x\right)๐\mu _{_k^n}\left(y\right)`$
$``$ $`\mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_k^n}}{\displaystyle _{_m^n}}d_{}(f\left(kx+k\epsilon +y\right),f\left(kx+y\right))^p๐\mu _{_m^n}\left(x\right)๐\mu _{_k^n}\left(y\right)๐\sigma \left(\epsilon \right)`$
$`=`$ $`\mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_{km}^n}}d_{}(f\left(z+k\epsilon \right),f\left(z\right))^p๐\mu _{_{km}^n}\left(z\right)๐\sigma \left(\epsilon \right)`$
$``$ $`\mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_{km}^n}}k^{p1}{\displaystyle \underset{s=1}{\overset{k}{}}}d_{}(f\left(z+s\epsilon \right),f\left(z+\left(s1\right)\epsilon \right))^pd\mu _{_{km}^n}\left(z\right)d\sigma \left(\epsilon \right)`$
$`=`$ $`\mathrm{\Gamma }^p\left(km\right)^pn^{1\frac{p}{q}}{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_{km}^n}}d_{}(f\left(z+\epsilon \right),f\left(z\right))^p๐\mu _{_{km}^n}\left(z\right)๐\sigma \left(\epsilon \right).\mathit{}`$
###### Lemma 2.5.
Let $`k,n`$ be integers such that $`kn`$, and let $`m`$ be an even integer. Then
$$\mathrm{\Gamma }_q^{(p)}(;k,m)\left(\frac{n}{k}\right)^{1\frac{p}{q}}\mathrm{\Gamma }_q^{(p)}(;n,m).$$
###### Proof.
Given an $`f:_m^k`$, we define an $``$-valued function on $`_m^n_m^k\times _m^{nk}`$ by $`g(x,y)=f(x)`$. Applying the definition $`\mathrm{\Gamma }_q^{(p)}(;n,m)`$ to $`g`$ yields the required inequality. โ
We end this section by recording some general inequalities which will be used in the ensuing arguments. In what follows $`(,d_{})`$ is an arbitrary metric space.
###### Lemma 2.6.
For every $`f:_m^n`$,
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}d_{}(f(x+e_j),f(x))^p๐\mu (x)\hfill \\ \hfill 32^{p1}n_{\{1,0,1\}^n}_{_m^n}d_{}(f(x+\epsilon ),f(x))^p๐\mu (x)๐\sigma (\epsilon ).\end{array}$$
###### Proof.
For every $`x_m^n`$ and $`\epsilon \{1,0,1\}^n`$,
$$d_{}(f(x+e_j),f(x))^p2^{p1}d_{}(f(x+e_j),f(x+\epsilon ))^p+2^{p1}d_{}(f(x+\epsilon ),f(x))^p.$$
Thus
$`{\displaystyle \frac{2}{3}}{\displaystyle _{_m^n}}`$ $`d_{}(f(x+e_j),f(x))^pd\mu (x)`$
$`=\sigma (\{\epsilon \{1,0,1\}^n:\epsilon _j1\}){\displaystyle _{_m^n}}d_{}(f(x+e_j),f(x))^p๐\mu (x)`$
$`2^{p1}{\displaystyle _{\{\epsilon \{1,0,1\}^n:\epsilon _j1\}}}{\displaystyle _{_m^n}}(d_{}(f(x+e_j),f(x+\epsilon ))^p`$
$`+d_{}(f(x+\epsilon ),f(x))^p)d\mu (x)d\sigma (\epsilon )`$
$`=2^{p1}{\displaystyle _{\{\epsilon \{1,0,1\}^n:\epsilon _j1\}}}{\displaystyle _{_m^n}}d_{}(f(y+\epsilon ),f(y))^p๐\mu (y)๐\sigma (\epsilon )`$
$`+2^{p1}{\displaystyle _{\{\epsilon \{1,0,1\}^n:\epsilon _j1\}}}{\displaystyle _{_m^n}}d_{}(f(x+\epsilon ),f(x))^p๐\mu (x)๐\sigma (\epsilon )`$
$`2^p{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}d_{}(f(x+\epsilon ),f(x))^p๐\mu (x)๐\sigma (\epsilon ).`$
Summing over $`j=1,\mathrm{},n`$ yields the required result. โ
###### Lemma 2.7.
Let $`(,d_{})`$ be a metric space. Assume that for an integer $`n`$ and an even integer $`m`$, we have that for every $`\mathrm{}n`$, and every $`f:_m^{\mathrm{}}`$,
$$\begin{array}{c}\underset{j=1}{\overset{\mathrm{}}{}}__m^{\mathrm{}}d_{}(f\left(x+\frac{m}{2}e_j\right),f\left(x\right))^p๐\mu (x)\hfill \\ \hfill C^pm^pn^{1\frac{p}{q}}(๐ผ_\epsilon __m^{\mathrm{}}d_{}(f(x+\epsilon ),f(x))^pd\mu (x)\\ \hfill +\frac{1}{\mathrm{}}\underset{j=1}{\overset{\mathrm{}}{}}__m^{\mathrm{}}d_{}(f(x+e_j),f(x))^pd\mu (x)).\end{array}$$
Then
$$\mathrm{\Gamma }_q^{(p)}(;n,m)5C.$$
###### Proof.
Fix $`f:_m^n`$ and $`\mathrm{}A\{1,\mathrm{},n\}`$. Our assumption implies that
$$\begin{array}{c}\underset{jA}{}_{_m^n}d_{}(f\left(x+\frac{m}{2}e_j\right),f\left(x\right))^p๐\mu (x)\hfill \\ \hfill C^pm^pn^{1\frac{p}{q}}(๐ผ_\epsilon _{_m^n}d_{}(f(x+\underset{jA}{}\epsilon _je_j),f(x))^pd\mu (x)\\ \hfill +\frac{1}{|A|}\underset{jA}{}_{_m^n}d_{}(f(x+e_j),f(x))^pd\mu (x)).\end{array}$$
Multiplying this inequality by $`\frac{2^{|A|}}{3^n}`$, and summing over all $`\mathrm{}A\{1,\mathrm{},n\}`$, we see that
(12)
$`{\displaystyle \frac{2}{3}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(f\left(x+{\displaystyle \frac{m}{2}}e_j\right),f\left(x\right))^p๐\mu (x)`$
$`={\displaystyle \underset{\mathrm{}A\{1,\mathrm{},n\}}{}}{\displaystyle \frac{2^{|A|}}{3^n}}{\displaystyle \underset{jA}{}}{\displaystyle _{_m^n}}d_{}(f\left(x+{\displaystyle \frac{m}{2}}e_j\right),f\left(x\right))^p๐\mu (x)`$
$`C^pm^pn^{1\frac{p}{q}}({\displaystyle \underset{\mathrm{}A\{1,\mathrm{},n\}}{}}{\displaystyle \frac{2^{|A|}}{3^n}}๐ผ_\epsilon {\displaystyle _{_m^n}}d_{}(f(x+{\displaystyle \underset{jA}{}}\epsilon _je_j),f(x))^pd\mu (x)`$
$`+{\displaystyle \underset{\mathrm{}A\{1,\mathrm{},n\}}{}}{\displaystyle \frac{2^{|A|}}{|A|3^n}}{\displaystyle \underset{jA}{}}{\displaystyle _{_m^n}}d_{}(f(x+e_j),f(x))^pd\mu (x))`$
(13) $`C^pm^pn^{1\frac{p}{q}}({\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}d_{}(f(x+\delta ),f(x))^pd\mu (x)d\sigma (\delta )`$
$`+{\displaystyle \frac{1}{n}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(f(x+e_j),f(x))^pd\mu (x))`$
$`C^pm^pn^{1\frac{p}{q}}\left(3^p+1\right){\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}d_{}(f\left(x+\delta \right),f(x))^p๐\mu (x)๐\sigma (\delta ),`$
where we used the fact that in (12), the coefficient of $`d_{}(f(x+e_j),f(x))^p`$ equals $`_{k=1}^n\frac{2^k}{k3^n}\left(\genfrac{}{}{0pt}{}{n1}{k1}\right)\frac{1}{n}`$, and in (13) we used Lemma 2.6. โ
## 3. Warmup: the case of Hilbert space
The fact that Hilbert spaces have metric cotype $`2`$ is particularly simple to prove. This is contained in the following proposition.
###### Proposition 3.1.
Let $`H`$ be a Hilbert space. Then for every integer $`n`$, and every integer $`m\frac{2}{3}\pi \sqrt{n}`$ which is divisible by $`4`$,
$$\mathrm{\Gamma }_2(H;n,m)\frac{\sqrt{6}}{\pi }.$$
###### Proof.
Fix $`f:Z_m^nH`$ and decompose it into Fourier coefficients:
$$f(x)=\underset{k_m^n}{}W_k(x)\widehat{f}(k).$$
For every $`j=1,2,\mathrm{},n`$ we have that
$$f\left(x+\frac{m}{2}e_j\right)f(x)=\underset{k_m^n}{}W_k(x)\left(e^{\pi ik_j}1\right)\widehat{f}(k).$$
Thus
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f(x)_H^2๐\mu (x)`$
$`=`$ $`{\displaystyle \underset{k_m^n}{}}\left({\displaystyle \underset{j=1}{\overset{n}{}}}\left|e^{\pi ik_j}1\right|^2\right)\widehat{f}(k)_H^2=4{\displaystyle \underset{k_m^n}{}}|\{j:k_j1\text{ mod}\mathrm{\hspace{0.17em}2}\}|\widehat{f}(k)_H^2.`$
Additionally, for every $`\epsilon \{1,0,1\}^n`$,
$$f(x+\epsilon )f(x)=\underset{k_m^n}{}W_k(x)(W_k(\epsilon )1)\widehat{f}(k).$$
Thus
$$\begin{array}{c}_{\{1,0,1\}^n}_{_m^n}f(x+\epsilon )f(x)_H^2๐\mu (x)๐\sigma (\epsilon )\hfill \\ \hfill =\underset{k_m^n}{}\left(_{\{1,0,1\}^n}\left|W_k(\epsilon )1\right|^2๐\sigma (\epsilon )\right)\widehat{f}(k)_H^2.\end{array}$$
Observe that
$`{\displaystyle _{\{1,0,1\}^n}}\left|W_k(\epsilon )1\right|^2๐\sigma (\epsilon )`$ $`=`$ $`{\displaystyle _{\{1,0,1\}^n}}|\mathrm{exp}({\displaystyle \frac{2\pi i}{m}}{\displaystyle \underset{j=1}{\overset{m}{}}}k_j\epsilon _j)1|^2d\sigma (\epsilon )`$
$`=`$ $`22\mathrm{Re}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{\{1,0,1\}^n}}\mathrm{exp}\left({\displaystyle \frac{2\pi i}{m}}k_j\epsilon _j\right)๐\sigma (\epsilon )`$
$`=`$ $`22{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{1+2\mathrm{cos}\left(\frac{2\pi }{m}k_j\right)}{3}}`$
$``$ $`22{\displaystyle \underset{j:k_j1mod2}{}}{\displaystyle \frac{1+2\left|\mathrm{cos}\left(\frac{2\pi }{m}k_j\right)\right|}{3}}.`$
Note that if $`m`$ is divisible by $`4`$ and $`\mathrm{}\{0,\mathrm{},m1\}`$ is an odd integer, then
$$\left|\mathrm{cos}\left(\frac{2\pi \mathrm{}}{m}\right)\right|\left|\mathrm{cos}\left(\frac{2\pi }{m}\right)\right|1\frac{\pi ^2}{m^2}.$$
Hence
$`{\displaystyle _{\{1,0,1\}^n}}\left|W_k(\epsilon )1\right|^2๐\sigma (\epsilon )`$ $``$ $`2\left(1\left(1{\displaystyle \frac{2\pi ^2}{3m^2}}\right)^{|\{j:k_j1mod2\}|}\right)`$
$``$ $`2\left(1e^{\frac{2|\{j:k_j1mod2\}|\pi ^2}{3m^2}}\right)`$
$``$ $`|\{j:k_j1mod2\}|{\displaystyle \frac{2\pi ^2}{3m^2}},`$
provided that $`m\frac{2}{3}\pi \sqrt{n}`$. โ
## 4. $`K`$-convex spaces
In this section we prove the โhard directionโ of Theorem 1.2 and Theorem 1.4 when $`X`$ is a $`K`$-convex Banach space; namely, we show that in this case Rademacher cotype $`q`$ implies metric cotype $`q`$. There are two reasons why we single out this case before passing to the proofs of these theorems in full generality. First of all, the proof for $`K`$-convex spaces is different and simpler than the general case. More importantly, in the case of $`K`$-convex spaces we are able to obtain optimal bounds on the value of $`m`$ in Definition 1.1 and Definition 1.3. Namely, we show that if $`X`$ is a $`K`$-convex Banach space of cotype $`q`$, then for every $`1pq`$, $`m_q^{(p)}(X;n,\mathrm{\Gamma })=O(n^{1/q})`$, for some $`\mathrm{\Gamma }=\mathrm{\Gamma }(X)`$. This is best possible due to Lemma 2.3. In the case of general Banach spaces we obtain worse bounds, and this is why we have the restriction that $`X`$ is $`K`$-convex in Theorem 1.9 and Theorem 1.11. This issue is taken up again in Section 8.
###### Theorem 4.1.
Let $`X`$ be a $`K`$-convex Banach space with cotype $`q`$. Then for every integer $`n`$ and every integer $`m`$ which is divisible by $`4`$,
$$m\frac{2n^{1/q}}{C_q^{(p)}(X)K_p(X)}\mathrm{\Gamma }_q^{(p)}(X;n,m)15C_q^{(p)}(X)K_p(X).$$
###### Proof.
For $`f:_m^nX`$ we define the following operators:
$`\stackrel{~}{}_jf(x)`$ $`=`$ $`f(x+e_j)f(xe_j),`$
$`_jf(x)`$ $`=`$ $`๐ผ_\epsilon f\left(x+{\displaystyle \underset{\mathrm{}j}{}}\epsilon _{\mathrm{}}e_{\mathrm{}}\right),`$
and for $`\epsilon \{1,0,1\}^n`$,
$$_\epsilon f(x)=f(x+\epsilon )f(x).$$
These operators operate diagonally on the Walsh basis $`\{W_k\}_{k_m^n}`$ as follows:
(14) $`\stackrel{~}{}_jW_k`$ $`=`$ $`\left(W_k(e_j)W_k(e_j)\right)W_k=2\mathrm{sin}\left({\displaystyle \frac{2\pi ik_j}{m}}\right)W_k,`$
(15) $`_jW_k`$ $`=`$ $`\left(๐ผ_\epsilon {\displaystyle \underset{\mathrm{}j}{}}e^{\frac{2\pi i\epsilon _{\mathrm{}}k_{\mathrm{}}}{m}}\right)W_k=\left({\displaystyle \underset{\mathrm{}j}{}}\mathrm{cos}\left({\displaystyle \frac{2\pi k_{\mathrm{}}}{m}}\right)\right)W_k,`$
and for $`\epsilon \{1,1\}^n`$,
$`_\epsilon W_k`$ $`=`$ $`\left(W(\epsilon )1\right)W_k`$
$`=`$ $`\left({\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi i\epsilon _jk_j}{m}}1\right)W_k`$
$`=`$ $`\left({\displaystyle \underset{j=1}{\overset{n}{}}}\left(\mathrm{cos}\left({\displaystyle \frac{2\pi \epsilon _jk_j}{m}}\right)+i\mathrm{sin}\left({\displaystyle \frac{2\pi \epsilon _jk_j}{m}}\right)\right)1\right)W_k`$
$`=`$ $`\left({\displaystyle \underset{j=1}{\overset{n}{}}}\left(\mathrm{cos}\left({\displaystyle \frac{2\pi k_j}{m}}\right)+i\epsilon _j\mathrm{sin}\left({\displaystyle \frac{2\pi k_j}{m}}\right)\right)1\right)W_k.`$
The last step was a crucial observation, using the fact that $`\epsilon _j\{1,1\}`$. Thinking of $`_\epsilon W_k`$ as a function of $`\epsilon \{1,1\}^n`$, equations (14), (15) and (4) imply that
$`\mathrm{๐๐๐}(_\epsilon W_k)`$ $`=`$ $`i\left({\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\mathrm{sin}\left({\displaystyle \frac{2\pi k_j}{m}}\right){\displaystyle \underset{\mathrm{}j}{}}\mathrm{cos}\left({\displaystyle \frac{2\pi k_{\mathrm{}}}{m}}\right)\right)W_k`$
$`=`$ $`{\displaystyle \frac{i}{2}}\left({\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\stackrel{~}{}_j_j\right)W_k.`$
Thus for every $`x_m^n`$ and $`f:_m^nX`$,
$$\mathrm{๐๐๐}(_\epsilon f(x))=\frac{i}{2}\left(\underset{j=1}{\overset{n}{}}\epsilon _j\stackrel{~}{}_j_j\right)f(x).$$
It follows that
(17) $`{\displaystyle _{_m^n}}๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j[_jf(x+e_j)_jf(xe_j)]_X^pd\mu (x)`$
$`={\displaystyle _{_m^n}}๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\stackrel{~}{}_j_jf(x)_X^pd\mu (x)`$
$`={\displaystyle _{_m^n}}๐ผ_\epsilon \mathrm{๐๐๐}(_\epsilon f(x))_X^pd\mu (x)`$
$`K_p(X)^p{\displaystyle _{_m^n}}๐ผ_\epsilon _\epsilon f(x)_X^pd\mu (x).`$
By (17) and the definition of $`C_q^{(p)}(X)`$, for every $`C>C_q^{(p)}(X)`$ we have that
(18) $`[K_p(X)C]^p๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`C^p๐ผ_\epsilon {\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j[_jf(x+e_j)_jf(xe_j)]_X^p๐\mu (x)`$
$`{\displaystyle _{_m^n}}\left({\displaystyle \underset{j=1}{\overset{n}{}}}_jf(x+e_j)_jf(xe_j)_X^q\right)^{p/q}๐\mu (x)`$
$`{\displaystyle \frac{1}{n^{1p/q}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}_jf(x+e_j)_jf(xe_j)_X^p๐\mu (x).`$
Now, for $`j\{1,\mathrm{},n\}`$,
$`{\displaystyle _{_m^n}}_jf\left(x+{\displaystyle \frac{m}{2}}e_j\right)_jf\left(x\right)_X^p๐\mu (x)`$
$`\left({\displaystyle \frac{m}{4}}\right)^{p1}{\displaystyle \underset{s=1}{\overset{m/4}{}}}{\displaystyle _{_m^n}}_jf\left(x+2se_j\right)_jf\left(x+2(s1)e_j\right)_X^p๐\mu (x)`$
$`=\left({\displaystyle \frac{m}{4}}\right)^p{\displaystyle _{_m^n}}_jf(x+e_j)_jf(xe_j)_X^p๐\mu (x).`$
Plugging (4) into (18) we get
$`\left({\displaystyle \frac{m}{4}}\right)^pn^{1\frac{p}{q}}[K_p(X)C]^p๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}_jf\left(x+{\displaystyle \frac{m}{2}}e_j\right)_jf\left(x\right)_X^p๐\mu (x)`$
$`{\displaystyle \frac{1}{3^{p1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`2{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}_jf\left(x\right)f\left(x\right)_X^p๐\mu (x)`$
$`={\displaystyle \frac{1}{3^{p1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`2{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}๐ผ_\epsilon \left(f\left(x+{\displaystyle \underset{\mathrm{}j}{}}\epsilon _{\mathrm{}}e_{\mathrm{}}\right)f\left(x\right)\right)_X^p๐\mu (x)`$
$`{\displaystyle \frac{1}{3^{p1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`2{\displaystyle \underset{j=1}{\overset{n}{}}}๐ผ_\epsilon {\displaystyle _{_m^n}}f\left(x+{\displaystyle \underset{\mathrm{}j}{}}\epsilon _{\mathrm{}}e_{\mathrm{}}\right)f\left(x\right)_X^p๐\mu (x)`$
$`{\displaystyle \frac{1}{3^{p1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`2^pn๐ผ_\epsilon {\displaystyle _{_m^n}}f\left(x+\epsilon \right)f\left(x\right)_X^p๐\mu (x)`$
$`2^p{\displaystyle \underset{j=1}{\overset{n}{}}}๐ผ_\epsilon {\displaystyle _{_m^n}}f\left(x+\epsilon _je_j\right)f\left(x\right)_X^p๐\mu (x).`$
Thus, the required result follows from Lemma 2.7. โ
The above argument actually gives the following generalization of Theorem 4.1, which holds for products of arbitrary compact Abelian groups.
###### Theorem 4.2.
Let $`G_1,\mathrm{},G_n`$ be compact Abelian groups, $`(g_1,\mathrm{},g_n)G_1\times \mathrm{}\times G_n`$, and let $`X`$ be a $`K`$-convex Banach space. Then for every integer $`k`$ and every $`f:G_1\times \mathrm{}\times G_nX`$,
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{G_1\times \mathrm{}\times G_n}f\left(x+2kg_je_j\right)f\left(x\right)_X^pd\left(\mu _{G_1}\mathrm{}\mu _{G_n}\right)\left(x\right)\hfill \\ \hfill C^p_{\{1,0,1\}^n}_{G_1\times \mathrm{}\times G_n}f\left(x+\underset{j=1}{\overset{n}{}}\epsilon _jg_je_j\right)f\left(x\right)_X^pd\left(\mu _{G_1}\mathrm{}\mu _{G_n}\right)\left(x\right)๐\sigma \left(\epsilon \right),\end{array}$$
where
$$C5\mathrm{max}\{C_q^{(p)}(X)K_p(X)kn^{\frac{1}{p}\frac{1}{q}},n^{\frac{1}{p}}\}.$$
Here $`\mu _G`$ denotes the normalized Haar measure on a compact Abelian group $`G`$. We refer the interested reader to the book , which contains the necessary background required to generalize the proof of Theorem 4.1 to this setting.
## 5. The equivalence of Rademacher cotype and metric cotype
We start by establishing the easy direction in Theorem 1.2 and Theorem 1.4, i.e. that metric cotype implies Rademacher cotype.
### 5.1. Metric cotype implies Rademacher cotype
Let $`X`$ be a Banach space and assume that $`\mathrm{\Gamma }_q^{(p)}(X)<\mathrm{}`$ for some $`1pq`$. Fix $`\mathrm{\Gamma }>\mathrm{\Gamma }_q^{(p)}(X)`$, $`v_1,\mathrm{},v_nX`$, and let $`m`$ be an even integer. Define $`f:_m^nX`$ by
$$f(x_1,\mathrm{},x_n)=\underset{j=1}{\overset{n}{}}e^{\frac{2\pi ix_j}{m}}v_j.$$
Then
(20) $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f(x)_X^p๐\mu (x)=2^p{\displaystyle \underset{j=1}{\overset{n}{}}}v_j_X^p,`$
and
(21)
$$\begin{array}{c}_{\{1,0,1\}^n}_{_m^n}f\left(x+\delta \right)f(x)_X^p๐\mu (x)๐\sigma (\delta )\hfill \\ \hfill =_{\{1,0,1\}^n}_{_m^n}\underset{j=1}{\overset{n}{}}e^{\frac{2\pi ix_j}{m}}\left(e^{\frac{2\pi i\delta _j}{m}}1\right)v_j_X^p๐\mu (x)๐\sigma (\delta ).\end{array}$$
We recall the contraction principle (see ), which states that for every $`a_1,\mathrm{},a_n`$,
$$๐ผ_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _ja_jv_j_X^p(\underset{1jn}{\mathrm{max}}|a_j|)^p๐ผ_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _jv_j_X^p.$$
Observe that for every $`\epsilon =(\epsilon _1,\mathrm{},\epsilon _n)\{1,1\}^n`$,
$`{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi ix_j}{m}}\left(e^{\frac{2\pi i\delta _j}{m}}1\right)v_j_X^p๐\mu (x)๐\sigma (\delta )`$
$`=`$ $`{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi i}{m}\left(x_j+\frac{m(1\epsilon _j)}{4}\right)}\left(e^{\frac{2\pi i\delta _j}{m}}1\right)v_j_X^p๐\mu (x)๐\sigma (\delta )`$
$`=`$ $`{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _je^{\frac{2\pi ix_j}{m}}\left(e^{\frac{2\pi i\delta _j}{m}}1\right)v_j_X^p๐\mu (x)๐\sigma (\delta ).`$
Taking expectation with respect to $`\epsilon `$, and using the contraction principle, we see that
(22) $`{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}e^{\frac{2\pi ix_j}{m}}\left(e^{\frac{2\pi i\delta _j}{m}}1\right)v_j_X^p๐\mu (x)๐\sigma (\delta )`$
$`={\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _je^{\frac{2\pi ix_j}{m}}(e^{\frac{2\pi i\delta _j}{m}}1)v_j_X^pd\mu (x)d\sigma (\delta )`$
$`{\displaystyle _{\{1,0,1\}^n}}{\displaystyle _{_m^n}}2^p\left(\underset{1jn}{\mathrm{max}}\right|e^{\frac{2\pi i\delta _j}{m}}1\left|\right)^p๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jv_j_X^pd\mu (x)d\sigma (\delta )`$
$`\left({\displaystyle \frac{4\pi }{m}}\right)^p๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jv_j_X^p,`$
where in the last inequality above we used the fact that for $`\theta [0,\pi ]`$, $`|e^{i\theta }1|\theta `$.
Combining (7), (20), (21), and (22), we get that
$`2^p{\displaystyle \underset{j=1}{\overset{n}{}}}v_j_X^p\mathrm{\Gamma }^pm^p\left({\displaystyle \frac{4\pi }{m}}\right)^pn^{1\frac{p}{q}}๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jv_j_X^p=\left(4\pi \mathrm{\Gamma }\right)^pn^{1\frac{p}{q}}๐ผ_\epsilon {\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jv_j_X^p.`$
If $`p=q`$ we see that $`C_q(X)2\pi \mathrm{\Gamma }_q(X)`$. If $`p<q`$ then when $`v_1_X=\mathrm{}=v_n_X=1`$ we get that
$$\left(๐ผ_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _jv_j_X^q\right)^{1/q}\left(๐ผ_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _jv_j_X^p\right)^{1/p}=\mathrm{\Omega }\left(\frac{n^{1/q}}{\mathrm{\Gamma }}\right).$$
This means that $`X`$ has โequal norm cotype $`q`$โ, implying that $`X`$ has cotype $`q^{}`$ for every $`q^{}>q`$ (see , , for quantitative versions of this statement). When $`q=2`$ this implies that $`X`$ has cotype $`2`$ (see and the references therein).
### 5.2. Proof of Theorem 1.2 and Theorem 1.4
The proof of Theorem 1.2 and Theorem 1.4 is based on several lemmas. Fix an odd integer $`k`$, with $`k<\frac{m}{2}`$, and assume that $`1pq`$. Given $`j\{1,\mathrm{},n\}`$, define $`S(j,k)_m^n`$ by
$$S(j,k):=\{y[k,k]^n_m^n:y_j0mod2\mathrm{and}\mathrm{}j,y_{\mathrm{}}1mod2\}.$$
For $`f:_m^nX`$ we define
(23) $`_j^{(k)}f(x)=\left(f{\displaystyle \frac{\mathrm{๐}_{S(j,k)}}{\mu (S(j,k))}}\right)(x)={\displaystyle \frac{1}{\mu (S(j,k))}}{\displaystyle _{S(j,k)}}f(x+y)๐\mu (y).`$
###### Lemma 5.1.
For every $`p1`$, every $`j\{1,\mathrm{},n\}`$, and every $`f:_m^nX`$,
$`{\displaystyle _{_m^n}}_j^{(k)}f(x)f(x)_X^p๐\mu (x)`$ $``$ $`2^pk^p๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`+2^{p1}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x).`$
###### Proof.
By convexity, for every $`x_m^n`$,
$`_j^{(k)}f(x)f(x)_X^p`$ $`=`$ $`{\displaystyle \frac{1}{\mu (S(j,k))}}{\displaystyle _{S(j,k)}}[f(x+y)f(x)]๐\mu (y)_X^p`$
$``$ $`{\displaystyle \frac{1}{\mu (S(j,k))}}{\displaystyle _{S(j,k)}}f(x)f(x+y)_X^p๐\mu (y).`$
Let $`x\{0,\mathrm{},k\}^n`$ be such that for all $`j\{1,\mathrm{},n\}`$, $`x_j`$ is a positive odd integer. Observe that there exists a geodesic $`\gamma :\{0,1,\mathrm{},x_{\mathrm{}}\}_m^n`$ such that $`\gamma (0)=0`$, $`\gamma (x_{\mathrm{}})=x`$ and for every $`t\{1,\mathrm{},x_{\mathrm{}}\}`$,$`\gamma (t)\gamma (t1)\{1,1\}^n`$. Indeed, we define $`\gamma (t)`$ inductively as follows: $`\gamma (0)=0`$, $`\gamma (1)=(1,1,\mathrm{},1)`$, and if $`t2`$ is odd then
$$\gamma (t)=\gamma (t1)+\underset{s=1}{\overset{n}{}}e_s\mathrm{and}\gamma (t+1)=\gamma (t1)+2\underset{\begin{array}{c}s\{1,\mathrm{},n\}\\ \gamma (t1)_s<x_s\end{array}}{}e_s.$$
Since all the coordinates of $`x`$ are odd, $`\gamma (x_{\mathrm{}})=x`$. In what follows we fix an arbitrary geodesic $`\gamma _x:\{0,1,\mathrm{},x_{\mathrm{}}\}_m^n`$ as above. For $`x(_m\{0\})^n`$ we denote $`|x|=(|x_1|,\mathrm{},|x_n|)`$ and $`\mathrm{sign}(x)=(\mathrm{sign}(x_1),\mathrm{},\mathrm{sign}(x_n))`$. If $`x[k,k]^n`$ is such that all of its coordinates are odd, then we define $`\gamma _x=\mathrm{sign}(x)\gamma _{|x|}`$ (where the multiplication is coordinate-wise).
If $`yS(j,k)`$ then all the coordinates of $`y\pm e_j`$ are odd. We can thus define two geodesic paths
$$\gamma _{x,y}^{+1}=x+e_j+\gamma _{ye_j}\mathrm{and}\gamma _{x,y}^1=xe_j+\gamma _{y+e_j},$$
where the addition is point-wise.
For $`z_m^n`$ and $`\epsilon \{1,1\}^n`$ define
$$\begin{array}{c}F^{+1}(z,\epsilon )=\{(x,y)_m^n\times S(j,k):t\{1,\mathrm{},ye_j_{\mathrm{}}\},\hfill \\ \hfill \gamma _{x,y}^{+1}(t1)=z,\gamma _{x,y}^{+1}(t)=z+\epsilon \},\end{array}$$
and
$$\begin{array}{c}F^1(z,\epsilon )=\{(x,y)_m^n\times S(j,k):t\{1,\mathrm{},y+e_j_{\mathrm{}}\},\hfill \\ \hfill \gamma _{x,y}^1(t1)=z,\gamma _{x,y}^1(t)=z+\epsilon \}.\end{array}$$
###### Claim 5.2.
For every $`z,w_m^n`$ and $`\epsilon ,\delta \{1,1\}^n`$,
$$|F^{+1}(z,\epsilon )|+|F^1(z,\epsilon )|=|F^{+1}(w,\delta )|+|F^1(w,\delta )|.$$
###### Proof.
Define $`\psi :_m^n\times S(j,k)_m^n\times S(j,k)`$ by
$$\psi (x,y)=(w\epsilon \delta z+\epsilon \delta x,\epsilon \delta y).$$
We claim that $`\psi `$ is a bijection between $`F^{+1}(z,\epsilon )`$ and $`F^{\epsilon _j\delta _j}(w,\delta )`$, and also $`\psi `$ is a bijection between $`F^1(z,\epsilon )`$ and $`F^{\epsilon _j\delta _j}(w,\delta )`$. Indeed, if $`(x,y)F^{+1}(z,\epsilon )`$ then there exists $`t\{1,\mathrm{},ye_j_{\mathrm{}}\}`$ such that $`\gamma _{x,y}^{+1}(t1)=z`$ and $`\gamma _{x,y}^{+1}(t)=z+\epsilon `$. The path $`w\epsilon \delta z+\epsilon \delta \gamma _{x,y}^{+1}`$ equals the path $`\gamma _{\psi (x,y)}^{\epsilon _j\delta _j}`$, which by definition goes through $`w`$ at time $`t1`$ and $`w+\delta `$ at time $`t`$. Since these transformations are clearly invertible, we obtain the required result for $`F^{+1}(z,\epsilon )`$. The proof for $`F^1(z,\epsilon )`$ is analogous. โ
###### Claim 5.3.
Denote $`N=|F^{+1}(z,\epsilon )|+|F^1(z,\epsilon )|`$, which is independent of $`z_m^n`$ and $`\epsilon \{1,1\}^n`$, by Claim 5.2. Then
$$N\frac{k|S(j,k)|}{2^{n1}}.$$
###### Proof.
We have that
$`Nm^n2^n`$ $`=`$ $`{\displaystyle \underset{(z,\epsilon )_m^n\times \{1,1\}^n}{}}\left(|F^{+1}(z,\epsilon )|+|F^1(z,\epsilon )|\right)`$
$`=`$ $`{\displaystyle \underset{(z,\epsilon )_m^n\times \{1,1\}^n}{}}\left({\displaystyle \underset{(x,y)_m^n\times S(j,k)}{}}{\displaystyle \underset{t=1}{\overset{ye_j_{\mathrm{}}}{}}}\mathrm{๐}_{\{\gamma _{x,y}^{+1}(t1)=z\gamma _{x,y}^{+1}(t)=z+\epsilon \}}\right)`$
$`+{\displaystyle \underset{(z,\epsilon )_m^n\times \{1,1\}^n}{}}\left({\displaystyle \underset{(x,y)_m^n\times S(j,k)}{}}{\displaystyle \underset{t=1}{\overset{y+e_j_{\mathrm{}}}{}}}\mathrm{๐}_{\{\gamma _{x,y}^1(t1)=z\gamma _{x,y}^1(t)=z+\epsilon \}}\right)`$
$`=`$ $`{\displaystyle \underset{(x,y)_m^n\times S(j,k)}{}}ye_j_{\mathrm{}}+{\displaystyle \underset{(x,y)_m^n\times S(j,k)}{}}y+e_j_{\mathrm{}}`$
$``$ $`2km^n|S(j,k)|.`$
We now conclude the proof of Lemma 5.1. Observe that for $`x_m^n`$ and $`yS(j,k)`$,
$`{\displaystyle \frac{f(x)f(x+y)_X^p}{2^{p1}}}`$ $``$ $`f(x)f(x+e_j)_X^p`$
$`+ye_j_{\mathrm{}}^{p1}{\displaystyle \underset{t=1}{\overset{ye_j_{\mathrm{}}}{}}}f(\gamma _{x,y}^{+1}(t))f(\gamma _{x,y}^{+1}(t1))_X^p`$
$``$ $`f(x)f(x+e_j)_X^p`$
$`+k^{p1}{\displaystyle \underset{t=1}{\overset{ye_j_{\mathrm{}}}{}}}f(\gamma _{x,y}^{+1}(t))f(\gamma _{x,y}^{+1}(t1))_X^p,`$
$`{\displaystyle \frac{f(x)f(x+y)_X^p}{2^{p1}}}`$ $``$ $`f(x)f(xe_j)_X^p`$
$`+y+e_j_{\mathrm{}}^{p1}{\displaystyle \underset{t=1}{\overset{y+e_j_{\mathrm{}}}{}}}f(\gamma _{x,y}^1(t))f(\gamma _{x,y}^1(t1))_X^p`$
$``$ $`f(x)f(xe_j)_X^p`$
$`+k^{p1}{\displaystyle \underset{t=1}{\overset{y+e_j_{\mathrm{}}}{}}}f(\gamma _{x,y}^1(t))f(\gamma _{x,y}^1(t1))_X^p.`$
Averaging inequalities (LABEL:eq:geodesic+) and (LABEL:eq:geodesic-), and integrating, we get that
(27) $`{\displaystyle \frac{1}{\mu (S(j,k))}}{\displaystyle _{_m^n}}{\displaystyle _{S(j,k)}}f(x)f(x+y)_X^p๐\mu (y)๐\mu (x)`$
$`2^{p1}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x)`$
$`+(2k)^{p1}{\displaystyle \frac{N2^n}{|S(j,k)|}}๐ผ_\epsilon {\displaystyle _{_m^n}}f(z+\epsilon )f(z)_X^p๐\mu (z)`$
(28) $`2^{p1}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x)`$
$`+(2k)^p๐ผ_\epsilon {\displaystyle _{_m^n}}f(z+\epsilon )f(z)_X^p๐\mu (z),`$
where in (27) we used Claim 5.2 and in (28) we used Claim 5.3. By (5.2), this completes the proof of Lemma 5.1. โ
Lemma 5.4 below is the heart of our proof. It contains the cancellation of terms which is key to the validity of Theorem 1.2 and Theorem 1.4.
###### Lemma 5.4.
For every $`f:_m^nX`$, every integer $`n`$, every even integer $`m`$, every $`\epsilon \{1,1\}^n`$, every odd integer $`k<m/2`$, and every $`p1`$,
$`{\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]_X^p๐\mu (x)`$
$`3^{p1}{\displaystyle _{_m^n}}f(x+\epsilon )f(x\epsilon )_X^p๐\mu (x)`$
$`+{\displaystyle \frac{24^pn^{2p1}}{k^p}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x).`$
We postpone the proof of Lemma 5.4 to Section 5.3, and proceed to prove Theorem 1.2 and Theorem 1.4 assuming its validity.
###### Proof of Theorem 1.2 and Theorem 1.4.
Taking expectations with respect to $`\epsilon \{1,1\}^n`$ in Lemma 5.4 we get that
(29)
$`๐ผ_\epsilon {\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]_X^p๐\mu (x)`$
$`3^{p1}๐ผ_\epsilon {\displaystyle _{_m^n}}2^{p1}\left(f(x+\epsilon )f(x)_X^p+f(x)f(x\epsilon )_X^p\right)๐\mu (x)`$
$`+{\displaystyle \frac{24^pn^{2p1}}{k^p}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x)`$
$`{\displaystyle \frac{6^p}{3}}๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`+{\displaystyle \frac{24^pn^{2p1}}{k^p}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x).`$
Fix $`x_m^n`$ and let $`m`$ be an integer which is divisible by $`4`$ such that $`m6n^{2+1/q}`$. Fixing $`C>C_q^{(p)}(X)`$, and applying the definition of $`C_q^{(p)}(X)`$ to the vectors $`\left\{_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right\}_{j=1}^n`$, we get
(30)
$$\begin{array}{c}๐ผ_\epsilon \underset{j=1}{\overset{n}{}}\epsilon _j[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)]_X^p\hfill \\ \hfill \frac{1}{C^pn^{1p/q}}\underset{j=1}{\overset{n}{}}_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)_X^p.\end{array}$$
Now, for every $`j\{1,\mathrm{},n\}`$,
(31)
$$\begin{array}{c}\underset{s=1}{\overset{m/4}{}}_j^{(k)}f\left(x+2se_j\right)_j^{(k)}f\left(x+2(s1)e_j\right)_X^p\hfill \\ \hfill \left(\frac{4}{m}\right)^{p1}_j^{(k)}f\left(x+\frac{m}{2}e_j\right)_j^{(k)}f\left(x\right)_X^p.\end{array}$$
Averaging (31) over $`x_m^n`$ we get that
(32)
$$\begin{array}{c}_{_m^n}_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)_X^p๐\mu (x)\hfill \\ \hfill \left(\frac{4}{m}\right)^p_{_m^n}_j^{(k)}f\left(x+\frac{m}{2}e_j\right)_j^{(k)}f\left(x\right)_X^p๐\mu (x).\end{array}$$
Combining (30) and (32) we get the inequality
(33)
$$\begin{array}{c}๐ผ_\epsilon _{_m^n}\underset{j=1}{\overset{n}{}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]_X^p๐\mu (x)\hfill \\ \hfill \frac{1}{C^pn^{1p/q}}\left(\frac{4}{m}\right)^p\underset{j=1}{\overset{n}{}}_{_m^n}_j^{(k)}f\left(x+\frac{m}{2}e_j\right)_j^{(k)}f\left(x\right)_X^p๐\mu (x).\end{array}$$
Now, for every $`j\{1,\mathrm{},n\}`$,
$`{\displaystyle _{_m^n}}_j^{(k)}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)_j^{(k)}f\left(x\right)_X^p๐\mu (x)`$
$`{\displaystyle \frac{1}{3^{p1}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`{\displaystyle _{_m^n}}_j^{(k)}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x+{\displaystyle \frac{m}{2}}e_j\right)_X^p๐\mu (x)`$
$`{\displaystyle _{_m^n}}_j^{(k)}f\left(x\right)f\left(x\right)_X^p๐\mu (x)`$
$`={\displaystyle \frac{1}{3^{p1}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`2{\displaystyle _{_m^n}}_j^{(k)}f\left(x\right)f\left(x\right)_X^p๐\mu (x)`$
$`{\displaystyle \frac{1}{3^{p1}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`2^{p+1}k^p๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`2^p{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x),`$
where we used Lemma 5.1.
Combining (5.2) with (33), we see that
(35)
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f\left(x+{\displaystyle \frac{m}{2}}e_j\right)f\left(x\right)_X^p๐\mu (x)`$
$`{\displaystyle \frac{(3Cm)^pn^{1\frac{p}{q}}}{34^p}}๐ผ_\epsilon {\displaystyle _{_m^n}}{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]_X^p๐\mu (x)`$
$`+6^pk^pn๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`+6^p{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x)`$
$`\left({\displaystyle \frac{(18Cm)^pn^{1\frac{p}{q}}}{4^p}}+6^pk^pn\right)๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^p๐\mu (x)`$
$`+\left({\displaystyle \frac{(3Cm)^pn^{1\frac{p}{q}}}{4^p}}{\displaystyle \frac{24^pn^{2p1}}{k^p}}+6^p\right){\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^p๐\mu (x)`$
(36)
$`(18Cm)^pn^{1\frac{p}{q}}(๐ผ_\epsilon {\displaystyle _{_m^n}}f(x+\epsilon )f(x)_X^pd\mu (x)d\sigma (\epsilon )`$
$`+{\displaystyle \frac{1}{n}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(x+e_j)f(x)_X^pd\mu (x)),`$
where in (35) we used (29), and (36) holds true when we choose $`4n^2k\frac{3m}{4n^{1/q}}`$ (which is possible if we assume that $`m6n^{2+1/q}`$). By Lemma 2.7, this completes the proof of Theorem 1.4. โ
### 5.3. Proof of Lemma 5.4
Fix $`\epsilon \{1,1\}^n`$, and $`x_m^n`$. Consider the following two sums:
$`A_f(x,\epsilon )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _j\left[_j^{(k)}f(x+e_j)_j^{(k)}f(xe_j)\right]`$
$`=`$ $`{\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{y_m^n}{}}a_y(x,\epsilon )f(y),`$
and
$`B_f(x,\epsilon )`$ $`=`$ $`{\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{zx(k,k)^n(2)^n}{}}[f(z+\epsilon )f(z\epsilon )]`$
$`=`$ $`{\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{y_m^n}{}}b_y(x,\epsilon )f(y),`$
where $`a_y(x,\epsilon ),b_y(x,\epsilon )`$ are appropriately chosen coefficients, which are independent of $`f`$.
For $`x_m^n`$ define $`S(x)_m^n`$,
$$\begin{array}{c}S(x)=\{yx+(2+1)^n:d_{_m^n}(y,x)=k,\hfill \\ \hfill \text{and}|\{j:|y_jx_j|kmodm\}|2\}.\end{array}$$
###### Claim 5.5.
For $`x_m^n`$ and $`yS(x)`$, $`a_y(x,\epsilon )=b_y(x,\epsilon )`$.
###### Proof.
If there exists a coordinate $`j\{1,\mathrm{},n\}`$ such that $`x_jy_j`$ is even, then it follows from our definitions that $`a_y(x,\epsilon )=b_y(x,\epsilon )=0`$. Similarly, if $`d_{_m^n}(x,y)>k`$ then $`a_y(x,\epsilon )=b_y(x,\epsilon )=0`$ (because $`k`$ is odd). Assume that $`xy(2+1)^n`$. If $`d_{_m^n}(y,x)<k`$ then for each $`j`$ the term $`f(y)`$ cancels in $`_j^{(k)}f(x+e_j)_j^{(k)}(xe_j)`$, implying that $`a_y(x,\epsilon )=0`$. Similarly, in the sum defining $`B_f(x,\epsilon )`$ the term $`f(y)`$ appears twice, with opposite signs, so that $`b_y(x,\epsilon )=0`$.
It remains to deal with the case $`|\{j:|y_jx_j|kmodm\}|=1`$. We may assume without loss of generality that
$$|y_1x_1|kmodm\text{ and for }j2,y_jx_j(k,k)modm.$$
If $`y_1x_1kmodm`$ then $`a_y(x,\epsilon )=\epsilon _1`$, since in the terms corresponding to $`j2`$ in the definition of $`A_f(x,\epsilon )`$ the summand $`f(y)`$ cancels out. We also claim that in this case $`b_y(x,\epsilon )=\epsilon _1`$. Indeed, if $`\epsilon _1=1`$ then $`f(y)`$ appears in the sum defining $`B_f(x,\epsilon )`$ only in the term corresponding to $`z=y\epsilon `$, while if $`\epsilon _1=1`$ then $`f(y)`$ appears in this sum only in the term corresponding to $`z=y+\epsilon `$, in which case its coefficient is $`1`$. In the case $`y_1x_1kmodm`$ the same reasoning shows that $`a_y(x,\epsilon )=b_y(x,\epsilon )=\epsilon _1`$. โ
By Claim 5.5 we have
(39) $`A_f(x,\epsilon )B_f(x,\epsilon )={\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{yS(x)}{}}[a_y(x,\epsilon )b_y(x,\epsilon )]f(y).`$
Thus,
$`{\displaystyle _{_m^n}}A_f(x,\epsilon )_X^p๐\mu (x)`$ $``$ $`3^{p1}{\displaystyle _{_m^n}}B_f(x,\epsilon )_X^p๐\mu (x)`$
$`+3^{p1}{\displaystyle _{_m^n}}{\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{yS(x)}{}}a_y(x,\epsilon )f(y)_X^p๐\mu (x)`$
$`+3^{p1}{\displaystyle _{_m^n}}{\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{yS(x)}{}}b_y(x,\epsilon )f(y)_X^p๐\mu (x).`$
Thus Lemma 5.4 will be proved once we establish the following inequalities
(40)
$$_{_m^n}B_f(x,\epsilon )_X^p๐\mu (x)_{_m^n}f(x+\epsilon )f(x\epsilon )_X^p๐\mu (x),$$
(41)
$$\begin{array}{c}_{_m^n}\frac{1}{k(k+1)^{n1}}\underset{yS(x)}{}a_y(x,\epsilon )f(y)_X^p๐\mu (x)\hfill \\ \hfill \frac{8^pn^{2p1}}{k^p}\underset{j=1}{\overset{n}{}}_{_m^n}f(x+e_j)f(x)_X^p,\end{array}$$
and
(42)
$$\begin{array}{c}_{_m^n}\frac{1}{k(k+1)^{n1}}\underset{yS(x)}{}b_y(x,\epsilon )f(y)_X^p๐\mu (x)\hfill \\ \hfill \frac{8^pn^{2p1}}{k^p}\underset{j=1}{\overset{n}{}}_{_m^n}f(x+e_j)f(x)_X^p.\end{array}$$
Inequality (40) follows directly from the definition of $`B_f(x,\epsilon )`$, by convexity. Thus, we pass to the proof of (41) and (42).
For $`j=1,2,\mathrm{},n`$ define for $`yS(x)`$,
$$\tau _j^x(y)=\{\begin{array}{cc}y2ke_j\hfill & y_jx_jkmodm,\hfill \\ y\hfill & \text{otherwise},\hfill \end{array}$$
and set $`\tau _j^x(y)=y`$ when $`yS(x)`$. Observe that the following identity holds true:
(43) $`\tau _j^x(y)=\tau _j^0(yx)+x.`$
###### Claim 5.6.
Assume that for every $`j\{1,2,\mathrm{},n\}`$, $`x,y_m^n`$ and $`\epsilon \{1,1\}^n`$, we are given a real number $`\eta _j(x,y,\epsilon )[1,1]`$. Then
$$\begin{array}{c}_{_m^n}\frac{1}{k(k+1)^{n1}}\underset{j=1}{\overset{n}{}}\underset{y_m^n}{}\eta _j(x,y,\epsilon )\left[f(y)f(\tau _j^x(y))\right]_X^p๐\mu (x)\hfill \\ \hfill \frac{8^pn^{2p1}}{2k^p}\underset{j=1}{\overset{n}{}}_{_m^n}f(x+e_j)f(x)_X^p๐\mu (x).\end{array}$$
###### Proof.
Denote by $`N(x,\epsilon )`$ the number of nonzero summands in
$$\underset{j=1}{\overset{n}{}}\underset{y_m^n}{}\eta _j(x,y,\epsilon )\left[f(y)f(\tau _j^x(y))\right].$$
For every $`\mathrm{}2`$ let $`S^{\mathrm{}}(x)`$ be the set of all $`yS(x)`$ for which the number of coordinates $`j`$ such that $`y_jx_j\{k,k\}modm`$ equals $`\mathrm{}`$. Then $`|S^{\mathrm{}}(x)|=\left(\genfrac{}{}{0pt}{}{n}{\mathrm{}}\right)2^{\mathrm{}}(k1)^n\mathrm{}`$. Moreover, for $`yS^{\mathrm{}}(x)`$ we have that $`y\tau _j^x(y)`$ for at most $`\mathrm{}`$ values of $`j`$. Hence
$`N(x,\epsilon ){\displaystyle \underset{\mathrm{}=2}{\overset{n}{}}}|S^{\mathrm{}}(x)|\mathrm{}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=2}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\mathrm{}}}\right)2^{\mathrm{}}(k1)^n\mathrm{}\mathrm{}`$
$`=`$ $`2n\left[(k+1)^{n1}(k1)^{n1}\right]{\displaystyle \frac{4n^2}{k^2}}k(k+1)^{n1}.`$
Now, using (43), we get
$`{\displaystyle _{_m^n}}{\displaystyle \frac{1}{k(k+1)^{n1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{y_m^n}{}}\eta _j(x,y,\epsilon )\left[f(y)f(\tau _j^x(y))\right]_X^p๐\mu (x)`$
$`=`$ $`{\displaystyle _{_m^n}}\left({\displaystyle \frac{N(x,\epsilon )}{k(k+1)^{n1}}}\right)^p{\displaystyle \frac{1}{N(x,\epsilon )}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{y_m^n}{}}\eta _j(x,y,\epsilon )\left[f(y)f(\tau _j^x(y))\right]_X^p๐\mu (x)`$
$``$ $`{\displaystyle _{_m^n}}{\displaystyle \frac{N(x,\epsilon )^{p1}}{k^p(k+1)^{(n1)p}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{y_m^n}{}}f(y)f(\tau _j^x(y))_X^pd\mu (x)`$
$``$ $`{\displaystyle \frac{4^{p1}n^{2p2}}{k^{2p1}(k+1)^{n1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{y_m^n}{}}{\displaystyle _{_m^n}}f(y)f(\tau _j^x(y))_X^p๐\mu (x)`$
$`=`$ $`{\displaystyle \frac{4^{p1}n^{2p2}}{k^{2p1}(k+1)^{n1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{z_m^n}{}}{\displaystyle _{_m^n}}f(z+x)f(\tau _j^0(z)+x)_X^p๐\mu (x).`$
Consider the following set:
$$E_j=\{z_m^n:\tau _j^0(z)=z2ke_j\}.$$
Observe that that for every $`j`$,
$`|E_j|`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=1}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{\mathrm{}}}\right)2^{\mathrm{}}(k1)^{n1\mathrm{}}`$
$``$ $`(k+1)^{n1}(k1)^{n1}{\displaystyle \frac{2n}{k}}(k+1)^{n1}.`$
Using the translation invariance of the Haar measure on $`_m^n`$ we get that
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{z_m^n}{}}{\displaystyle _{_m^n}}f(z+x)f(\tau _j^0(z)+x)_X^p๐\mu (x)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{zE_j}{}}{\displaystyle _{_m^n}}f(z+x)f(z+x2ke_j)_X^p๐\mu (x)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}|E_j|{\displaystyle _{_m^n}}f(w)f(w2ke_j)_X^p๐\mu (w)`$
$``$ $`{\displaystyle \frac{2n}{k}}(k+1)^{n1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(w)f(w2ke_j)_X^p๐\mu (w)`$
$``$ $`{\displaystyle \frac{2n}{k}}(k+1)^{n1}`$
$`\times {\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}((2k)^{p1}{\displaystyle \underset{t=1}{\overset{2k}{}}}f(w(t1)e_j)f(wte_j)_X^p)d\mu (w)`$
$``$ $`2^{p+1}nk^{p1}(k+1)^{n1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}f(z+e_j)f(z)_X^p๐\mu (z),`$
where in (5.3) we used (5.3). Combining (LABEL:eq:step1) and (5.3) completes the proof of Claim 5.6. โ
By Claim 5.6, inequalities (41) and (42), and hence also Lemma 5.4, will be proved once we establish the following identities:
(48) $`{\displaystyle \underset{yS(x)}{}}a_y(x,\epsilon )f(y)={\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{y_m^n}{}}\epsilon _j\left[f(y)f(\tau _j^x(y))\right].`$
and
(49) $`{\displaystyle \underset{yS(x)}{}}b_y(x,\epsilon )f(y)={\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{y_m^n}{}}\delta _j(x,y,\epsilon )\left[f(y)f(\tau _j^x(y))\right],`$
for some $`\delta _j(x,y,\epsilon )\{1,0,1\}`$.
Identity (48) follows directly from the fact that (5.3) implies that for every $`yS(x)`$,
$$a_y(x,\epsilon )=\underset{j:y_jx_jkmodm}{}\epsilon _j\underset{j:y_jx_jkmodm}{}\epsilon _j.$$
It is enough to prove identity (49) for $`x=0`$, since $`b_y(x,\epsilon )=b_{yx}(0,\epsilon )`$. To this end we note that it follows directly from (5.3) that for every $`yS(0)`$
$$b_y(0,\epsilon )=\{\begin{array}{cc}1\hfill & jy_j\epsilon _jkmodm\text{ and }\mathrm{}y_{\mathrm{}}\epsilon _jkmodm\hfill \\ 1\hfill & jy_j\epsilon _jkmodm\text{ and }\mathrm{}y_{\mathrm{}}\epsilon _jkmodm\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
For $`yS(0)`$ define
$$y_j^{}=\{\begin{array}{cc}y_j\hfill & y_j\{k,k\}modm\hfill \\ y_j\hfill & \text{otherwise}.\hfill \end{array}$$
Since $`b_y(0,\epsilon )=b_y^{}(0,\epsilon )`$ we get that
(50) $`{\displaystyle \underset{yS(0)}{}}b_y(0,\epsilon )f(y)={\displaystyle \frac{1}{2}}{\displaystyle \underset{yS(0)}{}}b_y(0,\epsilon )\left[f(y)f(y^{})\right].`$
Define for $`\mathrm{}\{1,\mathrm{},n+1\}`$ a vector $`y^{_{\mathrm{}}}_m^n`$ by
$$y_j^{_{\mathrm{}}}=\{\begin{array}{cc}y_j\hfill & j<\mathrm{}\text{ and }y_j\{k,k\}modm\hfill \\ y_j\hfill & \text{otherwise}.\hfill \end{array}$$
Then $`y^{_{n+1}}=y^{}`$, $`y^_1=y`$ and by (50)
$$\underset{yS(0)}{}b_y(0,\epsilon )f(y)=\frac{1}{2}\underset{\mathrm{}=1}{\overset{n}{}}\underset{yS(0)}{}b_y(0,\epsilon )\left[f(y^{_{\mathrm{}}})f(y^{_{\mathrm{}+1}})\right].$$
Since whenever $`y^{_{\mathrm{}}}y^{_{\mathrm{}+1}}`$, each of these vectors is obtained from the other by flipping the sign of the $`\mathrm{}`$-th coordinate, which is in $`\{k,k\}modm`$, this implies the representation (49). The proof of Lemma 5.4 is complete. โ
## 6. A nonlinear version of the Maurey-Pisier theorem
In what follows we denote by $`\mathrm{๐๐ข๐๐ }(_m^n)`$ the graph on $`_m^n`$ in which $`x,y_m^n`$ are adjacent if for every $`i\{1,\mathrm{},n\}`$$`x_iy_i\{\pm 1modm\}`$.
For technical reasons that will become clear presently, given $`\mathrm{},n`$ we denote by $`(;n,\mathrm{})`$ the infimum over $`>0`$ such that for every even $`m`$ and for every $`f:_m^n`$,
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(f\left(x+\mathrm{}e_j\right),f(x))^2๐\mu (x)^2\mathrm{}^2n๐ผ_\epsilon {\displaystyle _{_m^n}}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x).`$
###### Lemma 6.1.
For every metric space $`(,d_{})`$, every $`n,a`$, every even $`m,r`$ with $`0r<m`$, and every $`f:_m^n`$,
(51)
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}d_{}(f\left(x+(am+r)e_j\right),f(x))^2๐\mu (x)\hfill \\ \hfill \mathrm{min}\{r^2,(mr)^2\}n๐ผ_\epsilon _{_m^n}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x).\end{array}$$
In particular, $`(;n,\mathrm{})1`$ for every $`n`$ and every even $`\mathrm{}`$.
###### Proof.
The left-hand side of (51) depends only on $`r`$, and remains unchanged if we replace $`r`$ by $`mr`$. We may thus assume that $`a=0`$ and $`rmr`$. Fix $`x_m^n`$ and $`j\{1,\mathrm{}n\}`$. Observe that
$$\left\{x+\frac{1(1)^k}{2}\underset{rj}{}e_r+ke_j\right\}_{k=0}^r$$
is a path of length $`r`$ joining $`x`$ and $`x+re_j`$ in the graph $`\mathrm{๐๐ข๐๐ }(_m^n)`$. Thus the distance between $`x`$ and $`x+re_j`$ in the graph $`\mathrm{๐๐ข๐๐ }(_m^n)`$ equals $`r`$. If $`(x=w_0,w_1,\mathrm{},w_r=x+re_j)`$ is a geodesic joining $`x`$ and $`x+re_j`$ in $`\mathrm{๐๐ข๐๐ }(_m^n)`$, then by the triangle inequality
(52) $`d_{}(f(x+re_j),f(x))^2r{\displaystyle \underset{k=1}{\overset{r}{}}}d_{}(f(w_k),f(w_{k1}))^2.`$
Observe that if we sum (52) over all geodesics joining $`x`$ and $`x+re_j`$ in $`\mathrm{๐๐ข๐๐ }(_m^n)`$, and then over all $`x_m^n`$, then in the resulting sum each edge in $`\mathrm{๐๐ข๐๐ }(_m^n)`$ appears the same number of times. Thus, averaging this inequality over $`x_m^n`$ we get
$`{\displaystyle _{_m^n}}d_{}(f(x+re_j),f(x))^2d\mu (x)r^2๐ผ_\epsilon [d_{}(f(x+\epsilon ),f(x))]^2.`$
Summing over $`j=1,\mathrm{}n`$ we obtain the required result. โ
###### Lemma 6.2.
For every four integers $`\mathrm{},k,s,t`$,
$$(;\mathrm{}k,st)(;\mathrm{},s)(;k,t).$$
###### Proof.
Let $`m`$ be an even integer and take a function $`f:_m^\mathrm{}k`$. Fix $`x_m^\mathrm{}k`$ and $`\epsilon \{1,1\}^\mathrm{}k`$. Define $`g:_m^{\mathrm{}}`$ by
$$g(y)=f\left(x+\underset{r=1}{\overset{k}{}}\underset{j=1}{\overset{\mathrm{}}{}}\epsilon _{j+(r1)\mathrm{}}y_je_{j+(r1)\mathrm{}}\right).$$
By the definition of $`(;\mathrm{},s)`$, applied to $`g`$, for every $`_1>(;\mathrm{},s)`$,
$`{\displaystyle \underset{a=1}{\overset{\mathrm{}}{}}}{\displaystyle __m^{\mathrm{}}}d_{}(f(x+{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\epsilon _{j+(r1)\mathrm{}}y_je_{j+(r1)\mathrm{}}+s{\displaystyle \underset{r=1}{\overset{k}{}}}\epsilon _{a+(r1)\mathrm{}}e_{a+(r1)\mathrm{}}),`$
$`f(x+{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\epsilon _{j+(r1)\mathrm{}}y_je_{j+(r1)\mathrm{}}))^2d\mu __m^{\mathrm{}}(y)`$
$``$ $`_1^2s^2\mathrm{}๐ผ_\delta {\displaystyle __m^{\mathrm{}}}d_{}(f(x+{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\epsilon _{j+(r1)\mathrm{}}(y_j+\delta _j)e_{j+(r1)\mathrm{}}),`$
$`f(x+{\displaystyle \underset{r=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\epsilon _{j+(r1)\mathrm{}}y_je_{j+(r1)\mathrm{}}))^2d\mu __m^{\mathrm{}}(y).`$
Averaging this inequality over $`x_m^\mathrm{}k`$ and $`\epsilon \{1,1\}^\mathrm{}k`$, and using the translation invariance of the Haar measure, we get that
(53)
$$\begin{array}{c}๐ผ_\epsilon \underset{a=1}{\overset{\mathrm{}}{}}_{_m^\mathrm{}k}d_{}(f\left(x+s\underset{r=1}{\overset{k}{}}\epsilon _{a+(r1)\mathrm{}}e_{a+(r1)\mathrm{}}\right),f(x))^2๐\mu _{_m^\mathrm{}k}(x)\hfill \\ \hfill _1^2s^2\mathrm{}๐ผ_\epsilon _{_m^\mathrm{}k}d_{}(f\left(x+\epsilon \right),f\left(x\right))^2๐\mu _{_m^\mathrm{}k}(x).\end{array}$$
Next we fix $`x_m^\mathrm{}k`$, $`u\{1,\mathrm{},\mathrm{}\}`$, and define $`h_u:_m^k`$ by
$$h_u(y)=f\left(x+s\underset{r=1}{\overset{k}{}}y_re_{u+(r1)\mathrm{}}\right).$$
By the definition of $`(;k,t)`$, applied to $`h_u`$, for every $`_2>(;k,t)`$ we have
$`{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle _{_m^k}}d_{}(f(x+s{\displaystyle \underset{r=1}{\overset{k}{}}}y_re_{u+(r1)\mathrm{}}+ste_{u+(j1)\mathrm{}}),`$
$`f(x+s{\displaystyle \underset{r=1}{\overset{k}{}}}y_re_{u+(r1)\mathrm{}}))^2d\mu _{_m^k}(y)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle _{_m^k}}d_{}(h_u\left(y+te_j\right),h_u(y))^2๐\mu _{_m^k}(y)`$
$``$ $`_2^2t^2k๐ผ_\epsilon {\displaystyle _{_m^k}}d(h_u\left(y+\epsilon \right),h_u(y))^2๐\mu _{_m^k}(y)`$
$`=`$ $`_2^2t^2k๐ผ_\epsilon {\displaystyle _{_m^k}}d_{}(f(x+s{\displaystyle \underset{r=1}{\overset{k}{}}}(y_r+\epsilon _{u+(r1)\mathrm{}})e_{u+(r1)\mathrm{}}),`$
$`f(x+s{\displaystyle \underset{r=1}{\overset{k}{}}}y_re_{u+(r1)\mathrm{}}))^2d\mu _{_m^k}(y).`$
Summing this inequality over $`u\{1,\mathrm{},\mathrm{}\}`$ and averaging over $`x_m^\mathrm{}k`$, we get, using (53), that
$`{\displaystyle \underset{a=1}{\overset{\mathrm{}k}{}}}{\displaystyle _{_m^\mathrm{}k}}d_{}(f\left(x+ste_a\right),f(x))^2๐\mu (x)`$
$`_2^2t^2k๐ผ_\epsilon {\displaystyle \underset{u=1}{\overset{\mathrm{}}{}}}{\displaystyle _{_m^\mathrm{}k}}d_{}(f\left(x+s{\displaystyle \underset{r=1}{\overset{k}{}}}\epsilon _{u+(r1)\mathrm{}}e_{u+(r1)\mathrm{}}\right),f\left(x\right))^2๐\mu (x)`$
$`_2^2t^2k_1^2s^2\mathrm{}๐ผ_\epsilon {\displaystyle _{_m^\mathrm{}k}}d_{}(f\left(x+\epsilon \right),f\left(x\right))^2๐\mu (x).`$
This implies the required result. โ
###### Lemma 6.3.
Assume that there exist integers $`n_0,\mathrm{}_0>1`$ such that$`(;n_0,\mathrm{}_0)<1`$. Then there exists $`0<q<\mathrm{}`$ such that for every integer $`n`$,
$$m_q^{(2)}(;n,3n_0)2\mathrm{}_0n^{\mathrm{log}_{n_0}\mathrm{}_0}.$$
In particular, $`\mathrm{\Gamma }_q^{(2)}()<\mathrm{}`$.
###### Proof.
Let $`q<\mathrm{}`$ satisfy $`(,n_0,\mathrm{}_0)<n_0^{1/q}`$. Iterating Lemma 6.2 we get that for every integer $`k`$, $`(n_0^k,\mathrm{}_0^k)n_0^{k/q}`$. Denoting $`n=n_0^k`$ and $`m=2\mathrm{}_0^k`$, this implies that for every $`f:_m^n`$,
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}d_{}(f\left(x+\frac{m}{2}e_j\right),f(x))^2๐\mu (x)\hfill \\ \hfill \frac{1}{4}m^2n^{1\frac{2}{q}}๐ผ_\epsilon _{_m^n}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x).\end{array}$$
For $`f:_m^n^{}`$, where $`n^{}n`$, we define $`g:_m^n^{}\times _m^{nn^{}}`$ by $`g(x,y)=f(x)`$. Applying the above inqeuality to $`g`$ we obtain,
$$\begin{array}{c}\underset{j=1}{\overset{n^{}}{}}_{_m^n^{}}d_{}(f\left(x+\frac{m}{2}e_j\right),f(x))^2s\mu (x)\hfill \\ \hfill \frac{1}{4}m^2n^{1\frac{2}{q}}๐ผ_\epsilon _{_m^n^{}}๐(f(x+\epsilon ),f(x))^2๐\mu (x).\end{array}$$
Hence, by Lemma 2.7 we deduce that $`\mathrm{\Gamma }_q^{(2)}(;n_0^k,2\mathrm{}_0^k)3`$. For general $`n`$, let $`k`$ be the minimal integer such that $`nn_0^k`$. By Lemma 2.5 we get that $`\mathrm{\Gamma }(;n,2\mathrm{}_0^k)3n_0^{12/q}3n_0`$. In other words,
$`m_q^{(2)}(;n,3n_0)2\mathrm{}_0^k2\mathrm{}_0n^{\mathrm{log}_{n_0}\mathrm{}_0}.`$
###### Theorem 6.4.
Let $`n>1`$ be an integer, $`m`$ an even integer, and $`s`$ an integer divisible by $`4`$. Assume that $`\eta (0,1)`$ satisfies $`8^{sn}\sqrt{\eta }<\frac{1}{2}`$, and that there exists a mapping $`f:_m^n`$ such that
(54)
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}d_{}(f\left(x+se_j\right),f(x))^2๐\mu (x)\hfill \\ \hfill >(1\eta )s^2n๐ผ_\epsilon _{_m^n}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x).\end{array}$$
Then
$$c_{}\left(\left[s/4\right]_{\mathrm{}}^n\right)1+8^{sn}\sqrt{\eta }.$$
In particular, if $`(;n,s)=1`$ then $`c_{}\left([s/4]_{\mathrm{}}^n\right)=1`$.
###### Proof.
Observe first of all that (54) and Lemma 6.1 imply that $`m2s\sqrt{1\eta }>2s1`$, so that $`m2s`$. In what follows we will use the following numerical fact: If $`a_1,\mathrm{},a_r0`$ and $`0b\frac{1}{r}_{j=1}^ra_j`$, then
(55) $`{\displaystyle \underset{j=1}{\overset{r}{}}}\left(a_jb\right)^2{\displaystyle \underset{j=1}{\overset{r}{}}}a_j^2rb^2.`$
For $`x_m^n`$ let $`๐ข_j^+(x)`$ (resp. $`๐ข_j^{}(x)`$) be the set of all geodesics joining $`x`$ and $`x+se_j`$ (resp. $`xse_j`$) in the graph $`\mathrm{๐๐ข๐๐ }(_m^n)`$. As we have seen in the proof of Lemma 6.1, since $`s`$ is even, these sets are nonempty. Notice that if $`m=2s`$ then $`๐ข_j^+(x)=๐ข_j^{}(x)`$; otherwise $`๐ข_j^+(x)๐ข_j^{}(x)=\mathrm{}`$. Denote $`๐ข_j^\pm (x)=๐ข_j^+(x)๐ข_j^{}(x)`$, and for $`\pi ๐ข_j^\pm (x)`$,
$$\mathrm{sgn}(\pi )=\{\begin{array}{cc}+1\hfill & \text{ if }\pi ๐ข_j^+(x)\hfill \\ 1\hfill & \text{ otherwise}.\hfill \end{array}$$
Each geodesic in $`๐ข_j^\pm (x)`$ has length $`s`$. We write each $`\pi ๐ข_j^\pm (x)`$ as a sequence of vertices $`\pi =(\pi _0=x,\pi _1,\mathrm{},\pi _s=x+\mathrm{sgn}(\pi )se_j)`$. Using (55) with $`a_j=d_{}(f(\pi _j),f(\pi _{j1}))`$ and $`b=\frac{1}{s}d_{}(f\left(x+se_j\right),f(x))`$, which satisfy the conditions of (55) due to the triangle inequality, we get that for each $`\pi ๐ข_j^\pm (x)`$,
(56)
$$\begin{array}{c}\underset{\mathrm{}=1}{\overset{s}{}}\left[d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))\frac{1}{s}d_{}(f\left(x+\mathrm{sgn}(\pi )se_j\right),f(x))\right]^2\hfill \\ \hfill \underset{k=1}{\overset{s}{}}d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))^2\frac{1}{s}d_{}(f\left(x+\mathrm{sgn}(\pi )se_j\right),f(x))^2.\end{array}$$
By symmetry $`|๐ข_j^+(x)|=|๐ข_j^{}(x)|`$, and this value is independent of $`x_m^n`$ and $`j\{1,\mathrm{},n\}`$. Denote $`g=|๐ข_j^\pm (x)|`$, and observe that $`g22^{ns}`$. Averaging (56) over all $`x_m^n`$ and $`\pi ๐ข_j^\pm (x)`$, and summing over $`j\{1,\mathrm{},n\}`$, we get that
(57) $`{\displaystyle \frac{1}{g}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}{\displaystyle \underset{\pi ๐ข_j^\pm (x)}{}}{\displaystyle \underset{\mathrm{}=1}{\overset{s}{}}}[d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))`$
$`{\displaystyle \frac{1}{s}}d_{}(f(x+\mathrm{sgn}(\pi )se_j),f(x))]^2d\mu (x)`$
$`sn๐ผ_\epsilon {\displaystyle _{_m^n}}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x)`$
$`{\displaystyle \frac{1}{s}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(f\left(x+se_j\right),f(x))^2๐\mu (x)`$
$`<\eta sn๐ผ_\epsilon {\displaystyle _{_m^n}}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x).`$
Define $`\psi :_m^n`$ by
$$\begin{array}{c}\psi (x)=2\eta sn2^{sn}๐ผ_\epsilon [d_{}(f(x+\epsilon ),f(x))^2]\hfill \\ \hfill \underset{j=1}{\overset{n}{}}\underset{\pi ๐ข_j^\pm (x)}{}\underset{\mathrm{}=1}{\overset{s}{}}\left[d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))\frac{1}{s}d_{}(f\left(x+\mathrm{sgn}(\pi )se_j\right),f(x))\right]^2.\end{array}$$
Inequality (57), together with the bound on $`g`$, implies that
$$0<_{_m^n}\psi (x)๐\mu (x)=\frac{1}{(2s1)^n}_{_m^n}\underset{\begin{array}{c}y_m^n\\ d_{_m^n}(x,y)<s\end{array}}{}\psi (y)d\mu (x).$$
It follows that there exists $`x^0_m^n`$ such that
(58) $`{\displaystyle \underset{\begin{array}{c}y_m^n\\ d_{_m^n}(x^0,y)<s\end{array}}{}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{\pi ๐ข_j^+(x){\scriptscriptstyle ๐ข_j^{}(x)}}{}}`$
$`{\displaystyle \underset{\mathrm{}=1}{\overset{s}{}}}\left[d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1)){\displaystyle \frac{1}{s}}d_{}(f\left(y+\mathrm{sgn}(\pi )se_j\right),f(y))\right]^2`$
$`<2\eta sn2^{sn}{\displaystyle \underset{\begin{array}{c}y_m^n\\ d_{_m^n}(x^0,y)<s\end{array}}{}}๐ผ_\epsilon \left[d_{}(f(y+\epsilon ),f(y))^2\right].`$
By scaling the metric $`d_{}`$ we may assume without loss of generality that
(59) $`{\displaystyle \frac{1}{(2s1)^n}}{\displaystyle \underset{\begin{array}{c}y_m^n\\ d_{_m^n}(x^0,y)<s\end{array}}{}}๐ผ_\epsilon \left[d_{}(f(y+\epsilon ),f(y))^2\right]=1.`$
It follows that there exists $`y^0_m^n`$ satisfying $`d_{_m^n}(x^0,y^0)<s`$ such that
(60) $`๐ผ_\epsilon \left[d_{}(f(y^0+\epsilon ),f(y^0))^2\right]1.`$
By translating the argument of $`f`$, and multiplying (coordinate-wise) by an appropriate sign vector in $`\{1,1\}^n`$, we may assume that $`y^0=0`$ and all the coordinates of $`x^0`$ are nonnegative. Observe that this implies that every $`y\{0,1,\mathrm{},s1\}^n`$ satisfies $`d_{_m^n}(x^0,y)<s`$. Thus (58), and (59) imply that for every $`y\{0,1,\mathrm{},s1\}^n`$, every $`j\{1,\mathrm{},n\}`$, every $`\pi ๐ข_j^\pm (y)`$, and every $`\mathrm{}\{1,\mathrm{},s\}`$,
(61)
$$\begin{array}{c}\left|d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))\frac{1}{s}d_{}(f\left(y+\mathrm{sgn}(\pi )se_j\right),f(y))\right|\hfill \\ \hfill \sqrt{2\eta (2s1)^nsn2^{sn}}2^{2sn}\sqrt{\eta }.\end{array}$$
###### Claim 6.5.
For every $`\epsilon ,\delta \{1,1\}^n`$ and every $`x_m^n`$, such that $`x+\epsilon \{0,1,\mathrm{},s1\}^n`$,
$$\left|d_{}(f(x+\epsilon ),f(x))d_{}(f(x+\delta ),f(x))\right|2\sqrt{\eta }2^{2sn}.$$
###### Proof.
If $`\epsilon =\delta `$ then there is nothing to prove, so assume that $`\epsilon _{\mathrm{}}=\delta _{\mathrm{}}`$. Denote $`S=\{j\{1,\mathrm{},n\}:\epsilon _j=\delta _j\}`$ and define $`\theta ,\tau \{1,1\}^n`$ by
$$\theta _j=\{\begin{array}{cc}\epsilon _{\mathrm{}}\hfill & j=\mathrm{}\hfill \\ \epsilon _j\hfill & jS\{\mathrm{}\}\hfill \\ 1\hfill & jS\hfill \end{array}\mathrm{and}\tau _j=\{\begin{array}{cc}\epsilon _{\mathrm{}}\hfill & j=\mathrm{}\hfill \\ \epsilon _j\hfill & jS\{\mathrm{}\}\hfill \\ 1\hfill & jS.\hfill \end{array}$$
Consider the following path $`\pi `$ in $`\mathrm{๐๐ข๐๐ }(_m^n)`$: Start at $`x+\epsilon \{0,1,\mathrm{},s1\}^n`$, go in direction $`\epsilon `$ (i.e. pass to $`x`$), then go in direction $`\delta `$ (i.e. pass to $`x+\delta `$), then go in direction $`\theta `$ (i.e. pass to $`x+\delta +\theta `$), then go in direction $`\tau `$ (i.e. pass to $`x+\delta +\theta +\tau `$), and repeat this process $`s/4`$ times. It is clear from the construction that $`\pi ๐ข_{\mathrm{}}^\epsilon _{\mathrm{}}(x+\epsilon )`$. Thus, by (61) we get that
$$\begin{array}{c}\left|d_{}(f(x+\epsilon ),f(x))d_{}(f(x+\delta ),f(x))\right|\hfill \\ \hfill =\left|d_{}(f(\pi _1),f(\pi _0))d_{}(f(\pi _2),f(\pi _1))\right|2\sqrt{\eta }2^{2sn}.\end{array}$$
###### Corollary 6.6.
There exists a number $`A1`$ such that for every $`\epsilon \{1,1\}^n`$,
$$\left(14\sqrt{\eta }2^{2sn}\right)Ad_{}(f(\epsilon ),f(0))\left(1+4\sqrt{\eta }2^{2sn}\right)A.$$
###### Proof.
Denote $`e=_{j=1}^ne_j=(1,1,\mathrm{},1)`$ and take
$$A=\left(๐ผ_\delta \left[d_{}(f(\delta ),f(0))^2\right]\right)^{1/2}.$$
By (60), $`A1`$. By Claim 6.5 we know that for every $`\epsilon ,\delta \{1,1\}^{2^s}`$,
$$d_{}(f(\epsilon ),f(0))d_{}(f(e),f(0))+2\sqrt{\eta }2^{2sn}d_{}(f(\delta ),f(0))+4\sqrt{\eta }2^{2sn}.$$
Averaging over $`\delta `$, and using the Cauchy-Schwartz inequality, we get that
$`d_{}(f(\epsilon ),f(0))`$ $``$ $`\left(๐ผ_\delta \left[d_{}(f(\delta ),f(0))^2\right]\right)^{1/2}+4\sqrt{\eta }2^{2sn}`$
$`=`$ $`A+4\sqrt{\eta }2^{2sn}\left(1+4\sqrt{\eta }2^{2sn}\right)A.`$
In the reverse direction we also know that
$$A^2=๐ผ_\delta [d_{}(f(\delta ),f(0))^2]\left[d_{}(f(\epsilon ),f(0))+4\sqrt{\eta }2^{2sn}\right]^2,$$
which implies the required result since $`A1`$. โ
###### Claim 6.7.
Denote
(62) $`V=\{x_m^n:j0x_j{\displaystyle \frac{s}{2}}\mathrm{and}x_j\mathrm{is}\mathrm{even}\}.`$
Then the following assertions hold true:
1. For every $`x,yV`$ there is some $`z\{x,y\}`$, $`j\{1,\mathrm{},n\}`$, and a path $`\pi ๐ข_j^+(z)`$ of length $`s`$ which goes through $`x`$ and $`y`$. Moreover, we can ensure that if $`\pi =(\pi _0,\mathrm{},\pi _s)`$ then for all $`\mathrm{}\{1,\mathrm{},s\}`$, $`\{\pi _{\mathrm{}},\pi _\mathrm{}1\}\{0,\mathrm{},s1\}^n\mathrm{}`$.
2. For every $`x,yV`$, $`d_{\mathrm{๐๐ข๐๐ }(_m^n)}(x,y)=d_{_m^n}(x,y)=xy_{\mathrm{}}`$.
###### Proof.
Let $`j\{1,\mathrm{},n\}`$ be such that $`|y_jx_j|=xy_{\mathrm{}}:=t`$. Without loss of generality $`y_jx_j`$. We will construct a path of length $`s`$ in $`๐ข_j^+(x)`$ which goes through $`y`$. To begin with, we define $`\epsilon ^{\mathrm{}},\delta ^{\mathrm{}}\{1,1\}^n`$ inductively on $`\mathrm{}`$ as follows:
$`\epsilon _r^{\mathrm{}}`$ $`=`$ $`\{\begin{array}{cc}1\hfill & x_r+2_{k=1}^\mathrm{}1(\epsilon _r^k+\delta _r^k)<y_r\hfill \\ 1\hfill & x_r+2_{k=1}^\mathrm{}1(\epsilon _r^k+\delta _r^k)>y_r\hfill \\ 1\hfill & x_r+2_{k=1}^\mathrm{}1(\epsilon _r^k+\delta _r^k)=y_r\hfill \end{array}`$
$`\delta _r^{\mathrm{}}`$ $`=`$ $`\{\begin{array}{cc}1\hfill & x_r+2_{k=1}^\mathrm{}1(\epsilon _r^k+\delta _r^k)<y_r\hfill \\ 1\hfill & x_r+2_{k=1}^\mathrm{}1(\epsilon _r^k+\delta _r^k)>y_r\hfill \\ 1\hfill & x_r+2_{k=1}^\mathrm{}1(\epsilon _r^k+\delta _r^k)=y_r.\hfill \end{array}`$
If we define $`a_{\mathrm{}}=x+_{k=1}^{\mathrm{}}\epsilon ^k+_{k=1}^\mathrm{}1\delta ^k`$ and $`b_{\mathrm{}}=x+_{k=1}^{\mathrm{}}\epsilon ^k+_{k=1}^{\mathrm{}}\delta ^k`$ then the sequence
$$(x,a_1,b_1,a_2,b_2,\mathrm{},a_{t/21},b_{t/2}=y)$$
is a path of length $`t`$ in $`\mathrm{๐๐ข๐๐ }(_m^n)`$ joining $`x`$ and $`y`$. This proves the second assertion above. We extend this path to a path of length $`s`$ (in $`\mathrm{๐๐ข๐๐ }(_m^n)`$) from $`x`$ to $`x+se_j`$ as follows. Observe that for every $`1\mathrm{}t/2`$, $`\epsilon _j^{\mathrm{}}=\delta _j^{\mathrm{}}=1`$. Thus $`\epsilon ^{\mathrm{}}+2e_j,\delta ^{\mathrm{}}+2e_j\{1,1\}^n`$. If we define $`c_{\mathrm{}}=y+_{k=1}^{\mathrm{}}(\epsilon ^k+2e_j)+_{k=1}^\mathrm{}1(\delta ^k+2e_j)`$ and $`d_{\mathrm{}}=y+_{k=1}^{\mathrm{}}(\epsilon ^k+2e_j)+_{k=1}^{\mathrm{}}(\delta ^k+2e_j)`$, then $`d_{t/2}=x+2te_j`$. Observe that by the definition of $`V`$, $`2ts`$, and $`s2t`$ is even. Thus we can continue the path from $`x+2te_j`$ to $`x+se_j`$ by alternatively using the directions $`e_j+_\mathrm{}je_{\mathrm{}}`$ and $`e_j_\mathrm{}je_{\mathrm{}}`$. โ
###### Corollary 6.8.
Assume that $`xV`$. Then for $`A`$ as in Corollary 6.6, we have for all $`\epsilon \{1,1\}^n`$,
$$\left(110\sqrt{\eta }2^{2sn}\right)Ad_{}(f(x+\epsilon ),f(x))\left(1+10\sqrt{\eta }2^{2sn}\right)A.$$
###### Proof.
By Claim 6.7 (and its proof), there exist $`j\{1,\mathrm{},n\}`$ and $`\pi ๐ข_j^+(0)`$ such that $`\pi _1=e=(1,\mathrm{},1)`$ and for some $`k\{1,\mathrm{},s\}`$, $`\pi _k=x`$. Now, by (61) we have that
$$\left|d_{}(f(e),f(0))d_{}(f(\pi _{k1}),f(x))\right|2\sqrt{\eta }2^{2sn}.$$
Observe that since $`xV`$, $`x+e\{0,\mathrm{},s1\}^n`$. Thus by Claim 6.5
$`\left|d_{}(f(x+\epsilon ),f(x))d_{}(f(e),f(0))\right|`$
$`\left|d_{}(f(e),f(0))d_{}(f(\pi _{k1}),f(x))\right|`$
$`+\left|d_{}(f(\pi _{k1}),f(x))d_{}(f(x+e),f(x))\right|`$
$`+\left|d_{}(f(x+\epsilon ),f(x))d_{}(f(x+e),f(x))\right|`$
$`6\sqrt{\eta }2^{2sn},`$
so that the required inequalities follow from Corollary 6.6. โ
###### Corollary 6.9.
For every distinct $`x,yV`$,
$$\left(112\sqrt{\eta }2^{2sn}\right)A\frac{d_{}(f(x),f(y))}{xy_{\mathrm{}}}\left(1+12\sqrt{\eta }2^{2sn}\right)A,$$
where $`A`$ is as in Corollary 6.6.
###### Proof.
Denote $`t=xy_{\mathrm{}}`$; we may assume that there exists $`j\{1,\mathrm{},n\}`$ such that $`y_jx_j=t`$. By Claim 6.7 there is a path $`\pi ๐ข_j^+(x)`$ of length $`s`$ such that $`\pi _t=y`$. By (61) and Corollary 6.8 we have for every $`\mathrm{}\{1,\mathrm{},s\}`$
$`\left|d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1)){\displaystyle \frac{1}{s}}d_{}(f\left(x+se_j\right),f(x))\right|\sqrt{\eta }2^{2sn},`$
and
$$(110\sqrt{\eta }2^{2sn})Ad_{}(f(\pi _0),f((\pi _1))(1+10\sqrt{\eta }2^{2sn})A.$$
Thus, for all $`\mathrm{}\{1,\mathrm{},s\}`$,
$$\left(112\sqrt{\eta }2^{2sn}\right)Ad_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))\left(1+12\sqrt{\eta }2^{2sn}\right)A.$$
Thus
$`d_{}(f(x),f(y))`$ $``$ $`{\displaystyle \underset{\mathrm{}=1}{\overset{t}{}}}d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))t\left(1+12\sqrt{\eta }2^{2sn}\right)A`$
$`=`$ $`xy_{\mathrm{}}\left(1+12\sqrt{\eta }2^{2sn}\right)A.`$
On the other hand
$`d_{}(f(x),f(y))`$ $``$ $`d_{}(f(x+se_j),f(x))d_{}(f(x+se_j),f(y))`$
$``$ $`sd_{}(f(x),f(\pi _1))s\sqrt{\eta }2^{2sn}{\displaystyle \underset{\mathrm{}=t+1}{\overset{s}{}}}d_{}(f(\pi _{\mathrm{}}),f(\pi _\mathrm{}1))`$
$``$ $`s\left(110\sqrt{\eta }2^{2sn}\right)As\sqrt{\eta }2^{2sn}`$
$`\left(st\right)\left(112\sqrt{\eta }2^{2sn}\right)A`$
$``$ $`xy_{\mathrm{}}\left(112\sqrt{\eta }2^{2sn}\right)A.`$
This concludes the proof of Theorem 6.4, since the mapping $`xx/2`$ is a distortion $`1`$ bijection between $`(V,d_{_m^n})`$ and $`[s/4]_{\mathrm{}}^n`$. โ
We are now in position to prove Theorem 1.5.
###### Proof of Theorem 1.5.
We assume that $`\mathrm{\Gamma }_q^{(2)}()=\mathrm{}`$ for all $`q<\mathrm{}`$. By Lemma 6.3 it follows that for every two integers $`n,s>1`$, $`(;n,s)=1`$. Now the required result follows from Theorem 6.4. โ
###### Lemma 6.10.
Let $``$ be a metric space and $`K>0`$. Fix $`q<\mathrm{}`$ and assume that $`m:=m_q^{(2)}(;n,K)<\mathrm{}`$. Then
$$c_{}\left(_m^n\right)\frac{n^{1/q}}{2K}.$$
###### Proof.
Fix a bijection $`f:_m^n`$. Then
$`{\displaystyle \frac{nm^2}{4f^1_{\mathrm{Lip}}^2}}`$ $``$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(f\left(x+{\displaystyle \frac{m}{2}}e_j\right),f(x))^2๐\mu (x)`$
$``$ $`K^2m^2n^{1\frac{2}{q}}{\displaystyle _{\{1,01\}^n}}{\displaystyle _{_m^n}}d_{}(f(x+\epsilon ),f(x))^2๐\mu (x)๐\sigma (\epsilon )`$
$``$ $`K^2m^2n^{1\frac{2}{q}}f_{\mathrm{Lip}}^2.`$
It follows that $`\mathrm{dist}(f)\frac{n^{1/q}}{2K}`$. โ
###### Corollary 6.11.
Let $``$ be a family of metric spaces and $`0<q,K,c<\mathrm{}`$. Assume that for all $`n`$, $`\mathrm{\Gamma }_q^{(2)}(;n,n^c)K`$ for every $``$. Then for every integer $`N`$,
$$๐_N()\frac{1}{2cK}\left(\frac{\mathrm{log}N}{\mathrm{log}\mathrm{log}N}\right)^{1/q}.$$
We require the following simple lemma, which shows that the problems of embedding $`[m]_{\mathrm{}}^n`$ and $`_m^n`$ are essentially equivalent.
###### Lemma 6.12.
The grid $`[m]_{\mathrm{}}^n`$ embeds isometrically into $`_{2m}^n`$. Conversely, $`_{2m}^n`$ embeds isometrically into $`[m+1]_{\mathrm{}}^{2mn}`$. Moreover, for each $`\epsilon >0`$, $`_{2m}^n`$ embeds with distortion $`1+6\epsilon `$ into $`[m+1]^{(1/\epsilon +1)n}`$.
###### Proof.
The first assertion follows by consideration of only elements of $`_{2m}^n`$ whose coordinates are at most $`m1`$. Next, the Frรฉchet embedding
$$x(d_{_{2m}}(x,0),d_{_{2m}}(x,1),\mathrm{},d_{_{2m}}(x,2m1))[m+1]_{\mathrm{}}^{2m},$$
is an isometric embedding of $`_{2m}`$. Thus $`_{2m}^n`$ embeds isometrically into $`[m+1]_{\mathrm{}}^{2mn}`$. The final assertion is proved analogously by showing that $`_{2m}`$ embeds with distortion $`1+\epsilon `$ into $`[m+1]_{\mathrm{}}^{1/\epsilon +1}`$. This is done by consideration of the embedding
$$\begin{array}{c}x(d__m(x,0),d__m(x,2\epsilon m),d__m(x,4\epsilon m),d__m(x,6\epsilon m),\mathrm{}\hfill \\ \hfill \mathrm{},d__m(x,21/\epsilon \epsilon m)),\end{array}$$
which is easily seen to have distortion at most $`1+6\epsilon `$. โ
We are now in position to prove Theorem 1.6.
###### Proof of Theorem 1.6.
We first prove the implication $`1)2)`$. Let $`Z`$ be the disjoint union of all finite subsets of members of $``$, i.e.
$$Z=\{๐ฉ:|๐ฉ|<\mathrm{}\mathrm{and},๐ฉ\}.$$
For every $`k>1`$ we define a metric $`d_k`$ on $`Z`$ by
$$d_k(x,y)=\{\begin{array}{cc}\frac{d_๐ฉ(x,y)}{\mathrm{diam}(๐ฉ)}\hfill & ,๐ฉ\mathrm{s}.\mathrm{t}.|๐ฉ|<\mathrm{}\mathrm{and}x,y๐ฉ\hfill \\ k\hfill & \mathrm{otherwise}.\hfill \end{array}$$
Clearly $`d_k`$ is a metric. Moreover, by construction, for every $`K,k>1`$,
$$q_{(Z,d_k)}^{(2)}(K)q_{}^{(2)}(K).$$
Assume for the sake of contradiction that for every $`K,k>1`$, $`q_{(Z,d_k)}^{(2)}(K)=\mathrm{}`$. In other words, for every $`q<\mathrm{}`$, and $`k1`$, $`\mathrm{\Gamma }_q^{(2)}(Z,d_k)=\mathrm{}`$. By Lemma 6.3 it follows that for every $`k1`$, and every two integers $`n,s>1`$,
$$((Z,d_k);n,s)=1.$$
Theorem 6.4 implies that $`c_{(Z,d_k)}\left([m]_{\mathrm{}}^n\right)=1`$.
By our assumption there exists a metric space $`X`$ such that $`c_{}(X):=D>1`$. Define a metric space $`X^{}=X\times \{1,2\}`$ via $`d_X^{}((x,1),(y,1))=d_X^{}((x,2),(y,2))=d_X(x,y)`$ and $`d_X^{}((x,1),(y,2))=2\mathrm{diam}(X)`$. For large enough $`s`$ we have that $`c_{[2^{s3}]_{\mathrm{}}^{2^s}}(X^{})<D`$. Thus $`c_{(Z,d_k)}(X^{})<D`$ for all $`k`$. Define
$$k=\frac{4\mathrm{diam}(X)}{\mathrm{min}_{\begin{array}{c}x,yX\\ xy\end{array}}d_X(x,y)}.$$
Then there exists a bijection $`f:X^{}(Z,d_k)`$ with $`\mathrm{dist}(f)<\mathrm{min}\{2,D\}`$. Denote $`L=f_{\mathrm{Lip}}`$.
We first claim that there exists $``$, and a finite subset $`๐ฉ`$, such that $`|f(X^{})๐ฉ|2`$. Indeed, otherwise, by the definition of $`d_k`$, for all $`x^{},y^{}X^{}`$, $`d_k(f(x^{}),f(y^{}))=k`$. Choosing distinct $`x,yX`$, we deduce that
$$k=d_k(f(x,1),f(y,1))Ld_X(x,y)L\mathrm{diam}(X),$$
and
$`k`$ $`=`$ $`d_k(f(x,1),f(y,2)){\displaystyle \frac{L}{\mathrm{dist}(f)}}d_X^{}((x,1),(y,2))`$
$`>`$ $`{\displaystyle \frac{L}{2}}2\mathrm{diam}(X)=L\mathrm{diam}(X),`$
which is a contradiction.
Thus, there exists $``$ and a finite subset $`๐ฉ`$ such that $`|f(X^{})๐ฉ|2`$. We claim that this implies that $`f(X^{})๐ฉ`$. This will conclude the proof of 1) $``$ 2), since the metric induced by $`d_k`$ on $`๐ฉ`$ is a re-scaling of $`d_๐ฉ`$, so that $`X`$ embeds with distortion smaller than $`D`$ into $`๐ฉ`$, which is a contradiction of the definition of $`D`$.
Assume for the sake of a contradiction that there exists $`x^{}X^{}`$ such that $`f(x^{})๐ฉ`$. By our assumption there are distinct $`a^{},b^{}X^{}`$ such that $`f(a^{}),f(b^{})๐ฉ`$. Now,
$$1d_k(f(a^{}),f((b^{}))\frac{L}{\mathrm{dist}(f)}d_X^{}(a^{},b^{})>\frac{L}{2}\underset{\begin{array}{c}u,vX\\ uv\end{array}}{\mathrm{min}}d_X(u,v),$$
while
$`{\displaystyle \frac{4\mathrm{diam}(X)}{\mathrm{min}_{\begin{array}{c}u,vX\\ uv\end{array}}d_X(u,v)}}=k`$ $`=`$ $`d_k(f(x^{}),f((a^{}))`$
$``$ $`Ld(x^{},a^{})L\mathrm{diam}(X^{})=2L\mathrm{diam}(X),`$
which is a contradiction.
To prove the implication $`2)3)`$ observe that in the above argument we have shown that there exists $`k,q<\mathrm{}`$ such that $`\mathrm{\Gamma }_q^{(2)}(Z,d_k)<\mathrm{}`$. It follows that for some integer $`n_0`$, $`((Z,d_k);n_0,n_0)<1`$, since otherwise by Theorem 6.4 we would get that $`(Z,d_k)`$ contains, uniformly in $`n`$, bi-Lipschitz copies of $`[n]_{\mathrm{}}^n`$. Combining Lemma 6.12 and Lemma 6.10 we arrive at a contradiction. By Lemma 6.3, the fact that $`((Z,d_k);n_0,n_0)<1`$, combined with Corollary 6.11, implies that $`๐_n(Z,d_k)=\mathrm{\Omega }((\mathrm{log}n)^\alpha )`$ for some $`\alpha >0`$. By the definition of $`(Z,d_k)`$, this implies the required result. โ
We end this section by proving Theorem 1.8:
###### Proof of Theorem 1.8.
Denote $`|X|=n`$ and
$$\mathrm{\Phi }=\frac{\mathrm{diam}(X)}{\mathrm{min}_{xy}d(x,y)}.$$
Write $`t=4\mathrm{\Phi }/\epsilon `$ and let $`s`$ be an integer divisible by $`4`$ such that $`s\mathrm{max}\{n,t\}`$. Then $`c_{[s]_{\mathrm{}}^s}(X)1+\frac{\epsilon }{4}`$. Fix a metric space $`Z`$ and assume that $`c_Z(X)>1+\epsilon `$. It follows that $`c_Z([s]_{\mathrm{}}^s)1+\frac{\epsilon }{2}`$. By Theorem 6.4 we deduce that
$$(Z,s,4s)1\frac{\epsilon ^2}{2^{s^2}}.$$
By Lemma 6.3 we have that $`m_q^{(2)}(;n,3s)8sn^{\mathrm{log}_s(4s)}`$, where $`q\frac{10^s}{\epsilon ^2}`$. Thus by Lemma 6.10 and Lemma 6.12 we see that for any integer $`n8s`$,
$$c_Z\left(\left[n^5\right]_{\mathrm{}}^n\right)\frac{n^{1/q}}{4s}=\frac{n^{\epsilon ^2/10^s}}{4s}.$$
Choosing $`N(C\gamma )^{\frac{2^{4s}}{\epsilon ^2}}`$, for an appropriate universal constant $`C`$, yields the required result. โ
## 7. Applications to bi-Lipschitz, uniform, and coarse embeddings
Let $`(๐ฉ,d_๐ฉ)`$ and $`(,d_{})`$ be metric spaces. For $`f:๐ฉ`$ and $`t>0`$ we define
$$\mathrm{\Omega }_f(t)=sup\{d_{}(f(x),f(y));d_๐ฉ(x,y)t\},$$
and
$$\omega _f(t)=inf\{d_{}(f(x),f(y));d_๐ฉ(x,y)t\}.$$
Clearly $`\mathrm{\Omega }_f`$ and $`\omega _f`$ are nondecreasing, and for every $`x,y๐ฉ`$,
$$\omega _f\left(d_๐ฉ(x,y)\right)d_{}(f(x),f(y))\mathrm{\Omega }_f\left(d_๐ฉ(x,y)\right)$$
With these definitions, $`f`$ is uniformly continuous if $`lim_{t0}\mathrm{\Omega }_f(t)=0`$, and $`f`$ is a uniform embedding if $`f`$ is injective and both $`f`$ and $`f^1`$ are uniformly continuous. Also, $`f`$ is a coarse embedding if $`\mathrm{\Omega }_f(t)<\mathrm{}`$ for all $`t>0`$ and $`lim_t\mathrm{}\omega _f(t)=\mathrm{}`$.
###### Lemma 7.1.
Let $`(,d_{})`$ be a metric space, $`n`$ an integer, $`\mathrm{\Gamma }>0`$, and $`0<pqr`$. Then for every function $`f:\mathrm{}_r^n`$, and every $`s>0`$,
$$n^{1/q}\omega _f(2s)\mathrm{\Gamma }m_q^{(p)}(;n,\mathrm{\Gamma })\mathrm{\Omega }_f\left(\frac{2\pi sn^{1/r}}{m_q^{(p)}(;n,\mathrm{\Gamma })}\right).$$
###### Proof.
Denote $`m=m_q^{(p)}(;n,\mathrm{\Gamma })`$, and define $`g:_m^n`$ by
$$g(x_1,\mathrm{},x_n)=f\left(\underset{j=1}{\overset{n}{}}se^{\frac{2\pi ix_j}{m}}e_j\right).$$
Then
$$\begin{array}{c}_{\{1,0,1\}^n}_{_m^n}d_{}(g(x+\epsilon ),g(x))^p๐\mu (x)๐\sigma (\epsilon )\hfill \\ \hfill \underset{\epsilon \{1,0,1\}^n}{\mathrm{max}}\mathrm{\Omega }_f\left(s\left(\underset{j=1}{\overset{n}{}}\left|e^{\frac{2\pi i\epsilon _j}{m}}1\right|^r\right)^{1/r}\right)^p\mathrm{\Omega }_f\left(\frac{2\pi sn^{1/r}}{m}\right)^p.\end{array}$$
On the other hand,
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _{_m^n}}d_{}(g\left(x+{\displaystyle \frac{m}{2}}e_j\right),g(x))^p๐\mu (x)n\omega _f(2s)^p.`$
By the definition of $`m_q^{(p)}(;n,\mathrm{\Gamma })`$ it follows that
$$n\omega _f(2s)^p\mathrm{\Gamma }^pm^pn^{1\frac{p}{q}}\mathrm{\Omega }_f\left(\frac{2\pi sn^{1/r}}{m}\right)^p,$$
as required. โ
###### Corollary 7.2.
Let $``$ be a metric space and assume that there exist constants $`c,\mathrm{\Gamma }>0`$ such that for infinitely many integers $`n`$, $`m_q^{(p)}(;n,\mathrm{\Gamma })cn^{1/q}`$. Then for every $`r>q`$, $`\mathrm{}_r`$ does not uniformly or coarsely embed into $``$.
###### Proof.
To rule out the existence of a coarse embedding choose $`s=n^{\frac{1}{q}\frac{1}{r}}`$ in Lemma 7.1. Using Lemma 2.3 we get that
$$\omega _f\left(2n^{\frac{1}{q}\frac{1}{r}}\right)c\mathrm{\Gamma }\mathrm{\Omega }_f\left(2\pi \mathrm{\Gamma }\right).$$
Since $`q<r`$, it follows that $`lim\; inf_t\mathrm{}\omega _f(t)<\mathrm{}`$, so that $`f`$ is not a coarse embedding.
To rule out the existence of a uniform embedding, assume that $`f:\mathrm{}_rX`$ is invertible and $`f^1`$ is uniformly continuous. Then there exists $`\delta >0`$ such that for $`x,y\mathrm{}_r`$, if $`d_{}(f(x),f(y))<\delta `$ then $`xy_r<2`$. It follows that $`\omega _f(2)\delta `$. Choosing $`s=1`$ in Lemma 7.1, and using Lemma 2.3, we get that
$$0<\delta \omega _f(2)c\mathrm{\Gamma }\mathrm{\Omega }_f\left(2\pi \mathrm{\Gamma }n^{\frac{1}{r}\frac{1}{q}}\right).$$
Since $`r>q`$ it follows that $`lim\; sup_{t0}\mathrm{\Omega }_f(t)>0`$, so that $`f`$ is not uniformly continuous. โ
The following corollary contains Theorem 1.9, Theorem 1.10 and Theorem 1.11.
###### Corollary 7.3.
Let $`X`$ be a $`K`$-convex Banach space. Assume that $`Y`$ is a Banach space which coarsely or uniformly embeds into $`X`$. Then $`q_Yq_X`$. In particular, for $`p,q>0`$, $`L_p`$ embeds uniformly or coarsely into $`L_q`$ if and only if $`pq`$ or $`qp2`$.
###### Proof.
By the Maurey-Pisier theorem , for every $`\epsilon >0`$ and every $`n`$, $`Y`$ contains a $`(1+\epsilon )`$ distorted copy of $`\mathrm{}_{q_Y}^n`$. By Theorem 4.1, since $`X`$ is $`K`$-convex, for every $`q>q_X`$ there exists $`\mathrm{\Gamma }<\mathrm{}`$ such that $`m_q(;n,\mathrm{\Gamma })=O\left(n^{1/q}\right)`$. Thus, by the proof of Corollary 7.2, if $`Y`$ embeds coarsely or uniformly into $`X`$ then $`q_Yq`$, as required.
The fact that if $`pq`$ then $`L_p`$ embeds coarsely and uniformly into $`L_q`$ follows from the fact that in this case $`L_p`$, equipped with the metric $`xy_p^{p/q}`$, embeds isometrically into $`L_q`$ (for $`pq2`$ this is proved in , . For the remaining cases see Remark 5.10 in ). If $`2pq`$ then $`L_p`$ is linearly isometric to a subspace of $`L_q`$ (see e.g. ). It remains to prove that if $`p>q`$ and $`p>2`$ then $`L_p`$ does not coarsely or uniformly embed into $`L_q`$. We may assume that $`q2`$, since for $`q2`$, $`L_q`$ embeds coarsely and uniformly into $`L_2`$. But, now the required result follows from the fact that $`L_q`$ is $`K`$ convex and $`q_{L_q}=q`$, $`q_{L_p}=p`$ (see ). โ
We now pass to the proof of Theorem 1.12. Before doing so we remark that Theorem 1.12 is almost optimal in the following sense. The identity mapping embeds $`[m]_{\mathrm{}}^n`$ into $`\mathrm{}_q^n`$ with distortion $`n^{1/q}`$. By the Maurey-Pisier theorem , $`Y`$ contains a copy of $`\mathrm{}_{q_Y}^n`$ with distortion $`1+\epsilon `$ for every $`\epsilon >0`$. Thus $`c_Y([m]_{\mathrm{}}^n)n^{1/q_Y}`$. Additionally, $`[m]_{\mathrm{}}^n`$ is $`m`$-equivalent to an equilateral metric. Thus, if $`Y`$ is infinite dimensional then $`c_Y([m]_{\mathrm{}}^n)m`$. It follows that
$$c_Y([m]_{\mathrm{}}^n)\mathrm{min}\{n^{1/q_Y},m\}.$$
###### Proof of Theorem 1.12.
Assume that $`m`$ is divisible by $`4`$ and
$$m\frac{2n^{1/q}}{C_q(Y)K(Y)}.$$
By Theorem 4.1, for every $`f:_m^nY`$,
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}f\left(x+\frac{m}{2}\right)f(x)_Y^q๐\mu (x)\hfill \\ \hfill \left[15C_q(Y)K(Y)\right]^qm^q_{\{1,0,1\}^n}_{_m^n}f(x+\epsilon )f(x)_Y^q๐\mu (x)๐\sigma (\epsilon ).\end{array}$$
Thus, assuming that $`f`$ is bi-Lipschitz we get that
$$\frac{nm^q}{2^qf^1_{\mathrm{Lip}}^q}\left[15C_q(Y)K(Y)\right]^qm^qf_{\mathrm{Lip}}^q,$$
i.e.
$$\mathrm{dist}(f)\frac{n^{1/q}}{30C_q(Y)K(Y)}.$$
By Lemma 6.12 this shows that for $`m\frac{2n^{1/q}}{C_q(Y)K(Y)}`$, such that $`m`$ is divisible by $`4`$, $`c_Y([m]_{\mathrm{}}^n)=\mathrm{\Omega }\left(n^{1/q}\right)`$. If $`m<\frac{2n^{1/q}}{C_q(Y)K(Y)}`$ then the required lower bound follows from the fact that $`[m]_{\mathrm{}}^n`$ contains an isometric copy of $`[m_1]_{\mathrm{}}^{n_1}`$, where $`m_1`$ is an integer divisible by $`4`$, $`m_1\frac{2n_1^{1/q}}{C_q(Y)K(Y)}`$, and $`m_1=\mathrm{\Theta }(m)`$, $`n_1=\mathrm{\Theta }(m^q)`$. Passing to integers $`m`$ which are not necessarily divisible by $`4`$ is just as simple. โ
###### Remark 7.4.
Similar arguments yield bounds on $`c_Y([m]_p^n)`$, which strengthen the bounds in .
###### Remark 7.5.
Although $`L_1`$ is not $`K`$-convex, we can still show that
$$c_1([m]_{\mathrm{}}^n)=\mathrm{\Theta }\left(\mathrm{min}\{\sqrt{n},m\}\right).$$
This is proved as follows. Assume that $`f:_m^nL_1`$ is bi-Lipschitz. If $`m`$ is divisible by $`4`$, and $`m\pi \sqrt{n}`$, then the fact that $`L_1`$, equipped with the metric $`\sqrt{xy_1}`$, is isometric to a subset of Hilbert space , , together with Proposition 3.1, shows that
$$\begin{array}{c}\underset{j=1}{\overset{n}{}}_{_m^n}f\left(x+\frac{m}{2}\right)f(x)_1๐\mu (x)\hfill \\ \hfill m^2_{\{1,0,1\}^n}_{_m^n}f(x+\epsilon )f(x)_1๐\mu (x)๐\sigma (\epsilon ).\end{array}$$
Arguing as in the proof of Theorem 1.12, we see that for $`m\sqrt{n}`$, $`c_1([m]_{\mathrm{}}^n)=\mathrm{\Omega }\left(\sqrt{n}\right)`$. This implies the required result, as in the proof of Theorem 1.12.
## 8. Discussion and open problems
1. Perhaps the most important open problem related to the nonlinear cotype inequality on Banach spaces is whether for every Banach space $`X`$ with cotype $`q<\mathrm{}`$, for every $`1pq`$ there is a constant $`\mathrm{\Gamma }<\mathrm{}`$ such that $`m_q^{(p)}(X;n,\mathrm{\Gamma })=O\left(n^{1/q}\right)`$. By Lemma 2.3 this is best possible. In Theorem 4.1 we proved that this is indeed the case when $`X`$ is $`K`$-convex, while our proof of Theorem 1.2 only gives $`m_q^{(p)}(X;n,\mathrm{\Gamma })=O\left(n^{2+1/q}\right)`$.
2. $`L_1`$ is not $`K`$-convex, yet we do know that $`m_2^{(1)}(L_1;n,4)=O\left(\sqrt{n}\right)`$. This follows directly from Remark 7.5, Lemma 2.4 and Lemma 2.5. It would be interesting to prove the same thing for $`m_2(L_1;n,\mathrm{\Gamma })`$.
3. We conjecture that the $`K`$-convexity assumption in Theorem 1.9 and Theorem 1.11 is not necessary. Since $`L_1`$ embeds coarsely and uniformly into $`L_2`$, these theorems do hold for $`L_1`$. It seems to be unknown whether any Banach space with finite cotype embeds uniformly or coarsely into a $`K`$-convex Banach space. The simplest space for which we do not know the conclusion of these theorems is the Schatten trace class $`C_1`$ (see . In it is shown that this space has cotype $`2`$). The fact that $`C_1`$ does not embed uniformly into Hilbert space follows from the results of , together with , . For more details we refer to the discussion in (a similar argument works for coarse embeddings of $`C_1`$ into Hilbert space, by use of ). We remark that the arguments presented here show that a positive solution of the first problem stated above would yield a proof of Theorem 1.9 and Theorem 1.11 without the $`K`$-convexity assumption.
## 9. Acknowledgments
We are grateful to Keith Ball for several valuable discussions. We also thank Yuri Rabinovich for pointing out the connection to Matouลกekโs BD Ramsey theorem. |
warning/0506/quant-ph0506241.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Quantum entanglement has been identified as an important resource for quantum computing, quantum communication, and other applications . A fundamental theoretical problem is to understand the types of entanglement that composite quantum systems can achieve.
Defining entanglement types as equivalence classes of quantum states under local unitary (LU) equivalence is perhaps the most natural way to proceed in classifying entanglement . The group of LU transformations acts on the space of quantum states, partitioning it into LU orbits. Each orbit is a collection of locally equivalent quantum states that forms a differentiable manifold with a certain dimension.
The classification of entanglement types has turned out to be a difficult problem. Most of the progress in understanding multipartite entanglement has been for systems of only two or three qubits . Few results exist concerning the classification of $`n`$-qubit entanglement types for arbitrary $`n`$.
A promising approach to the difficult problem of characterizing entanglement types is to break the problem into two parts. First, identify the possible dimensions of LU orbits in the state space. Then, identify the orbits that have each possible dimension. In , the present authors addressed the first part of this program by identifying the allowable range of orbit dimensions for $`n`$-qubit quantum states to be $`3n/2`$ to $`3n`$ for even $`n`$ and $`(3n+1)/2`$ to $`3n`$ for odd $`n`$. The present paper begins the second part of the program by completely identifying all orbits (entanglement types) with minimum dimension.
In this paper, we identify all $`n`$-qubit entanglement types that have minimum orbit dimension. States that have the minimum orbit dimension are, in some sense, the โrarestโ quantum states. We show that the only quantum states to achieve minimum orbit dimension are tensor products of singlet states (with a single unentangled qubit for $`n`$ odd) and their LU equivalents. This suggests a special role for the 2-qubit singlet state in the theory of $`n`$-qubit quantum entanglement.
## 2 Notation and previous results
In this section we establish notation, some definitions, and state some results from our previous paper needed for the present paper. For the convenience of the reader, we give a list (Appendix A) of equation and statement numbers in the present paper with their matching numbers in .
Let $`|0`$, $`|1`$ denote the standard computational basis for $`^2`$ and write $`|i_1i_2\mathrm{}i_n`$ for $`|i_1|i_2\mathrm{}|i_n`$ in $`(^2)^n`$. For a multi-index $`I=(i_1i_2\mathrm{}i_n)`$ with $`i_k=0,1`$ for $`1kn`$, we will write $`|I`$ to denote $`|i_1i_2\mathrm{}i_n`$. Let $`i_k^c`$ denote the bit complement
$$i_k^c=\{\begin{array}{cc}0\hfill & \text{if }i_k=1\hfill \\ 1\hfill & \text{if }i_k=0\hfill \end{array}$$
and let $`I_k`$ denote the multi-index
$$I_k:=(i_1i_2\mathrm{}i_{k1}i_k^ci_{k+1}\mathrm{}i_n)$$
obtained from $`I`$ by taking the complement of the $`k`$th bit for $`1kn`$. Similarly, let $`I_{kl}`$ denote the multi-index
$$I_{kl}:=(i_1i_2\mathrm{}i_{k1}i_k^ci_{k+1}\mathrm{}i_{l1}i_l^ci_{l+1}\mathrm{}i_n)$$
obtained from $`I`$ by taking the complement of the $`k`$th and $`l`$th bits for $`1k<ln`$.
Let $`H=(^2)^n`$ denote the Hilbert space for a system of $`n`$ qubits and let $`(H)`$ denote the projectivization of $`H`$ which is the state space of the system. We take the local unitary group to be $`G=\text{SU}(2)^n`$. Its Lie algebra $`LG=\text{su}(2)^n`$ is the space of $`n`$-tuples of traceless skew-Hermitian $`2\times 2`$ matrices. We choose $`A=i\sigma _z=\left[\begin{array}{cc}i& 0\\ 0& i\end{array}\right]`$, $`B=i\sigma _y=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]`$, and $`C=i\sigma _x=\left[\begin{array}{cc}0& i\\ i& 0\end{array}\right]`$ as a basis for $`\text{su}(2)`$, where $`\sigma _x,\sigma _y,\sigma _z`$ are the Pauli spin matrices. We define elements $`A_k,B_k,C_k`$ of $`LG`$ for $`1kn`$ to have $`A,B,C`$, respectively, in the $`k`$th coordinate and zero elsewhere.
$`A_k`$ $`=`$ $`(0,\mathrm{},0,\left[\begin{array}{cc}i& 0\\ 0& i\end{array}\right],0,\mathrm{},0)`$
$`B_k`$ $`=`$ $`(0,\mathrm{},0,\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],0,\mathrm{},0)`$
$`C_k`$ $`=`$ $`(0,\mathrm{},0,\left[\begin{array}{cc}0& i\\ i& 0\end{array}\right],0,\mathrm{},0)`$
Given a state vector $`|\psi =c_I|I`$, we have the following.
$`A_k|\psi `$ $`=`$ $`{\displaystyle \underset{I}{}}i(1)^{i_k}c_I|I`$ (4)
$`B_k|\psi `$ $`=`$ $`{\displaystyle \underset{I}{}}(1)^{i_k}c_{I_k}|I`$ (5)
$`C_k|\psi `$ $`=`$ $`{\displaystyle \underset{I}{}}ic_{I_k}|I`$ (6)
Given a state vector $`|\psi H`$, we define the $`2^n\times (3n+1)`$ complex matrix $`M`$ to be the $`(3n+1)`$-tuple of column vectors
$$M=(A_1|\psi ,B_1|\psi ,C_1|\psi ,\mathrm{},A_n|\psi ,B_n|\psi ,C_n|\psi ,i|\psi ).$$
(7)
We view the matrix $`M`$ and also its column vectors both as real and also as complex via the standard identification
$`^N`$ $``$ $`^{2N}`$
$`(z_1,z_2,\mathrm{},z_N)`$ $``$ $`(a_1,b_1,a_2,b_2,\mathrm{},a_N,b_N)`$ (8)
where $`z_j=a_j+ib_j`$ for $`1jN`$. The real rank of $`M`$ gives the dimension of the $`G`$ orbit of a state.
Proposition 2.1. Orbit dimension as rank of $`M`$: Let $`x(H)`$ be a state, let $`|\psi `$ be a Hilbert space representative for $`x`$, and let $`M`$ be the associated matrix constructed from the coordinates of $`|\psi `$ as described above. Then we have
$$\text{rank}_{}M=dim๐ช_x+1.$$
We denote by $`๐`$ the set
$$๐=\{A_1|\psi ,B_1|\psi ,C_1|\psi ,\mathrm{},A_n|\psi ,B_n|\psi ,C_n|\psi ,i|\psi \}$$
of columns of $`M`$, and for $`1kn`$ we define the triple $`T_k๐`$ to be the set of vectors
$$T_k=\{A_k|\psi ,B_k|\psi ,C_k|\psi \}.$$
(9)
For a subset $`๐`$, we write $``$ to denote the real span of the column vectors contained in $``$. For convenience we write $`T_{i_1},T_{i_2},\mathrm{},T_{i_r}`$ to denote the real span $`T_{i_1}T_{i_2}\mathrm{}T_{i_r}`$ of a set of triples. We have the following facts about the dimensions of spans of sets of triples.
Proposition 2.2. Let $`T_k=\{A_k|\psi ,B_k|\psi ,C_k|\psi \}`$ be a triple of columns of $`M`$. The three vectors in the triple are orthogonal when viewed as real vectors.
Proposition 2.3. Main orthogonality proposition: Suppose that
$$dimT_{j_1},T_{j_2},\mathrm{},T_{j_m}<3m$$
for some $`1j_1<j_2<\mathrm{}<j_mn`$. Then there is a nonempty subset $`K\{j_1,j_2,\mathrm{},j_m\}`$ containing an even number of elements such that there are two orthogonal vectors $`|\zeta _k,|\eta _k`$ in $`T_k`$, both of which are orthogonal to $`i|\psi `$, $`A_j|\psi `$, $`B_j|\psi `$ and to $`C_j|\psi `$ for all $`kK,jK`$.
More can be said about a pair of triples which together span fewer than five dimensions.
Lemma 2.4. Suppose that for some $`1l<l^{}n`$ we have $`A_l|\psi =A_l^{}|\psi `$ and $`C_l|\psi =C_l^{}|\psi `$. Then $`A_k|\psi `$, $`B_k|\psi `$, and $`C_k|\psi `$ are each orthogonal to $`i|\psi `$ and to $`A_j|\psi `$, $`B_j|\psi ,C_j|\psi `$ for all $`k\{l,l^{}\},j\{l,l^{}\}`$.
Proposition 2.5. Generalization of 2: Suppose that $`dimT_l,T_l^{}4`$ for some $`1l<l^{}n`$. Then $`A_k|\psi `$, $`B_k|\psi `$, and $`C_k|\psi `$ are each orthogonal to $`i|\psi `$ and to $`A_j|\psi `$, $`B_j|\psi ,C_j|\psi `$ for all $`k\{l,l^{}\},j\{l,l^{}\}`$.
We use propositions 22 and 2 to establish the lower bound for the rank of $`M`$, and therefore also for orbit dimension.
Proposition 2.6. Minimum rank of $`M`$: Let $`|\psi H`$ be a state vector, and let $`M`$ be the matrix associated to $`|\psi `$. We have
$$\text{rank}_{}M\{\begin{array}{cc}\frac{3n}{2}+1& n\text{ even}\\ \frac{3n+1}{2}+1& n\text{ odd}\end{array}.$$
A product of singlets (together with an unentangled qubit for $`n`$ odd) achieves the lower bound established in 2. This establishes the minimum orbit dimension.
Theorem 2.7. For the local unitary group action on state space for $`n`$ qubits, the smallest orbit dimension is
$$\mathrm{min}\{dim๐ช_x:x(H)\}=\{\begin{array}{cc}\frac{3n}{2}& n\text{ even}\\ \frac{3n+1}{2}& n\text{ odd}\end{array}.$$
This concludes the statements of definitions and results from to be used in the sequel.
## 3 Further results on ranks of sets of triples
In this section we develop more facts about ranks of sets of triples of the columns of the matrix $`M`$ associated to an $`n`$-qubit state vector $`|\psi `$ as described in the previous section. These include strengthened versions of 2, 2 and 2. We begin with a general fact about local unitary invariance.
Proposition 3.1. The dimension of the span of a union of triples with or without the rightmost column $`i|\psi `$ is a local unitary invariant.
Proof. We prove the proposition for the case โwith the rightmost column.โ The same proof works for the case โwithout the rightmost columnโ by making the obvious changes.
Let $`T_{j_1},T_{j_2},\mathrm{},T_{j_m}`$ be a set of triples of columns of the matrix $`M`$ associated to the state vector $`|\psi `$. Let $`U=_{i=1}^nU_i`$ be an element of the local unitary group and let $`M^{}`$ with triples $`T_k^{}=\{A_kU|\psi ,B_kU|\psi ,C_kU|\psi \}`$ be associated to the state $`|\psi ^{}=U|\psi `$. Since $`U`$ is unitary, the dimension of the span
$$dimT_{j_1}^{}T_{j_2}^{}\mathrm{}T_{j_m}^{}\{i|\psi ^{}\}$$
is equal to the dimension of the span of the set
$$()\underset{i=1}{\overset{m}{}}\{U^{}A_{j_i}U|\psi ,U^{}B_{j_i}U|\psi ,U^{}C_{j_i}U|\psi \}\{i|\psi \}.\text{}$$
Observe that
$`U^{}A_kU`$ $`=`$ $`(0,\mathrm{},0,U_k^{}AU_k,0,\mathrm{},0)=(0,\mathrm{},0,\text{Ad}(U_k^{})(A),0,\mathrm{},0)`$
$`U^{}B_kU`$ $`=`$ $`(0,\mathrm{},0,U_k^{}BU_k,0,\mathrm{},0)=(0,\mathrm{},0,\text{Ad}(U_k^{})(B),0,\mathrm{},0)`$
$`U^{}C_kU`$ $`=`$ $`(0,\mathrm{},0,U_k^{}CU_k,0,\mathrm{},0)=(0,\mathrm{},0,\text{Ad}(U_k^{})(C),0,\mathrm{},0)`$
where the zeros occur in all but the $`k`$th coordinate and $`\text{Ad}:\text{SU}(2)\text{SO}(\text{su}(2))`$ denotes the adjoint representation of $`\text{SU}(2)`$ on its Lie Algebra. It follows that the span of the set $`()`$ lies inside the span of the set
$$()\underset{i=1}{\overset{m}{}}\{A_{j_i}|\psi ,B_{j_i}|\psi ,C_{j_i}|\psi \}\{i|\psi \}.\text{}$$
and hence that
$$dimT_{j_1}^{}T_{j_2}^{}\mathrm{}T_{j_m}^{}\{i|\psi ^{}\}dimT_{j_1}T_{j_2}\mathrm{}T_{j_m}\{i|\psi \}.$$
Reversing the roles of $`|\psi `$ and $`|\psi ^{}`$ yields that the span of the set $`()`$ lies inside the span of the set $`()`$, and hence that
$$dimT_{j_1}^{}T_{j_2}^{}\mathrm{}T_{j_m}^{}\{i|\psi ^{}\}dimT_{j_1}T_{j_2}\mathrm{}T_{j_m}\{i|\psi \}.$$
This concludes the proof. $`\mathrm{}`$
Next we consider the case of two triples which together span five dimensions.
Proposition 3.2. Suppose that $`dimT_l,T_l^{}=5`$ for some $`1l<l^{}n`$. Then there are four independent vectors $`|\zeta _l,|\eta _lT_l`$, $`|\zeta _l^{},|\eta _l^{}T_l^{}`$, which are orthogonal to $`i|\psi `$ and to $`A_j|\psi `$, $`B_j|\psi ,C_j|\psi `$ for all $`j\{l,l^{}\}`$.
Proof. Applying the main orthogonality proposition 2, the only new claim made in the statement of 3 is that the vectors $`|\zeta _k,|\eta _kT_k`$, $`k\{l,l^{}\}`$ are independent. In the proof of 2 there is a unitary transformation $`U:HH`$ such that
$`U|\zeta _k`$ $`=`$ $`B_kU|\psi `$
$`U|\eta _k`$ $`=`$ $`C_kU|\psi `$
for $`k\{l,l^{}\}`$ and $`A_lU|\psi `$ is collinear with $`A_l^{}U|\psi `$. Since $`dimT_l,T_l^{}=5`$, the span of $`\{A_kU|\psi `$, $`B_kU|\psi ,C_kU|\psi :k=l,l^{}\}`$ is also five dimensional by 3. Thus the collinearity of $`A_lU|\psi `$ and $`A_l^{}U|\psi `$ implies that $`U|\zeta _l`$, $`U|\eta _l`$, $`U|\zeta _l^{}`$, and $`U|\eta _l^{}`$ are independent, and therefore that $`|\zeta _l`$, $`|\eta _l`$, $`|\zeta _l^{}`$, and $`|\eta _l^{}`$ are independent. $`\mathrm{}`$
Next, a small observation proves a stronger version of 2.
Lemma 3.3. Suppose that for some $`1l<l^{}n`$ we have $`A_l|\psi =A_l^{}|\psi `$ and $`C_l|\psi =C_l^{}|\psi `$. Then $`B_l|\psi =B_l^{}|\psi `$, the dimension of $`T_l,T_l^{}`$ is three, and $`T_l,T_l^{}`$ is orthogonal to $`i|\psi `$ and to $`A_j|\psi `$, $`B_j|\psi ,C_j|\psi `$ for all $`j\{l,l^{}\}`$.
Proof. We only need to prove that $`B_l|\psi =B_l^{}|\psi `$. The rest of the statement follows from 2. Equation (4) and the hypothesis $`A_l|\psi =A_l^{}|\psi `$ imply that $`(1)^{i_l}c_I=(1)^{i_l^{}}c_I`$, so if $`c_I0`$ then $`i_l=i_l^{}`$ mod 2. It follows that if $`c_{I_l}0`$ then $`i_l=i_l^{}+1`$ mod 2. Equation (6) and the hypothesis $`C_l|\psi =C_l^{}|\psi `$ imply that $`c_{I_l}=c_{I_l^{}}`$ for all $`I`$. Hence, for all $`I`$ we have
$$I\left|B_l\right|\psi =(1)^{i_l}c_{I_l}=(1)^{i_l^{}+1}c_{I_l^{}}=I\left|B_l^{}\right|\psi .$$
This establishes the claim. $`\mathrm{}`$
The strengthened Lemma 3 yields the following strengthened version of 2. We omit the proof because it requires only a minor change to apply 3 in the proof of 2.
Proposition 3.4. Strengthened version of 2: Let $`|\psi `$ be a state vector and let $`M`$ be its associated matrix. Suppose that $`dimT_l,T_l^{}4`$ for some $`1l<l^{}n`$. Then $`|\psi `$ is local unitary equivalent to a state vector $`|\psi ^{}`$ such that $`A_l|\psi ^{}=A_l^{}|\psi ^{}`$, $`B_l|\psi ^{}=B_l^{}|\psi ^{}`$, $`C_l|\psi ^{}=C_l^{}|\psi ^{}`$, the dimension of $`T_l,T_l^{}`$ is three, and $`T_l,T_l^{}`$ is orthogonal to $`i|\psi `$ and to $`A_j|\psi `$, $`B_j|\psi ,C_j|\psi `$ for all $`j\{l,l^{}\}`$.
Next we state and prove a stronger version of 2.
Proposition 3.5. Minimum rank for submatrices of $`M`$: Let $`๐ฎ=T_{i_1}T_{i_2}\mathrm{}T_{i_q}\{i|\psi \}`$ be a union of $`q`$ triples together with the rightmost column $`i|\psi `$ of $`M`$. Then
$$dim๐ฎ\{\begin{array}{cc}\frac{3q}{2}+1& q\text{ even}\\ \frac{3q+1}{2}+1& q\text{ odd}\end{array}.$$
Proof. Let $`๐ฎ_0๐ฎ`$ be a union of some number $`p`$ of triples, maximal with respect to the property that $`๐ฎ_0`$ contains a subspace $`W`$ for which
1. $`dimW\{\begin{array}{cc}\frac{3p}{2}& p\text{ even}\\ \frac{3p+1}{2}& p\text{ odd}\end{array}`$, and
2. $`W๐๐ฎ_0`$.
We separate the argument into cases. We show that in every case, either 3 holds or we can derive a contradiction by constructing a superset $`๐ฎ_1`$ such that $`๐ฎ_0๐ฎ_1๐ฎ`$, and $`๐ฎ_1`$ is the union of some number $`p^{}>p`$ triples and contains a subspace $`W^{}`$ satisfying properties (i) and (ii) with $`p^{}`$ in place of $`p`$. The construction of $`๐ฎ_1`$ violates the maximality of $`๐ฎ_0`$ and therefore rules out the case in question.
Case 1: Suppose that $`p=q`$. Then 3 holds.
Case 2: Suppose that $`p<q`$ and that the triples $`T_{j_1},T_{j_2},\mathrm{},T_{j_{qp}}`$ in $`๐ฎ๐ฎ_0`$ have the maximum possible span, that is,
$$dimT_{j_1},T_{j_2},\mathrm{},T_{j_{qp}}=3(qp).$$
Properties (i) and (ii) imply that
$`dim๐ฎ`$ $``$ $`dimW+dim๐ฎ๐ฎ_0`$
$``$ $`{\displaystyle \frac{3p}{2}}+3(qp)`$
$`=`$ $`{\displaystyle \frac{6q3p}{2}}`$
$``$ $`{\displaystyle \frac{6q(3q3)}{2}}\text{(since }pq1\text{)}`$
$`=`$ $`{\displaystyle \frac{3q+3}{2}}`$
$`=`$ $`{\displaystyle \frac{3q+1}{2}}+1`$
and so 3 holds. Note that if $`p=q1`$, the hypothesis of full span is met by 2. Therefore in the remaining cases we need only consider $`pq2`$.
Case 3: Suppose $`pq2`$ and that there is a pair of triples $`T_l,T_l^{}`$ in $`๐ฎ๐ฎ_0`$ with $`dimT_l,T_l^{}4`$. Let $`๐ฎ_1=๐ฎ_0T_lT_l^{}`$, let $`p^{}=p+2`$, and let $`W^{}=WT_lT_l^{}`$, where โ$``$โ denotes the orthogonal direct sum. That the sum is orthogonal is guaranteed by property (ii) for $`W`$. Proposition 2 implies that property (ii) also holds for the pair $`(๐ฎ_1,W^{})`$ and that $`dimW^{}dimW+3`$ (in fact, we have equality by 3). It follows that if $`p`$ is even, so is $`p^{}`$ and we have
$$dimW^{}\frac{3p}{2}+3=\frac{3p+6}{2}=\frac{3(p+2)}{2}=\frac{3p^{}}{2}$$
and similarly if $`p`$ and $`p^{}`$ are odd we have
$$dimW^{}\frac{3p+1}{2}+3=\frac{3p^{}+1}{2}$$
so $`(๐ฎ_1,W^{})`$ satisfies property (i). Thus $`๐ฎ_1`$ violates the maximality of $`๐ฎ_0`$, so we conclude that the hypothesis of case 3 is impossible.
Case 4: Suppose $`pq2`$ and that there is a pair of triples $`T_l,T_l^{}`$ in $`๐ฎ๐ฎ_0`$ such that $`dimT_l,T_l^{}=5`$. Applying 3 we have four independent vectors
$$|\zeta _l,|\eta _lT_l,|\zeta _l^{},|\eta _l^{}T_l^{}$$
orthogonal to all column vectors of $`M`$ in columns outside of triples $`T_l,T_l^{}`$, so once again $`๐ฎ_1=๐ฎ_0T_lT_l^{}`$ with the subspace
$$W^{}=W|\zeta _l,|\eta _l,|\zeta _l^{},|\eta _l^{}$$
violates the maximality of $`๐ฎ_0`$. We conclude that the hypothesis of case 4 is impossible.
Case 5: The only remaining possibility is that $`pq3`$. Let $`๐ฏ=\{T_{j_1},T_{j_2},\mathrm{},T_{j_m}\}`$ be a set of triples in $`๐ฎ๐ฎ_0`$ with $`m3`$ minimal with respect to the property
$$dimT_{j_1},T_{j_2},\mathrm{},T_{j_m}<3m.$$
Applying 2 we have 2 vectors
$$|\zeta _k,|\eta _kT_k$$
for each of the $`m^{}2`$ elements $`kK`$. Let
$$๐ฎ_1=๐ฎ_0\left(\underset{kK}{}T_k\right),$$
let $`p^{}=p+m^{}`$, and let
$$W^{}=W\{|\zeta _k,|\eta _k\}_{kK}.$$
Note that property (ii) holds for $`(๐ฎ_1,W^{})`$. If $`m^{}<m`$, then the $`2m^{}`$ vectors in $`\{|\zeta _k,|\eta _k\}_{kK}`$ are independent by the minimality of $`๐ฏ`$, so we have
$$dimW^{}dimW+2m^{}\frac{3p}{2}+2m^{}=\frac{3p^{}+m^{}}{2}\frac{3p^{}+1}{2}$$
so property (i) holds for $`(๐ฎ_1,W^{})`$, but this contradicts the maximality of $`๐ฎ_0`$. Finally, if $`m^{}=m`$, then $`m4`$ (since $`m^{}`$ is even) and at least $`2(m1)`$ of the vectors in $`\{|\zeta _k,|\eta _k\}_{kK}`$ must be independent, again by the minimality of $`๐ฏ`$. If $`p`$ is even, then $`p^{}=p+m`$ is also even and we have
$$dimW^{}dimW+2(m1)\frac{3p}{2}+2(m1)=\frac{3p^{}+m4}{2}\frac{3p^{}}{2}.$$
If $`p`$ is odd, then $`p^{}=p+m`$ is odd and we have
$$dimW^{}dimW+2(m1)\frac{3p+1}{2}+2(m1)=\frac{3p^{}+m3}{2}\frac{3p^{}+1}{2}.$$
Thus $`๐ฎ_1`$ with the subspace $`W^{}`$ violates the maximality of $`๐ฎ_0`$. We conclude that the hypothesis of case 5 is impossible.
Having exhausted all possible cases, this completes the proof of 3. $`\mathrm{}`$
Next we state and prove a general statement about additivity of ranks for bipartite states.
Proposition 3.6. Let $`|\psi =|\psi _1|\psi _2`$ be a state vector for a bipartite state, where $`|\psi _j`$ is an $`n_j`$-qubit state vector for $`j=1,2`$. Let $`M`$ be the matrix associated to $`|\psi `$, and let $`M_j`$ denote the associated matrix for $`|\psi _j`$ for $`j=1,2`$. Let $`๐ฎ_1`$ and $`๐ฎ_2`$ be the following submatrices of $`M`$
$`๐ฎ_1`$ $`=`$ $`T_1T_2\mathrm{}T_{n_1}\{i|\psi \}`$
$`๐ฎ_2`$ $`=`$ $`T_{n_1+1}T_{n_2+2}\mathrm{}T_{n_1+n_2}\{i|\psi \}`$
and let $`^{}=T_{j_1}^{}T_{j_2}^{}\mathrm{}T_{j_m}^{}`$ be a union of triples in $`M_1`$ with corresponding union $`=T_{j_1}T_{j_2}\mathrm{}T_{j_m}`$ in $`M`$. We have
1. $`\text{rank}_{}M1=(\text{rank}_{}M_11)+(\text{rank}_{}M_21)`$,
2. $`\text{rank}_{}M_j=dim๐ฎ_j`$ for $`j=1,2`$, and
3. $`dim^{}=dim.`$
Proof. Let $`x`$ denote the state represented by $`|\psi `$ and let $`G`$ denote the local unitary group. Let $`x_j`$ denote the state represented by $`|\psi _j`$, and let $`G_j`$ denote the local unitary group for $`j=1,2`$, so we have $`G=G_1\times G_2`$.
It is easy to see that the $`G`$-orbit of $`x`$ is diffeomorphic to the product of the $`G_j`$-orbits of the $`x_j`$, so dimensions add.
$$dimGx=dimG_1x_1+dimG_2x_2$$
Applying 2, it follows immediately that (i) holds.
For (ii), observe that the columns of $`๐ฎ_j`$ are simply the comlumns of $`M_j`$ tensored with $`|\psi _2`$. The same reasoning applied to the subset $`๐ฎ_1`$ establishes (iii). $`\mathrm{}`$
We end this section with statements about factoring singlets and unentangled qubits.
Proposition 3.7. There are two triples $`T_l,T_l^{}`$ with $`dimT_l,T_l^{}=3`$ if and only if the state factors as a product of a singlet in qubits $`l,l^{}`$ times a state in the remaining qubits.
Proof. Without loss of generality, let us renumber the qubits so that $`l=1,l^{}=2`$.
First we prove the โifโ part of the statement. Let $`|\psi =|s|\varphi `$ where $`|s`$ is a singlet in qubits 1 and 2, and $`|\varphi `$ is an $`(n2)`$-qubit state. Then $`|\psi `$ is local unitary equivalent to $`|\psi ^{}=|s^{}|\varphi `$, where $`|s^{}=|00+|11`$. A simple calculation shows that the two triples in the matrix for a 2-qubit singlet state vector $`|s^{}`$ together span three dimensions. Therefore by 3 (iii), the dimension of the span of triples 1 and 2 in the matrix for $`|\psi ^{}`$ is also three. Since the rank of unions of triples is local unitary invariant by 3, we conclude that the dimension of the span of triples 1 and 2 in the matrix for $`|\psi `$ is also three.
Next we prove โonly if.โ Let $`|\psi `$ be a state vector for which $`dimT_1,T_2=3`$. By 3, $`|\psi `$ is local unitary equivalent to state vector $`|\psi ^{}`$ for which $`A_1|\psi ^{}=A_2|\psi ^{}`$ and $`C_1|\psi ^{}=C_2|\psi ^{}`$. Equation (4) and the hypothesis $`A_1|\psi =A_2|\psi `$ imply that $`(1)^{i_1}c_I=(1)^{i_2}c_I`$, so if $`c_I0`$ then $`i_1=i_2`$ mod 2, so every $`I`$ for which $`c_I0`$ has either both zeros or both ones in the first two indices. Equation (6) and the hypothesis $`C_1|\psi =C_2|\psi `$ imply that $`c_{I_1}=c_{I_2}`$ for all $`I`$. Apply this to $`I`$ for which $`i_1i_2`$ and we get
$$c_{(00i_3i_4\mathrm{}i_n)}=c_{(11i_3i_4\mathrm{}i_n)}$$
for all $`(i_3i_4\mathrm{}i_n)`$. It follows that $`|\psi ^{}`$ factors as a product
$$|\psi ^{}=(|00+|11)|\varphi $$
where $`|\varphi `$ is an $`(n2)`$-qubit state. Therefore $`|\psi `$ is a product of a singlet state in the first two qubits times a state in the remaining qubits. $`\mathrm{}`$
Lemma 3.8. If $`A_j|\psi =i|\psi `$ then the $`j`$th qubit is unentangled.
Proof. Renumber the qubits, if necessary, so that $`j=1`$. By (4), the hypothesis $`A_1|\psi =i|\psi `$ implies that if $`c_I0`$ then $`i_1=0`$. Therefore $`|\psi `$ factors as a product
$$|\psi =|0|\varphi $$
where $`|\varphi `$ is an $`(n1)`$-qubit state vector. $`\mathrm{}`$
Proposition 3.9. If the dimension of the span of a triple together with the rightmost column $`i|\psi `$ is three, then the state factors as an unentangled qubit times a state in the remaining qubits.
Proof. Let $`T_j`$ be a triple such that $`dimT_j\{i|\psi \}=3`$. Since $`i|\psi `$ lies in the span of $`T_j`$, we may write
$$i|\psi =\alpha A_j|\psi +\beta B_j|\psi +\gamma C_j|\psi $$
for some real $`\alpha ,\beta ,\gamma `$. Choose $`R\text{SO}(\text{su}(2))`$ such that
$$R(A)=\alpha A+\beta B+\gamma C.$$
Since the adjoint representation $`\text{Ad}:\text{SU}(2)\text{SO}(\text{su}(2))`$ is surjective, we can choose $`U_j\text{SU}(2)`$ such that $`\text{Ad}(U_j^{})=R`$, that is, $`U_j^{}XU_j=R(X)`$ for all $`X\text{su}(2)`$. For $`1kn`$, $`kj`$, set $`U_k`$ equal to the identity. Finally, let $`UG=\text{SU}(2)^n`$ be $`U=_{i=1}^nU_i`$. We have
$$U^{}A_jU|\psi =i|\psi .$$
Applying $`U`$ to both sides, we obtain
$$A_jU|\psi =iU|\psi .$$
Applying Lemma 3 to the matrix $`M^{}`$ for the state $`U|\psi `$ shows that the $`j`$th qubit is unentangled for the state $`U|\psi `$. Since unentanglement of a particular qubit is local unitary invariant, the proposition is established. $`\mathrm{}`$
Proposition 3.10. Rank of unentangled triples: If a state has $`k`$ unentangled qubits, the rank of the union of the triples corresponding to those qubits together with the rightmost column $`i|\psi `$ of $`M`$ is $`2k+1`$.
Proof. Let $`|\psi `$ be a state vector for a state with $`k`$ unentangled qubits, and let us renumber the qubits, if necessary, so that the unentangled qubits are numbered 1 through $`k`$. The state vector $`|\psi `$ is local unitary equivalent to a state vector
$$|\psi ^{}=|00\mathrm{}0|\varphi $$
where $`|00\mathrm{}0`$ is the product of $`k`$ unentangled qubits and $`|\varphi `$ is an $`(nk)`$-qubit state. Let $`M^{}`$ be the matrix associated to $`|\psi ^{}`$, $`M_1`$ the matrix for $`|00\mathrm{}0`$, and $`M^{\prime \prime }`$ the matrix for the single qubit state $`|0`$. Apply 3 (i) to $`|00\mathrm{}0=|0\mathrm{}|0`$ to get $`\text{rank}_{}M_1=2k+1`$ (using the fact that the single qubit state $`|0`$ has $`\text{rank}_{}M^{\prime \prime }=3`$). Then apply 3 (ii) to the set $`๐ฎ_1`$ which is the union of the first $`k`$ triples of the matrix $`M^{}`$ together with the column $`i|\psi ^{}`$ to get $`dim๐ฎ_1=\text{rank}_{}M_1=2k+1`$. Finally, apply 3 to conclude that the desired statement holds for $`|\psi `$. $`\mathrm{}`$
## 4 Minimum dimension orbit classification
Now we prove that any state with minimum orbit dimension is a product of singlets for $`n`$ even, times an unentangled qubit for $`n`$ odd.
Main Lemma 4.1. For $`n2`$, if $`M`$ has minimum rank, then there is some pair of triples whose span is three dimensional.
Proof. Suppose not. A consequence of 3 is that every pair of triples spans either three, five or six dimensions. We can rule out the possibility that some pair of triples spans five dimensions, as follows. If there is a pair $`T_l,T_l^{}`$ of triples which spans five dimensions, then by 3, the pair contributes four independent column vectors which are orthogonal to every column vector in the set $`๐ฎ=๐(T_lT_l^{})`$. Applying 3 to $`๐ฎ`$ with with $`q=n2`$, we have
$`\text{rank}M`$ $``$ $`4+\{\begin{array}{cc}\frac{3q}{2}+1& n\text{ even}\\ \frac{3q+1}{2}+1& n\text{ odd}\end{array}`$
$`=`$ $`\{\begin{array}{cc}\frac{3n+2}{2}+1& n\text{ even}\\ \frac{3n+3}{2}+1& n\text{ odd}\end{array}`$
which is greater than minimum, and therefore impossible.
Thus we need only consider the case where every pair of triples spans six dimensions.
Let $`๐ฏ=\{T_{j_1},T_{j_2},\mathrm{},T_{j_m}\}`$ be a set of $`m`$ triples minimal with respect to the property
$$dimT_{j_1},T_{j_2},\mathrm{},T_{j_m}<3m$$
(โminimalโ means that $`๐ฏ`$ contains no proper subset of triples which satisfy the given property; thus, any subset of $`m^{}<m`$ triples of $`๐ฏ`$ has โfullโ span of $`3m^{}`$ dimensions). We know such a $`๐ฏ`$ exists since the rank of $`M`$ is minimum. Apply 2 to get a subset $`K\{j_1,j_2,\mathrm{},j_m\}`$ with some even positive number $`m^{}`$ of elements and vectors $`|\zeta _k,|\eta _k`$ in $`T_k`$ for $`kK`$. Let $`๐ฏ^{}={\displaystyle \underset{kK}{}}T_k`$ and let $`๐ฎ=๐๐ฏ^{}`$ be the union of the $`q=nm^{}`$ triples not in $`๐ฏ^{}`$ together with the rightmost column $`i|\psi `$ of $`M`$.
If $`m^{}<m`$, the minimality of $`๐ฏ`$ guarantees that the $`2m^{}`$ vectors $`\{|\zeta _k,|\eta _k\}_{kK}`$ are independent, so the rank of $`M`$ is at least (apply 3 to $`๐ฎ`$ with $`q=nm^{}`$)
$`\text{rank}M`$ $``$ $`2m^{}+\{\begin{array}{cc}\frac{3q}{2}+1& n\text{ even}\\ \frac{3q+1}{2}+1& n\text{ odd}\end{array}`$
$`=`$ $`\{\begin{array}{cc}\frac{3n+m^{}}{2}+1& n\text{ even}\\ \frac{3n+1+m^{}}{2}+1& n\text{ odd}\end{array}`$
which is greater than minimum, so this case cannot occur.
If $`m^{}=m`$, the minimality of $`๐ฏ`$ guarantees that at least $`2(m^{}1)`$ of the vectors in $`\{|\zeta _k,|\eta _k\}_{kK}`$ are independent, so the rank of $`M`$ is at least
$`\text{rank}M`$ $``$ $`2(m^{}1)+\{\begin{array}{cc}\frac{3q}{2}+1& n\text{ even}\\ \frac{3q+1}{2}+1& n\text{ odd}\end{array}`$
$`=`$ $`\{\begin{array}{cc}\frac{3n+m^{}4}{2}+1& n\text{ even}\\ \frac{3n+1+m^{}4}{2}+1& n\text{ odd}\end{array}.`$
If $`m^{}>4`$, this is greater than minimum, so we may assume that $`m^{}=m4`$.
We rule out the possibility $`m=2`$ since every pair of triples spans six dimensions, so the only remaining case is $`m^{}=m=4`$. We may further assume that every minimal set of triples that has less than full span consists of $`m=4`$ triples, and that applying 2 to such a set yields $`m^{}=4`$. Thus the columns of $`M`$ decompose into a disjoint union
$$๐=๐ฎ_0๐ฏ_1๐ฏ_2\mathrm{}๐ฏ_t$$
where each $`๐ฏ_i`$ is a union of four triples for which applying 2 yields six independent vectors orthogonal to $`๐๐ฏ_i`$, and $`๐ฎ_0`$ is a union of $`q=n4t`$ triples together with the rightmost column $`i|\psi `$ of $`M`$ such that the span of the union of the triples in $`๐ฎ_0`$ is $`3q`$ dimensional.
We consider the cases $`q2`$, $`q=0`$, and finally $`q=1`$.
If $`q2`$ then we have
$$\text{rank}M6t+3q=6\left(\frac{nq}{4}\right)+3q=\frac{3n+3q}{2}$$
which is greater than minimum, so this case cannot occur.
If $`q=0`$ or $`q=1`$, let $`๐ฎ_1=๐ฎ_0๐ฏ_1`$ and consider the disjoint union
$$๐=๐ฎ_1๐ฏ_2\mathrm{}๐ฏ_t$$
of the set $`๐ฎ_1`$ with $`t^{}=t1`$ unions of four triples $`\{๐ฏ_i\}_{i=2}^t`$.
If $`q=0`$, $`๐ฎ_1`$ has a span of at least nine dimensions since any three of the triples in $`๐ฏ_1`$ have full span. We have
$$\text{rank}M9+6t^{}=9+6\left(\frac{n4}{4}\right)=\frac{3n+6}{2}$$
which is greater than minimum, so this case cannot occur.
If $`q=1`$, then $`๐ฎ_1`$ is the union of five triples together with the rightmost column $`i|\psi `$. If any subset of four of those five triples has full span, then we have
$$\text{rank}M12+6t^{}=12+6\left(\frac{n5}{4}\right)=\frac{3n+9}{2}$$
which is greater than minimum, so it must be the case that all subsets of four triples have less than full span, so any subset of four triples contributes six independent vectors orthogonal to the remaining triple and the rightmost column $`i|\psi `$ of $`M`$. If any one of the qubits corresponding to one of the five triples is not unentangled, then by 3 we have
$$\text{rank}M6+4+6t^{}=10+6\left(\frac{n5}{4}\right)=\frac{3n+5}{2}$$
which is greater than minimum, so it must be the case that all five qubits are unentangled. But then by 3 we have
$$\text{rank}M11+6t^{}=11+6\left(\frac{n5}{4}\right)=\frac{3n+7}{2}$$
which is greater than minimum.
Since all possible cases lead to contradictions, we conclude that some pair of triples must have a three dimensional span. $`\mathrm{}`$
Corollary 4.2. Any state which has minimum orbit dimension is a product of singlets when $`n`$ is even, together with an unentangled qubit when $`n`$ is odd.
Proof. Let $`|\psi `$ be a state vector for a state with minimum orbit dimension, with associated matrix $`M`$. Apply Lemma 4 to $`M`$ to get a pair of triples whose span is three dimensional. By 3, $`|\psi `$ factors as a product of a singlet in those two qubits times an $`(n2)`$-qubit state, say, with state vector $`|\psi _1`$, in the remaining qubits. Let $`M_1`$ be the matrix associated to $`|\psi _1`$. By 3 (i), $`M_1`$ also has minimum rank. Repeating this reasoning yields a sequence $`|\psi ,|\psi _1,|\psi _2,\mathrm{}`$ which eventually exhausts all the qubits of $`|\psi `$ unless $`n`$ is odd, in which case a single unentangled qubit remains. $`\mathrm{}`$
Next we classify states with minimum orbit dimension up to local unitary equivalence.
Proposition 4.3. Separation of singlet products: Products of singlets are local unitary equivalent if and only if the choices of entangled pairs are the same in each product.
Proof. If $`|\psi `$ is a product of singlets for which some pair of qubits, say qubits 1 and 2, forms a singlet, then $`|\psi `$ is of the form $`|\alpha |\beta `$, where $`|\alpha `$ is a singlet in qubits 1 and 2. Clearly, any state local unitary equivalent to $`|\psi `$ is also of this form. $`\mathrm{}`$
The following classification theorem summarizes the results of this section.
Theorem 4.4. Classification of states with minimum orbit dimension: An $`n`$-qubit pure state has minimum orbit dimension $`3n/2`$ ($`n`$ even) or $`(3n+1)/2`$ ($`n`$ odd) if and only if it is a product of singlets, together with an unentangled qubit for $`n`$ odd. Furthermore, two of these states which do not have the same choices for pairs of entangled qubits are not local unitary equivalent.
## 5 Conclusion
The classification of types of quantum entanglement is a difficult but central problem in the field of quantum information. Entanglement types partially distinguish themselves by their local unitary orbit and their orbit dimension. As an integer that can be readily calculated for a given quantum state, orbit dimension is a convenient LU invariant. It provides a useful โfirst stratificationโ of quantum state space, which suggests a two-step program for entanglement classification. The first step is to understand the possible orbit dimensions for a composite quantum system, and the second step is to understand the types of entanglement that occur in each orbit dimension. For pure states of $`n`$-qubits, the first step was achieved in . The present work completes the second step for the orbits of minimum dimension. In particular, states with minimum orbit dimension are precisely products of pairs of qubits, each pair in a singlet state (or LU equivalent). While there is much work left to be done in the second step of the classification program, it is worth remarking that the present results, dealing with an arbitrary number of qubits, give some hope that a meaningful classification of entanglement for $`n`$ qubits is possible.
Reduced orbit dimension states appear to be the most interesting states. Carteret and Sudbery pointed out that reduced orbit dimension states (which they call โexceptional statesโ) must have extreme values of the local unitary invariants that are used in the construction of all measures of entanglement. They concluded that reduced orbit dimension states should be expected to be particularly interesting and important. Subsequent studies have confirmed this expectation. All of the โfamousโ states that theorists use for examples and that experimentalists try to exhibit in the laboratory have reduced orbit dimension (examples include the EPR singlet state, the GHZ state, unentangled states, the W state, and $`n`$-cat states). Both the most entangled and the least entangled states states have reduced orbit dimension.
Orbit dimension provides a useful first step in entanglement classification, but the numerical value of orbit dimension, beyond providing a sense of how rare an entanglement type is, does not carry a simple physical meaning. For example, in the case of pure three-qubit states , the minimum orbit dimension is 5 and the maximum orbit dimension is 9. States with orbit dimension 5 consist only of products of a singlet pair and a qubit, orbit dimension 6 contains the unentangled states, orbit dimension 7 contains GHZ states as well as products of generic two-qubit states with an unentangled qubit, orbit dimension 8 contains the W state, and orbit dimension 9 contains all generic states.
Minimum orbit dimension states have the most symmetry with respect to local transformations in that they remain invariant under a larger class of transformations than any other states. They are maximal symmetry generalizations of the spin singlet state, which is invariant (as a quantum state, not an entanglement type) to any rotation applied identically to both spins. The present result, that the $`n`$-qubit maximal symmetry generalizations of the two-qubit singlet state are themselves products of singlets, shows a special role for the two-qubit singlet state in the theory of $`n`$-qubit quantum entanglement.
Linden, Popescu, and Wootters have shown that almost all pure $`n`$-qubit quantum states lack essential $`n`$-qubit quantum entanglement in the sense that they can be reconstructed from their reduced density matrices. It appears likely (and is known in the three-qubit case) that states with essential $`n`$-qubit entanglement are found among the reduced orbit dimension states. The present work shows that minimum orbit dimension states do not exhibit essential $`n`$-qubit entanglement for $`n3`$. Rather, minimum orbit dimension states maximize pairwise entanglement. We conjecture that states with minimum orbit dimension among *non-product* states have essential $`n`$-qubit entanglement in the sense of .
## Acknowledgments
We thank the anonymous referees for their helpful suggestions. Co-author Walck thanks the Research Corporation for their support.
## A Equation and statement numbers in
Table 1 gives a list of equation and statement numbers in the present paper with their matching numbers in . |
warning/0506/astro-ph0506545.html | ar5iv | text | # A Multi-wavelength study of the Pulsar PSR B1929+10 and its X-ray trail
## 1 INTRODUCTION
Currently, a consistent scenario for the evolution of the X-ray emission properties of aging rotation-powered pulsars does not exist yet. This surprising fact is largely due to the lack of sufficient observational data. Young<sup>1</sup><sup>1</sup>1In standards of high energy astronomy rotation-powered pulsars are called young, middle aged and old if their spin-down age is of the order of few times $`10^310^4`$ yrs, $`10^510^6`$ yrs and $`10^6`$ yrs, respectively. This classification is diffuse, though, with a smooth transition in between the different groups. and middle aged neutron stars, which emit strong pulsed non-thermal and/or surface hot-spot plus cooling emission, were studied reasonably well in the X-ray band. In contrast, until recently, most old radio pulsars were too faint for a detailed examination of their X-ray emission (cf. Sun et al. 1993; Manning & Willmore 1994; Becker & Trรผmper 1997; Saito 1998). However, especially old rotation-powered non-recycled pulsars are of particular interest for the study of particle acceleration and high energy radiation processes near the neutron starโs surface and in its magnetosphere. This is because their ages are intermediate between those of the well-studied young and cooling neutron stars, whose surface may produce copious thermal X-ray photons, and those of very old recycled millisecond pulsars, in which thermal hot-spot and non-thermal magnetospheric X-ray production mechanisms are believed to dominate. Old, non-recycled pulsars (outside globular clusters) therefore aid in answering questions such as how do the emission properties of the younger pulsars, like Geminga, PSR B0656+14 and PSR B1055-52, change as they age from $`10^5`$ to $`10^7`$ years? Will the thermal emission simply fade away due to cooling with increasing age or will the star be kept hot (at about $`0.51\times 10^5`$ K) over millions of years due to energy dissipation by processes such as internal frictional heating ($`\dot{E}_{diss}10^{28}10^{30}`$ erg/s) and crust cracking, as proposed by vortex creeping and pinning models? What happens to the non-thermal, hard-tail emission seen in the X-ray spectra of the middle-aged field pulsars? (See e.g. Saito 1988; Becker & Aschenbach 2002; De Luca et al. 2005). Will this emission become the dominant source or will this component also decay with time and will only thermal emission from the hot and heated polar-caps remain?
In order to address these questions, we initiated a program to study the X-ray emission properties of old rotation-powered pulsars with XMM-Newton, aiming to probe and identify the origin of their X-ray radiation. First results from parts of this project have been presented recently, reporting on the pulsars B0950+08, B0823+26, J2043+2740 (Becker et al. 2004; Zavlin & Pavlov 2004) and PSR B0628-28 (Becker et al. 2005) all of which have a spin-down age in the range from about one million to seventeen million years.
If one extrapolates the X-ray emission properties of young and cooling neutron stars to this age bracket, one may expect that the cooling emission fades away and thermal emission from heated polar caps dominates the X-rays. Surprisingly, the X-ray emission from old pulsars is largely dominated by non-thermal radiation processes. None of the pulsarsโ X-ray spectra required the addition of a thermal component, consisting of hot polar cap emission, to model their energy spectra. Further support for an emission scenario dominated by non-thermal mechanisms is given by the observed temporal emission properties. The pulse profiles of PSRs B0950+08 and B0628-28 are not broad and sinusoidal as would have been expected for spin-modulated thermal X-ray emission from heated polar caps, but are double peaked with narrow pulse components and pulsed fractions in the range of $`3050\%`$.
Some models, such as those by Harding & Muslimov (2001; 2002; 2003), predicted in the framework of their revised space-charge-limited flow model that polar cap heating, as a fraction of the spin-down luminosity, increases with pulsar age and should be most efficient for pulsars of spin-down age $`\tau `$$`10^7`$ yrs, if they are in fact producing pairs from curvature radiation photons. However, according to the Harding & Muslimov model, B0950+08 and B0823+26 cannot produce pairs from curvature radiation of primary electrons since they both lie below the curvature radiation pair death line in the $`P`$-$`\dot{P}`$ diagram of radio pulsars (cf. Becker et al. 2004 for a more detailed discussion).
A good candidate object to test these models is PSR B1929$`+`$10 which according to its X-ray emission properties can be considered to be prototypical of an old pulsar. With a pulse period of $`P=226.5`$ ms and a period derivative of $`\dot{P}=1.16\times 10^{15}`$, its characteristic age is determined to be $`3\times 10^6`$ years. These spin parameters imply a spin-down luminosity of $`\dot{E}=3.9\times 10^{33}\text{erg s}^1`$ and a magnetic field at the neutron star magnetic poles of $`B_{}5\times 10^{11}`$ G. With a radio dispersion measure of $`\mathrm{\hspace{0.17em}3.176}\text{pc cm}^3`$, the NE2001 Galactic free electron density model of Cordes & Lazio (2002) predicts a distance of 170 pc. However, the recent astrometric measurements by Chatterjee et al. (2004) yielded a precise proper motion and parallax determination that translates into an accurate distance measurement of $`d=361_8^{+10}`$ pc and a transverse speed of $`V_{}=177_5^{+4}`$ km s<sup>-1</sup>.
Thus, PSR B1929$`+`$10 is among the closest pulsars known. In addition to its relatively young age it appears to be the brightest among all old non-recycled X-ray detected rotation-powered pulsars. Its X-ray emission was discovered with the EINSTEIN observatory by Helfand (1983). Pulsed X-ray emission was discovered using a deep ROSAT observation (Yancopoulos et al. 1994). The pulse profile was found to be very broad with a single pulse stretching across the entire phase cycle, markedly different from the sharp peak observed in the radio band. The fraction of pulsed photons in the $`0.12.4`$ keV band was determined to be $`30\%`$. The ROSAT data were not able to constrain the nature of the pulsar emission as a blackbody spectrum (representing thermal polar-cap emission) and a power law model (representing non-thermal magnetospheric emission) fitted the data equally well (Becker & Trรผmper 1997). Similar results were obtained in the analysis of ASCA spectral/timing data by Wang and Halpern (1997) and Saito (1998). The broad sinusoidal pulse profile together with the higher column absorption of $`(0.61.1)\times 10^{21}\text{cm}^2`$ deduced from power law fits were taken as indirect arguments and strong indicator for a thermal polar-cap origin of the X-rays<sup>2</sup><sup>2</sup>2Fitting a blackbody model resulted in a column density of $`(13)\times 10^{20}\text{cm}^2`$ which was believed to be in better agreement with the close pulsar distance of $`170`$ pc as believed at that time..
In contrast, Slowikowska et al. (2005) recently found that a single blackbody spectral model cannot describe the pulsar spectrum if the ROSAT and ASCA observed spectra are modeled in a joint analysis. In their work it is shown that a single power law model or a composite model consisting of a two component blackbody spectrum can successfully describe the energy spectrum up to $`7`$ keV. The higher column density fitted by these models is found to be in agreement with that observed for other sources located near to the pulsarโs line of sight and at comparable distances (Slowikowska et al. 2005).
Regardless of the nature of its X-ray emission, PSR B1929$`+`$10 seems to be special for its extended X-ray emission which was discovered by Wang, Li & Begelman (1993) in their archival study of the Galactic soft X-ray background using deep ROSAT PSPC images. They found that the orientation of the diffuse X-ray emission is almost aligned with the pulsarโs proper motion direction, suggesting an interpretation in terms of an X-ray emitting trail behind the pulsar. If indeed associated with the pulsar, the trail could account for $`3\times 10^4`$ of the pulsarโs spin-down luminosity although the effective brightness may depend strongly on the density of the ambient interstellar matter. In recent years, the lack of confirmation of the trail from a subsequent $`350`$ ksec deep ROSAT HRI observation casts some doubt on its existence.
Near-UV emission from PSR B1929$`+`$10 has been detected in three partly overlapping spectral bands using the Hubble Space Telescopeโs Faint Object Camera (Pavlov et al. 1996) and the NUV-MAMA detector (Mignani et al. 2002). The nature of the optical emission is uncertain since the paucity of color information makes any spectral fit based on the optical data only merely tentative.
PSR B1929$`+`$10 has also been extensively observed at radio frequencies. The main peak of the radio profile, although much smaller than the X-ray pulse, has a substructure which can be modeled by six separate components (Kramer et al. 1994). Low-level emission connects this main pulse with an interpulse about 180 deg in longitude apart (see e.g. Everett & Weisberg 2001). The pulsarโs viewing geometry has been studied by many authors via the observed polarization angle swing, applying a rotating vector model. Everett & Weisberg (2001) reviewed the various results and concluded that both, the main and interpulse, are produced by nearly aligned rotation and magnetic axes and are emitted from nearly opposite sides of a wide, hollow cone. They derive an inclination of the magnetic axis with the spin axis of $`\alpha 36^o`$ while the impact angle of the line of sight was determined to be $`\beta 26^o`$ (cf. Fig.1).
In this paper we report on X-ray, optical H-alpha and radio observations of PSR B1929$`+`$10 which were made with XMM-Newton, the ESO New Technology Telescope (NTT) in La Silla (Chile) and the Jodrell Bank Radio Observatory in order to explore the spectral and timing emission properties of this interesting pulsar and its environment. The paper is organized in the following manner: in ยง2 we describe the radio, optical H-alpha and XMM-Newton observations of PSR B1929$`+`$10 and its X-ray trail and provide the details of the data processing and data filtering. The results of the spatial, spectral and timing analysis are given in ยง3. A summary and concluding discussion is presented in ยง4.
## 2 OBSERVATIONS AND DATA REDUCTION
### 2.1 RADIO OBSERVATIONS OF PSR B1929$`+`$10
The ephemerides for the analysis of the X-ray data were obtained from radio observations and the measurement of pulse timesโofโarrival (TOAs) using the 76-m Lovell radio telescope at Jodrell Bank Observatory. Table 1 summarizes the radio ephemerides of PSR B1929$`+`$10. A dual-channel cryogenic receiver system sensitive to two orthogonal polarizations was used predominantly at frequencies close to 1400 MHz. The signals of each polarization were mixed to an intermediate frequency, fed through a multichannel filter-bank and digitized. The data were de-dispersed in hardware and folded onโline according to the pulsarโs dispersion measure and topocentric period. The folded pulse profiles were stored for subsequent analysis. In a later offโline processing step, any sub-integrations corrupted by RFI were removed, the polarizations combined and the remaining sub-integrations averaged to produce a single totalโintensity profile for the observation. TOAs were subsequently determined by convolving, in the time domain, the averaged profile with a template corresponding to the observing frequency. The uncertainty on the TOA was found using the method, described by Downs & Reichley (1983) which incorporates the offโpulse RMS noise and the โsharpnessโ of the template. These TOAs were transferred to an arrival time at the solar system barycenter using the Jet Propulsion Laboratory DE200 solar system ephemeris (Standish 1982). More details can be found in Hobbs et al. (2004). Spectral data from PSR B1929$`+`$10 were obtained from the compilation of Maron et al. (2000).
### 2.2 OPTICAL OBSERVATIONS OF PSR B1929$`+`$10
We have performed optical H<sub>ฮฑ</sub> and R-band observations of the $`4.25\times 5.41`$ arcmin sky region around PSR B1929$`+`$10 using the SUperb-Seeing Imager (SUSI2) at the focus of the ESO New Technology Telescope<sup>3</sup><sup>3</sup>3See http://www.ls.eso.org/lasilla/sciops/ntt/index.html for a description of the ESO NTT and its instrumentation. in La Silla (Chile).
19 exposures of 600 s each ($``$3.2 hours total integration time) have been taken on July, 18 and 20 2004 through the H<sub>ฮฑ</sub> filter (central wavelength $`\lambda =6555.28\AA `$; $`\mathrm{\Delta }\lambda =69.76\AA `$). Additional 15 shorter exposures of $`3060`$ s in the R-band filter (central wavelength $`\lambda =6415.8\AA `$; $`\mathrm{\Delta }\lambda =1588.9\AA `$) were taken in order to discriminate H<sub>ฮฑ</sub> line emission from the continuum. In order to compensate for the $`8^{\prime \prime }`$ gap between the two SUSI2 CCD chips, the exposure sequences were taken with a jitter pattern with typical offset steps of $`20^{\prime \prime }`$ in RA. All observations have been performed with average airmass of 1.35, clear sky conditions and good seeing ($`0.6^{\prime \prime }`$). Single exposures were corrected for the instrumental effects (bias and dark subtraction, flat-fielding), cleaned of bad columns and of cosmic rays hits through median filter combination. To account for the uneven exposure map due to the jitter pattern, both the final H<sub>ฮฑ</sub> and R-band images were exposure-corrected.
The image astrometry was recomputed using as a reference the positions of a number of stars selected from the Guide Star Catalogue II (GSC-II), which has an intrinsic absolute astrometric accuracy of $`0\stackrel{}{\mathrm{.}}35`$ per coordinate<sup>4</sup><sup>4</sup>4http://www-gsss.stsci.edu/gsc/gsc2/GSC2home.htm. The pixel coordinates of the reference stars have been computed by a two-dimensional Gaussian fitting procedure, and transformation from pixel to sky coordinates was then computed using the programme ASTROM<sup>5</sup><sup>5</sup>5http://star-www.rl.ac.uk/Software/software.htm, yielding an rms of $``$ 0$`\stackrel{}{\mathrm{.}}`$08 in both Right Ascension and Declination, which we assume representative of the accuracy of our astrometric solution.
### 2.3 XMM-NEWTON OBSERVATIONS OF PSR B1929$`+`$10
PSR B1929$`+`$10 was observed by XMM-Newton<sup>6</sup><sup>6</sup>6For a description of XMM-Newton, its instrumentation and the various detector modes available for observations see http://xmm.vilspa.esa.es/. as part of the European Photon Imaging Camera (EPIC) guaranteed time program. Observations were performed on 2003 November 10 (XMM rev. 718) with a duration of $`11`$ ksec and five months later on 2004 April 27 (XMM rev. 803) and April 29 (XMM rev. 804) for a duration of $`22`$ ksec and $`23`$ ksec, respectively. In all three observations the EPIC Positive-Negative charge depleted Silicon Semiconductor (PN camera) was used as the prime instrument. The two Metal Oxide Semiconductor cameras (MOS1 & MOS2) were operated in PrimeFullWindow mode to obtain imaging and spectral data. The EPIC-PN camera was set up to operate in SmallWindow read-out mode which provides imaging, spectral and timing information with a temporal resolution of 5.67 ms which is more than sufficient to resolve the 226 ms period of PSR B1929$`+`$10. The SmallWindow mode was preferred over other EPIC-PN imaging modes because of its higher temporal resolution, albeit its $`30\%`$ higher dead-time caused a decrease in the net exposure by $`1/3`$. The medium filter was used for the EPIC-PN and MOS1/2 cameras in all observations to block optical stray light. Given the target flux, both the RGS and optical monitor are of limited use. A summary of exposure times, instrument modes and filters used for the X-ray observations of PSR B1929$`+`$10 is given in Table 2.
XMM-Newton data have been known to show timing discontinuities in the photon arrival times with positive and negative jumps of the order of one to several seconds (Becker & Aschenbach 2002; Kirsch et al. 2004). While an inspection of the data processing log-files reveals that none of the EPIC-PN data taken from PSR B1929$`+`$10 exhibit such discontinuities, we nevertheless used a release track Version of the XMM-Newton Standard Analysis Software (SAS) version 6.5 (released in August 2005) to process the EPIC-PN data. This software detects and corrects most of the timing discontinuities during data processing. In addition, known timing offsets due to ground station and spacecraft clock propagation delays are corrected by this software in using newly reconstructed time correlation (TCX) data (cf. Becker et al. 2006). Barycenter correction of the EPIC-PN data and all other analysis steps were performed by using SAS Version 6.1. Data screening for times of high sky background was done by inspecting the light-curves of the EPIC-MOS1/2 and PN data at energies above 10 keV. Apart from having a rather high sky background contribution in both April observations, very strong X-ray emission from soft proton flares is covering about half of these April data sets. The data quality of the shorter November 2003 observation is not reduced by these effects.
Using light-curves with 100 s bins, we rejected time intervals where the MOS1/2 had more than 130, 140 and 175 cts/bin in the 2003 November, 2004 April 27 and 29 observations, respectively. For the EPIC-PN data sets, we rejected times with more than 7, 12 and 40 cts/bin. The data screening reduced the effective exposure time for the MOS1/2 and PN-camera to a total of 64 ksec and 23.8 ksec, respectively.
For the spectral analysis based on the MOS1/2 data we used only those events with a detection pattern between $`012`$ (i.e. single, double and triple events) and the flag parameter set to less than, or equal to, 1. The latter criterion excludes events which are located near a hot pixel, or a bright CCD column, or which are near the edge of the CCD. For the EPIC-PN timing and spectral analysis, we used only single and double events, i.e. those which have a pattern parameter of less than, or equal to, 4 and a flag parameter equal to zero. The energy range of the MOS1/2 and EPIC-PN CCDs was restricted to $`0.310`$ keV for the spectral analysis due to calibration issues towards softer spectral channels and to $`0.210`$ keV for the timing analysis.
## 3 ANALYSIS OF THE MULTI-WAVELENGTH DATA OF PSR B1929$`+`$10
The X-ray counterpart of PSR B1929$`+`$10 is detected with high significance in both the MOS1/2 and EPIC-PN data. The count rates are $`0.0283\pm 0.0001`$ cts/s (MOS1/2) and $`0.078\pm 0.003`$ cts/s (EPIC-PN) within the $`0.210`$ keV band. Inspection of the MOS1/2 and EPIC-PN images revealed diffuse extended emission at the position of the putative X-ray trail seen in the ROSAT data by Wang, Li and Begelman (1993). Figure 2 shows the MOS1/2 image made with the 2003 November and 2004 April data. Contour lines indicate the diffuse X-ray trail lying in the direction opposite to the transverse motion of the pulsar ($`\mu _{\mathrm{RA}}=17.00\pm 0.27`$ mas yr<sup>-1</sup> and $`\mu _{\mathrm{DEC}}=9.48\pm 0.37`$ mas yr<sup>-1</sup>, Chaterjee et al. 2004). Contour lines, obtained from a re-analysis of 45 ksec ROSAT PSPC and 350 ksec ROSAT HRI observations, are overlaid. Owing to a sensitivity ratio between the PSPC and HRI of up to a factor of $`5`$ (depending on the source spectrum) and a much higher detector noise the faint trail like emission is much harder to detect in the HRI data. Only after applying an adaptive kernel smoothing procedure a source structure which matches with the shape of the trail seen in the PSPC and MOS1/2 images becomes visible in this data. An interesting difference, though, is that in the HRI data the trail seems to break up into two separate pieces. In the MOS1/2 images this region is partly covered by a CCD gap. The angular resolution of the ROSAT detectors are 25 arcsec in the PSPC and 5 arcsec in the HRI, respectively, while that of XMM-Newton is 15 arcsec (Half Energy Width).
Inspecting the MOS1/2 full field of view (cf. Figure 3a) there is some indication for trail emission beyond the edge of the inner MOS CCD. To reduce the impact on the trail detection due to the presence of other X-ray point sources in the field (though no bright sources are located along the trail) we have removed point source contributions from the MOS1/2 data and corrected the image for telescope vignetting effects. The result is displayed in Figure 3b. Clearly, there is a whiff of emission which probably extends down to the edge of the detectorโs field of view, but its significance at locations more distant from the pulsar is gradually fading into the background. In order to estimate the significance of the emission along the tail we extracted four circles at increased distance from the pulsar, at 1.5โ, 4.5โ, 7.5โ, and 10.5โ, respectively. The significance of the trail emission in these circles is computed to be (C1): $`493\text{scts}/\sqrt{1417\text{bcts}}=13\sigma `$, (C2): $`183\text{scts}/\sqrt{1197\text{bcts}}=5.3\sigma `$, (C3): $`161\text{scts}/\sqrt{1234\text{bcts}}=4.6\sigma `$, and (C4): $`125\text{scts}/\sqrt{1100\text{bcts}}=3.8\sigma `$, where scts and bcts denote the source counts and background counts measured for these four circles in the energy band $`0.210`$ keV.
In the first half of the 80s the Effelsberg 100-m radio telescope has mapped the entire Galactic Plane in the latitude range $`\pm 5^{}`$ at 11 cm wavelength (Reich et al. 1990). Inspecting this data for a possible radio counterpart of the pulsarโs X-ray trail we discovered an elongated feature which roughly matches with that observed in the X-ray band. Figure 4a shows the $`30\times 30`$ arcmin region centered on PSR B1929$`+`$10 as observed in the 11 cm Effelsberg survey. The extracted map has 140 x 140 pixels. Its resolution (HPBW) is 4.3โ with fluxes ranging from 0 to 250 mKTb (1 mKTb = 0.398 mJy). The rms of the survey was given as 20 mKTb (= 8 mJy). In the vicinity of the source we found a diffuse background of 25 mJy. The pulsar itself was evident with 25 mJy above the local background level. The radio tail was visible with fluxes of up to 15 mJy above local background and unresolved in the transverse direction. The tail was not aligned with the low level residual scanning effects. The distance to the bright node on the tail was found to be 8.8โ, but more aligned features were visible further in the beam out to 12โ. In order to estimate the probability of a spurious detection we derived a probability of a pixel exceeding 34 mJy as $`p_x(34\mathrm{m}\mathrm{J}\mathrm{y})=0.374`$ from a histogram of all pixels in our map. We made another histogram from a 33โ (78 pixel) line along the tail direction from the pulsar, excluding the first 6โ of the point source. 18 of the remaining 64 pixels were found to be above 34 mJ. The binominal distribution gives a probability of $`p_{rand}=0.032`$ for a random occurrence of such a result along an arbitrary line of similar length in our map. Using the resolution (4.3โ) for a conservative estimate of the tail width, we find that the probability of a random alignment in the direction of proper motion is about $`p_{align}=0.078`$. Hence the probability of a coincidence of weak random radio sources with the observed psr tail is about $`p_{radio}=p_{align}p_x=2.510^3`$
Nevertheless, inspecting the NRAO VLA Sky Survey (NVSS) which is a 1.4 GHz continuum survey mapping the entire sky north of -40 deg declination (Condon et al. 1998) we found faint not fully resolved radio sources which match in position with the brighter end of the radio tail seen at 11cm (2.72 GHz). There is no radio emission seen bridging the pulsar and these faint sources, however, this absence is likely a function of the VLA telescope configuration which affects the sensitivity of the observations to diffuse emission. The NVSS image has an angular resolution of 45 arcsec (FWHM) and is shown in Figure 4b. An elongated emission feature seen at RA(2000) 19<sup>h</sup> 31<sup>m</sup> 36<sup>s</sup>, DEC +10<sup>d</sup> 51โ 54โ, and thus still located within the segment indicated in Figure 3b, roughly lines up with the pulsarโs proper motion direction. Making use of such an alignment to claim a relation with PSR B1929$`+`$10, however, seems premature. As there are other bright sources in the 11cm Effelsberg data which are not seen in the NVSS image (e.g. the bright circularly shaped source located in the north-east of Figure 4a) the NVSS data cannot help to constrain the nature of the elongated radio feature seen at 11cm. A comparison of the Effelsberg and NVSS radio images with the ROSAT PSPC and HRI X-ray images is shown in Figure 5. Interestingly, the part of the tail which in the HRI appears to be separated from the pulsar matches with the position of the three not fully resolved radio sources in the NVSS, suggesting that both could be related.
The NTT H<sub>ฮฑ</sub> image which shows the $`4.25\times 5.41`$ arcmin sky region around PSR B1929$`+`$10 is shown in Figure 6. Contour lines from the MOS1/2 image are overlaid. The inspection of the H<sub>ฮฑ</sub> image and of the ratio of the H<sub>ฮฑ</sub> to $`R`$-band images does not reveal any diffuse structure that could be related to the brighter parts of the pulsarโs X-ray trail. We performed a rough flux calibration of the H<sub>ฮฑ</sub> image by computing the H<sub>ฮฑ</sub> fluxes corresponding to the R magnitudes of a sample of GSC-II stars which are included in the field assuming blackbody spectra with the appropriate stellar temperatures. We then derived the relation between the instrumental magnitude and the H<sub>ฮฑ</sub> fluxes. A flux upper limit of $`10^{16}\text{erg s}^1\text{cm}^2\text{arcsec}^2`$ in the H<sub>ฮฑ</sub> band can be set for the sky region around the pulsar.
As can be seen in Figure 6, there is a bright star near to the position of PSR B1929$`+`$10. A second bright star seems to coincide with the brighter part of the trail near to the pulsar. In view of the density of stars in the observed field the chance probability to have a star near to the position of an X-ray source is quite high. We investigated therefore whether these stars contribute to the flux recorded from PSR B1929$`+`$10 and the brighter part of its trail. The bright star close to PSR B1929$`+`$10 has been identified to be of K(4-6)III-I type (Kouwenhoven & van der Berg 2001) and has a magnitude in B and R of $`14.85`$ and $`13.44`$, respectively. The bright star within the trail is possibly of K2-4 class and according to the GSC-II and the 2MASS catalogues has B, R, J and K magnitudes of $`13.97,12.13,10.6`$ and $`10.0`$, respectively. With $`\mathrm{log}(F_x/F_{opt})=2.77\pm 1.0`$ (Krautter et al. 1999), we find for a mekal plasma model with kT=0.35, solar abundances, and a column absorption of $`N_H10^{21}\text{cm}^2`$ a possible contribution from these stars to the soft channels in XMM-Newton of $`1020\%`$. Given the very large spread in the emission properties of K stars and taking its colors into account the star near to PSR B1929$`+`$10 is possibly a bright giant, and thus could lie towards the lower end of the $`F_x/F_{opt}`$ range (e.g. Zickgraf et al. 2005) which then would imply a negligible flux contribution in the pulsar extraction region. As far as the pulsarโs X-ray trail is concerned it is clear from its length and from its hard X-ray spectrum (cf. ยง3.2.3) that this trail is not due to unresolved foreground or background sources.
### 3.1 TIMING ANALYSIS
We used all the EPIC-PN SmallWindow mode data for the timing analysis, including those times of high sky background which were excluded for spatial and spectral analysis. Experience has shown that this does not affect the results from the timing analysis if the sky background is properly taken into account in all pulsed-fraction estimates.
Events were selected from a circle of 20 arcsec radius centered on the pulsar radio timing position (cf. Table 1). For the barycenter correction we applied the standard procedures for XMM-Newton data using barycen-1.17.3 and the JPL DE200 Earth ephemeris (Standish 1982) to convert photon arrival times from the spacecraft to the solar system barycenter (SSB) and the barycentric dynamical time (TDB). The pulsar radio timing position (cf. Table 1) was used for the barycenter correction. The spin-parameters $`f`$ and $`\dot{f}`$ of PSR B1929$`+`$10 are known with high precision from our contemporaneous radio observations, covering all XMM-Newton orbits relevant for our analysis. PSR B1929$`+`$10 is not known to show timing irregularities (glitches) so that we can fold the photon arrival times using the pulsarโs radio ephemeris. The statistical significance for the presence of a periodic signal was obtained from a $`Z_n^2`$-test with $`110`$ harmonics in combination with the H-Test to determine the optimal number of harmonics (De Jager 1987; Buccheri & De Jager 1989). The optimal number of phase bins for the representation of the pulse profile was determined by taking into account the signalโs Fourier power and the optimal number of harmonics deduced from the H-Test (see Becker & Trรผmper 1999 and references therein).
Within the $`0.210`$ keV energy band, 5736 events were available for the timing analysis of which $`51\%`$ are estimated to be instrument and sky background. The $`Z_n^2`$-test gave 82.14 for $`n=6`$ harmonics ($`Z_1^2=48.61`$). According to the H-Test, the probability of measuring $`Z_6^2=82.14`$ by chance is $`1.6\times 10^{12}`$. The significance of the pulsed signal thus is comparable with that found recently in the other old pulsars B0950+08 (Becker et al. 2004) and PSR B0628-28 (Becker et al. 2005).
The $`0.210`$ keV pulse profile is shown together with a radio profile observed at 1.4 GHz in Figure 7. It reveals a significant deviation from a sinusoidal pulse shape. The X-ray pulse profile is found to consist of at least two pulse peaks, a broader component of $`273^{}`$ width, plus a narrow component of $`44^{}`$ width. Taking the center of mass of the pulse as a reference point, the narrow X-ray pulse appears to be slightly phase shifted (by $`22^{}`$) from the location of the main radio peak, although both components overlap well in phase. The broader pulse appears to have substructures hinting the presence of two narrower pulse peaks which are not fully resolved by the available data. Indeed, modeling the profile with three Gaussians yields acceptable results, supporting such an interpretation. The fitted functions together with the post-fit residuals are shown in Figure 8. The centers and widths (FWHM) of each of these three components obtained from the fit are, in pulse longitudes, $`(112^{},80^{})`$, $`(220^{},75^{})`$, and $`(346^{},22^{})`$, respectively. The uncertainty in the center position of each component is estimated as $`8^{}`$ $`(1\sigma `$).
In order to search for any energy dependence in the pulsarโs temporal X-ray emission properties, we restricted the timing analysis to the $`0.21.0`$ keV, $`1.02.1`$ keV, and $`2.110`$ keV energy bands. This analysis shows that the pulsed signal appears to be most significant if we consider only events which are recorded at energies below $`2.1`$ keV. The maximum $`Z_n^2`$-values found for the pulsed signal in the three energy bands are $`Z_1^2=32.12`$, $`Z_3^2=52.1`$, and $`Z_6^2=28.4`$, respectively. The fact that the pulsed signal is less significant beyond $`2.1`$ keV can be explained by an increase of instrumental background along with a decrease of pulsar signal, yielding a lower signal-to-noise ratio compared to the $`0.21.0`$ keV and $`1.02.1`$ keV energy bands. An inspection of the pulse profiles for these two energy bands reveals that the emission from the narrow pulse peak appears only in the profile beyond $`1`$ keV. Computing the fraction of pulsed photons (pulsed fractions, hereafter), we find $`32\pm 4\%`$ for the total $`0.210`$ keV energy band. The pulsed fractions in the sub-bands $`0.21.0`$ keV, $`1.02.1`$ keV, and $`2.110`$ keV are $`24\pm 5\%`$, $`44\pm 6\%`$, and $`17\pm 17\%`$, respectively (errors represent the $`1\sigma `$ confidence limits). For the narrow peak which appears at energies $`1`$ keV, the pulsed fraction in the $`1.02.1`$ keV band is slightly higher than within the $`0.21.0`$ keV band, although the significance for this is at the $`2\sigma `$ level only.
Computing the TOAs of the X-ray pulse we note that uncertainties of the XMM-Newton clock against UTCs are not relevant as those are on a scale of $`100\mu s`$ (Becker et al. 2005 in prep.) and thus are a factor of $`140`$ smaller than the bin width of the X-ray pulse profile shown in Figure 7. The largest uncertainty from comparing radio with X-ray pulse arrival is the definition of a suitable reference point which we choose to be the center of mass of a pulse peak.
### 3.2 SPECTRAL ANALYSIS
#### 3.2.1 PSR B1929$`+`$10
The energy spectrum of PSR B1929$`+`$10 was extracted from the MOS1/2 data by selecting all events detected in a circle of 40 arcsec radius, centered on the pulsar position. According to the XMM-Newton/EPIC-MOS model point spread function, $`87\%`$ of all events of a point source are within this region. The background spectrum was extracted from a source-free circular region of 30 arcsec radius, northwest from the pulsar at RA(2000) $`19^h\mathrm{\hspace{0.17em}32}^m\mathrm{\hspace{0.17em}08.59}^s`$, DEC $`10^{}\mathrm{\hspace{0.17em}59}^{}\mathrm{\hspace{0.17em}36.49}\mathrm{}`$. A second background spectrum was extracted from a circle of the same size but centered east from the pulsar at the position RA(2000) $`19^h\mathrm{\hspace{0.17em}32}^m\mathrm{\hspace{0.17em}19.5}^s`$, DEC $`10^{}\mathrm{\hspace{0.17em}59}^{}\mathrm{\hspace{0.17em}56.1}\mathrm{}`$. This second background spectrum was used to check the independence of the spectral results with respect to the selected background region.
For the EPIC-PN data we used an extraction radius of 30 arcsec centered on PSR B1929$`+`$10. This selection region includes $`85\%`$ of the point source flux. Out-of-time events and a gradient of decreasing background towards the PN-CCD readout node requires that the background spectrum be selected from regions located at about the same CCD row level as the location of the pulsar. We therefore extracted two background spectra from source-free circular regions of radius 30 arcsec about 1.5 arc-minute east and west of the pulsar position. The spectral results were found to be independent of the specific background region used.
Because of its unique soft response we made use of archival ROSAT PSPC data and extracted the pulsar spectrum from a 60 arcsec circular region in the PSPC. The background spectrum was extracted from an annulus of 70 arcsec and 100 arcsec inner and outer radius, respectively, centered on the pulsar.
In total, the extracted spectra include 2405 EPIC-PN source counts and 1536 EPIC-MOS1/2 source counts. The PN and MOS1/2 spectral data were dynamically binned so as to have at least 30 counts per bin. 462 additional source counts, recorded within $`0.12.4`$ keV by ROSAT, were available for the joined ROSAT plus XMM-Newton spectral analysis of PSR B1929$`+`$10. As the energy resolution of the ROSAT PSPC was only $`30\%`$ at 1 keV (Briel et al. 1989) we binned the ROSAT data so as to have 4 independent spectral bins. Owing to its unique soft response, ROSAT data contribute information primarily to the lowest spectral channels.
Model spectra were then simultaneously fitted to the ROSAT and XMM-Newton pulsar data. An anomaly in the MOS2 spectral data below $`0.8`$ keV, which probably is related to the MOS redistribution problem in which events from higher energies incorrectly are redistributed downwards (cf. Sembay et al. 2004), required exclusion of those MOS2 spectral bins from the analysis.
Amongst the single component spectral models, a power law model was found to give the statistically best representation ($`\chi ^2`$=119.6 for 121 dof) of the observed energy spectrum. A single blackbody model ($`\chi ^2`$=246.9 for 121 dof) which was used by Yancopoulos et al. (1994) and Wang & Halpern (1997) to describe the ROSAT and ASCA observed pulsar spectrum did not give acceptable fits and is finally rejected. The best-fit power law spectrum and residuals are shown in Figure 10. Contour plots showing the relationship between the photon index and the column absorption for various confidence levels are shown in Figure 11.
In the following, we describe various models fitted to the energy spectrum of PSR B1929$`+`$10. The resulting spectral parameters are summarized in Tables 3 and 4 where all errors represent the $`1\sigma `$ confidence range computed for one parameter of interest. The power law model yields spectral parameters which are in good agreement with the results obtained by Slowikowska et al. (2005) based on their joint analysis of ROSAT plus ASCA data. The unabsorbed energy fluxes and luminosities (see Table 4) imply a rotational energy to X-ray energy conversion efficiency $`L_x/\dot{E}`$ of $`\mathrm{\hspace{0.17em}1.1}\times 10^3`$ within $`0.510`$ keV and of $`3.4\times 10^3`$ within the ROSAT band (cf. Becker & Trรผmper 1997).
In addition to the single component spectral models we tested two composite models consisting of either a blackbody plus power law component or of two blackbody components, respectively. The first model represents the scenario in which the X-ray emission of PSR B1929$`+`$10 originates from heated polar caps and from magnetospheric radiation processes. The double blackbody model implies that its X-ray emission would be entirely of thermal origin, e.g. with all the X-rays being emitted from heated polar caps with an anisotropic temperature distribution for which the model is approximated by two Planckian spectra with different emitting areas and temperatures.
Clearly, the statistical motivation to include an additional thermal component to the already excellent power law fit is very small. All combinations of blackbody normalizations and temperatures that were fitted gave reduced $`\chi ^2`$-values which are not better than the fits to a power law model of the pulsar spectrum. The F-test statistic for adding the extra blackbody spectral component to the power law model, thus, is very low. Based upon the errors of the fitted spectral components all fits of this composite model yield only upper limits for a maximum thermal component which the power law model โacceptsโ before the fits become statistically unacceptable. The parameters obtained for the thermal component are thus intrinsically upper limits. The fitted parameters are shown in Tables 3 and 4. The resulting blackbody temperature and the size of the projected emitting area are $`kT0.28`$ keV and $`R_{bb}110`$ m, assuming a pulsar distance of 361 pc.
In computing the relative contributions of the thermal and non-thermal spectral components and stretching the errors to the limits for the composite blackbody plus power law model, we find that no more than $`40\%`$ of the detected X-ray flux could come from heated polar caps. For the best fit parameters only $`7\%`$ could be emitted from these caps. However, the radius of $`110`$ m which we computed for the maximum projected emitting area is not too different from the expected size of a polar cap. For comparison, defining the size of the presumed polar cap as the foot points of the neutron starโs dipolar magnetic field, the radius of the polar cap area is given by $`\rho =\sqrt{2\pi R^3/cP}`$ with $`R`$ being the neutron star radius, c the velocity of light and P the pulsar rotation period (see e.g. Michel 1991). For PSR B1929$`+`$10 with a rotation period of 226 ms this yields a polar cap radius of $`\rho 300\text{m}`$. A plot of the thermal and non-thermal spectral components and the combined model is shown in Figure 12.
With a spin-down age of $`3\times 10^6`$ years PSR B1929$`+`$10 should still have some residual heat content from its birth event. Depending on the equation of state the surface temperature could still be in the range $`13\times 10^5`$ K (cf. Becker & Pavlov 2001 and references therein) and could contribute on a low level to the detected soft X-ray emission. To estimate the upper limit for the surface temperature of PSR B1929$`+`$10 we have fixed the blackbody normalization so that the emitting area corresponds to the surface of a $`R=10`$ km neutron star and calculated the confidence ranges of the blackbody temperature by leaving all other parameters free. We find a $`3\sigma `$ surface temperature upper limit of $`T_s^{\mathrm{}}<4.5\times 10^5`$ K which is somewhat above the temperatures predicted by cooling models (e.g. Page & Applegate 1992; Tsuruta 1998; Yakovlev et al. 1999) and, thus, may constrain only those thermal evolution models which predict extreme reheating.
The spectral model consisting of two thermal components describes the observed spectrum formally with comparable goodness of fit as the blackbody plus power law model. In comparison with a single blackbody fit the F-test statistic thus supports an addition of a second thermal component to the single blackbody model. However, inspection of the fit residuals shows that the double blackbody model falls off rapidly beyond $`5`$ keV which causes the residuals beyond that energy to systematically lie above the zero line, albeit error bars are large. The fitted model parameters are a column absorption of $`N_H=0.35_{0.08}^{+0.11}\times 10^{21}\text{cm}^2`$, temperatures $`kT_1=0.59_{0.05}^{+0.06}`$, $`kT_2=0.20_{0.2}^{+0.2}`$, and projected emitting areas $`R_1=9.8_{1.8}^{+2.1}\text{m}`$ and $`R_2=92_{11}^{+14}\text{m}`$, respectively.
#### 3.2.2 PHASE RESOLVED SPECTRAL ANALYSIS
In order to investigate a possible variation of the pulsar emission spectrum as a function of pulse phase $`\varphi `$ we selected the events from phase intervals $`\varphi _1=]\mathrm{\hspace{0.17em}0.15}0.65[`$ and from $`\varphi _2=[\mathrm{\hspace{0.17em}0.0}0.15;0.651.0]`$ (cf. Fig.7) for a spectral analysis. The source counts are from a circle of $`20^{\prime \prime }`$ radius centered on the target coordinates (the extraction region was smaller than the one used for the phase-integrated case in order to improve the signal-to-noise in the lower-statistics phase-resolved spectra). The spectra were binned in order to have at least 30 counts per channel. Allowing both the $`N_\mathrm{H}`$ and the photon index to vary we obtained for the phase $`\varphi _1`$ $`N_\mathrm{H}=(0.22\pm 0.03)\times 10^{22}`$ cm<sup>-2</sup>, $`\mathrm{\Gamma }=3.1\pm 0.2`$, ($`\chi _\nu ^2=1.2`$; 44 dof) and for $`\varphi _2`$ $`N_\mathrm{H}=(0.17\pm 0.03)\times 10^{22}`$ cm<sup>-2</sup>, $`\mathrm{\Gamma }=2.7\pm 0.2`$ ($`\chi _\nu ^2=1.03`$, 36 dof). Fixing the interstellar column to the phase-averaged value, the best fit values for the photon indices are $`\mathrm{\Gamma }=3.05\pm 0.07`$ ($`\chi _\nu ^2=1.18`$, 45 dof) and $`\mathrm{\Gamma }=3.00\pm 0.09`$ ($`\chi _\nu ^2=1.05`$, 37 dof) for the peak and off-pulse, respectively. Using composite (blackbody + power law) models for both $`\varphi _1`$ and $`\varphi _2`$ did not yield better results. Within the statistical uncertainties there are therefore no spectral changes observed as a function of pulse phase.
#### 3.2.3 Diffuse emission from the X-ray trail of PSR B1929$`+`$10
The large collecting area of XMM-Newton allows, for the first time, a spectral analysis of the emission from the pulsarโs X-ray trail. The satellite roll-angle and the EPIC-PN reduced field of view of $`4.4\times 4.4`$ arcmin in SmallWindow mode caused only a small portion of the trail to be observed by the PN CCD #4 in the short November 2003 observation. During the April 2004 observations the trail is outside the PNs field of view. Due to the increased instrument background towards the PNโs readout node which is near to the location of the trail region in the CCD #4 we did not include the PN data in the analysis of the diffuse emission. Both MOS cameras have covered a large portion of the trail in all observations (cf. Figure 3) so the analysis was restricted to the MOS1/2 data only.
The photon statistic of the diffuse trail emission is sufficient for a detailed spectral modeling only near to the pulsar. We, therefore, extracted the X-ray spectrum from a circular region of 1 arcmin radius centered at RA(2000) $`19^h\mathrm{\hspace{0.17em}32}^m\mathrm{\hspace{0.17em}07.5}^s`$, DEC $`10^{}\mathrm{\hspace{0.17em}58}^{}\mathrm{\hspace{0.17em}32}\mathrm{}`$. The background spectrum was extracted from a source-free region located at RA(2000) $`19^h\mathrm{\hspace{0.17em}32}^m\mathrm{\hspace{0.17em}01.6}^s`$, DEC $`10^{}\mathrm{\hspace{0.17em}59}^{}\mathrm{\hspace{0.17em}36}\mathrm{}`$. About 800 counts ($`60\%`$ background contribution) were available for the spectral analysis of the diffuse trail emission. We binned the spectrum dynamically so as to have at least 30 counts per bin.
The pulsarโs X-ray trail is likely formed by a ram-pressure confined pulsar wind. Its X-ray emission should arise from synchrotron radiation of relativistic electrons with a spectral shape characterized by a power law. To test this hypothesis we fitted a power law to the extracted trail spectrum and found the model describes the observed spectrum well ($`\chi ^2`$=20.4 for 23 dof). The best-fit power law spectrum and residuals are shown in Figure 13. Contour plots showing the relationship between the photon index and the column absorption for various confidence levels are shown in Figure 14. Details of the spectral fits are again listed in Tables 3 and 4. Accordingly, the 1 arcmin size portion of the diffuse nebula radiates $`2.1\times 10^4\dot{E}`$ into the $`0.510`$ keV X-ray band and $`2.3\times 10^4\dot{E}`$ into the $`0.12.4`$ keV soft X-ray band.
For a second circular region which is located $`2.5`$ arcmin behind the pulsar along its proper motion direction we converted the background and vignetting corrected counting rate to an energy flux by assuming that the power law photon index is similar to the one we modeled from region one. This assumption may not be justified. For example, Willingal et al. (2001) found in a detailed spectral analysis of the Crab nebula that its outer regions show the steepest spectrum. This indicates enhanced synchrotron losses of the electrons during their passage from the pulsar to the outskirts of the nebula. It is conceivable that a similar behavior is valid for the electrons radiating in the trail of PSR B1929$`+`$10. The energy fluxes and luminosities we obtain for the second region are upper limits in this respect. For the $`0.510`$ keV and $`0.12.4`$ keV bands we find $`f_x1.8\times 10^{14}\text{erg s}^1\text{cm}^2`$, $`f_x1\times 10^{14}\text{erg s}^1\text{cm}^2`$ and $`L_x2.8\times 10^{29}\text{erg s}^1`$, $`L_x1.6\times 10^{29}\text{erg s}^1`$, respectively, which is $`47\times 10^5\dot{E}`$.
Wang, Li & Begelman (1993) investigated the trail emission by assuming a synchrotron emission spectral shape for which they simulated a spectrum and compared model predictions vs. ROSAT observed counting rates. Our findings based on a more detailed spectral modeling agree well with their results.
### 3.3 MULTI-WAVELENGTH SPECTRUM
In order to construct a broadband spectrum combining all spectral information available from PSR B1929$`+`$10 we adopted the radio spectrum from Maron et al. (2000) and plotted it in Figure 15 together with the XMM-Newton observed pulsar spectrum and the optical fluxes obtained from the Hubble Space Telescope observations by Pavlov et al. (1996) and Mignani et al. (2002).
PSR B1929$`+`$10 belongs to the small group of pulsars for which radio emission was detected up to 43 GHz (Kramer et al. 1997). The radio spectrum of PSR B1929$`+`$10 is a power law with a photon index of $`\alpha =2.6\pm 0.04`$ in the frequency range $`0.424`$ GHz (Maron et al. 2000). Its energy flux at 100 MHz is 950 mJy$`\pm `$ 600 mJy, but at 43 GHz it is still 0.18 mJy $`\pm `$ 0.05 mJy (which is slightly more than extrapolated from the lower frequency data). The flux density in the radio part of the spectrum, which should be due to coherent radiation, is therefore several orders of magnitude greater than the extrapolated optical or X-ray flux densities.
Extrapolating the power law spectrum which describes the XMM-Newton data to the optical bands yields a photon flux which exceeds the flux measured in the near-UV bands by more than an order of magnitude. While the pulsar was clearly detected in the U-band (m$`{}_{342W}{}^{}=25.7`$) and with the F130LP (m$`{}_{130LP}{}^{}=26.9`$) and STIS F25QTZ filters, only an upper limit is available in B (m$`{}_{430W}{}^{}26.2`$) (Pavlov et al. 1996; Mignani et al. 2002). This suggests that the broadband spectrum, if entirely non-thermal, has to break somewhere before or in the soft channels of the X-ray spectrum. To test this hypothesis we have fitted a broken power law model to the XMM-Newton and optical data. A broken power law model, with $`E_{break}=0.83_{0.03}^{+0.05}\text{keV}`$, provides an excellent description ($`\chi ^2`$= 121.16 for 120 dof) of both spectral data sets. The photon index for energies below and above the break $`E_{break}`$, is found to be $`\alpha _1=1.12_{0.03}^{+0.02}`$ and $`\alpha _2=2.48_{0.07}^{+0.08}`$, respectively (with a normalization of $`8.2_{0.7}^{+1.0}\times 10^5`$ photons cm<sup>-2</sup> s<sup>-1</sup> keV<sup>-1</sup> at 1 keV). The model predicted column absorption value of $`N_H=0.52_{0.06}^{+0.12}\times 10^{21}\text{cm}^2`$, somewhat smaller than the single power law model prediction, but is still compatible with our fit to the diffuse trail emission.
## 4 SUMMARY & DISCUSSION
We have investigated the emission properties of PSR B1929$`+`$10 and its X-ray trail in a multi-wavelength study using XMM-Newton, the ESO New Technology Telescope (NTT) in La Silla (Chile), the Hubble Space Telescope, the Effelsberg 100-m Radio Telescope and the Jodrell Bank Radio Observatory.
In X-rays, PSR B1929$`+`$10 is the brightest pulsar among the old rotation-powered pulsars. Therefore, for about 10 years it was the only member of its class detected in X-rays and for which some details on its emission properties were known โ though this picture has been revised thoroughly in our multi-wavelength study of the old rotation-powered pulsars (cf. Becker et al. 2004; 2005).
Clearly, the study of old rotation-powered pulsars is in the domain of XMM-Newton which was designed and built to study faint sources in the X-ray sky. Its sensitivity allowed, for the first time, a more detailed study of this interesting group of pulsars. As for PSR B0628-28 (Becker et al. 2005), PSR B0950+08, PSR B0823+26 and PSR J2043+2740 (Becker et al 2004; Pavlov & Zavlin 2004) recently observed by XMM-Newton, we found the temporal and spectral X-ray emission properties of PSR B1929$`+`$10 to be in excellent agreement with a non-thermal, magnetospheric emission scenario. Its X-ray spectrum is best described by a single power law model with a photon index of $`2.72_{0.09}^{+0.12}`$. Using the best fit composite Planckian power law model, the contribution from thermal emission of heated polar caps is inferred to be at most $`7\%`$. However, a pure thermal emission spectrum consisting of two Planckian spectra is regarded as unlikely. Taking the optical spectral data into account a broken power law with $`E_{break}=0.83_{0.03}^{+0.05}\text{keV}`$ and the photon-index $`\alpha _1=1.12_{0.03}^{+0.02}`$ and $`\alpha _2=2.48_{0.07}^{+0.08}`$ is able to describe the emission in both spectral ranges entirely in terms of a non-thermal magnetospheric origin.
The X-ray pulse profile observed in the $`0.210`$ keV band is found to be markedly different from the broad sinusoidal pulse profile seen in the low statistic ROSAT data. Fitting Gaussians to the X-ray light curve indicates the possible existence of three pulse components. A small narrow pulse, characterized by energies greater than 1 keV, is found to lead the radio main pulse by $`20^{}`$. Two larger pulses, observed in all three energy bands, follow this small pulse. These three pulses are roughly separated by about the same phase cycle (cf. Figure 8). The fraction of pulsed photons in the $`0.210`$ keV band is $`32\pm 4\%`$. For the sub-bands $`0.21.0`$ keV and $`1.02.1`$ keV the pulsed fraction is $`24\pm 5\%`$ and $`44\pm 6\%`$, respectively, indicating a mild energy dependence at a $`2\sigma `$ level.
Various theoretical models have been developed to explain the observed non-thermal high energy emission properties of younger pulsars like those in the Crab and Vela supernova remnants. They all appear to be seen not only in X-rays but also in the gamma-ray band. Can the generic features of these models also explain the emission properties we observe from PSR B1929$`+`$10?
It is commonly believed that the non-thermal X-ray photons are emitted by relativistic charged particles in the pulsar magnetosphere. These relativistic particles could be accelerated in a polar cap region (e.g. Harding 1981; Zhang & Harding 2000) or in the outer-magnetosphere (e.g. Cheng, Ho & Ruderman 1986). In order to calculate the non-thermal spectrum and the energy dependent light curves, the amount of current flow and detailed three dimensional geometry of the accelerator are required. To fix these parameters, gamma-ray data is normally required. However, if the inclination and viewing angles of a specific pulsar are known, the qualitative features of the light curve can be predicted according to the three dimensional magnetospheric models (e.g. Yadigaroglu & Romani 1995; Cheng, Ruderman & Zhang 2000; Dyks, Harding & Rudak 2004).
Everett & Weisberg (2001) have reported from radio polarization data that the inclination and viewing angles of PSR B1929$`+`$10 are $`36^{}`$ and $`26^{}`$, respectively. For such small inclination and viewing angles, at most two pulses with a phase separation of $`180^{}`$ can be produced if there is only outgoing current. In this case, one pulse arises from a region near one polar cap and another pulse arises from a region near the light cylinder, but associated with another magnetic pole. In Figure 9, the first narrow peak only appears in the band $`1.02.1`$ keV and could be associated with a heated polar cap of kT $``$ 1keV. However, this is very unlikely by its narrowness as strong gravitational light bending would smear out the thermal pulse from the surface (e.g. Page 1995). Cheng, Ruderman & Zhang (2000) have shown that incoming current must exist from null charge surface to the stellar surface due to pair creation in the outer magnetosphere accelerator.
Here, we can simulate the light curves of PSR B1929$`+`$10 by using its observed inclination and viewing angles, and the three dimensional outer gap model (Cheng, Ruderman & Zhang 2000). Since the radiation is expected to be emitted from open field lines, the coordinate values ($`x_0,y_0,z_0`$) of the last closed field lines at the stellar surface must be determined. The coordinate values ($`x_0^,,y_0^,,z_0^,`$), where $`x_0^,`$=$`a_1x_0`$, $`y_0^,`$=$`a_1y_0`$ and $`z_0^,`$=$`[1(x_0^{,2}+y_0^{,2})]^{1/2}`$, then represent an open field line surface for a given value of $`a_1`$. For simplicity, we choose $`a_1=0.97`$, which is very close to the first open field lines ($`a_1=1.0`$). In Figure 16, we show that two large pulses and one small pulse can be simulated. The first large pulse is the result of the incoming current toward the south pole while another pulse is produced by the outgoing current from the north pole. The small pulse is produced by outgoing current near the light cylinder. We note that the phase separation of these three pulses roughly corresponds to the phase separation of the three pulse components observed in the X-ray light curve. Here, the (x,z) plane is chosen to lie at the zero-phase position in Figure 16. On the other hand, in Figure 7, the (x,z) plane should lie in the middle between the two radio pulses. Therefore, the simulated X-ray light curves roughly reproduce the phase relative to the radio pulses within a phase error of about $`\pm 0.1`$ which is well within the observed uncertainties.
#### THE TRAIL OF PSR B1929$`+`$10
The existence of diffuse emission with a trail morphology lying in the direction opposite to the motion of the pulsar is confirmed in our XMM Newton observation and provides a unique opportunity to probe the pulsar and its environment. The extended diffuse emission associated with this old pulsar (with spin down energy less than $`10^{34}`$ ergs s<sup>-1</sup>) indicates that spin down power is not the sole criterion for its detection. In addition to distance, the detectability may also be dependent on the pulsarโs transverse velocity (see Chatterjee & Cordes 2002). The existence of a possible radio counterpart in the 11 cm Effelsberg data is exciting and, if confirmed in subsequent observations, can provide important information on the trail properties. No diffuse emission from the pulsar trail is detected in $`H_\alpha `$ perhaps suggesting that the neutral component of the interstellar medium is low (see Chatterjee & Cordes 2002) in the environment surrounding PSR B1929$`+`$10.
Since the trails X-ray emission near to the pulsar has a hard spectrum characterized by a power law photon index of $`2`$, the emission is non-thermal and is likely to be produced from the synchrotron process of highly relativistic electrons in the shocked region between the pulsar wind and surrounding interstellar medium. A physical description of PSR B1929$`+`$10โs X-ray trail based on the ROSAT findings has been discussed by Wang, Li, & Begelman (1993) in terms of an outflow collimated within the pulsarโs cavity created by its motion.
Alternatively, the properties of the distorted wind nebula can be inferred under the assumption that the electron lifetime due to synchrotron losses, $`\tau _{syn}`$, is comparable to the timescale for the passage of the pulsar over the length of its X-ray trail. Such a hypothesis has also been considered, for consistency, by Caraveo et al. (2003) in their model analysis of a similar X-ray trail observed in the Geminga pulsar. For PSR 1929+10, the angular extent of the trail can not be definitively determined with the present data, but it is likely greater than 4 arcmin. For an assumed distance of 361 pc, the linear scale, $`d`$, corresponding to this angular scale is greater than about 0.4 pc. An estimate of the flow time, $`t_{flow}`$, can be obtained, for the proper motion measured by Chatterjee et al. (2004) and resulting velocity of $`v_p`$ of 177 km s<sup>-1</sup>, leading to $`t_{flow}2200`$ yrs. This lower limit is about a factor of 2.3 times longer than the comparable timescale found for the Geminga pulsar (see Caraveo et al. 2003). To determine the consistency of our interpretation, the magnetic field in the shocked region can be estimated by equating $`\tau _{syn}`$ to $`t_{flow}`$. Here $`\tau _{syn}=6\pi m_ec/\gamma \sigma _TB^210^5B_{\mu G}^{3/2}(h\nu _X/\mathrm{keV})^{1/2}`$ yr where $`\gamma `$ is the Lorentz factor of the wind, taken to be equal to $`10^6`$, $`\sigma _T`$ is the Thompson cross section, and $`B_{\mu G}`$ is the magnetic field in the emission region in micro gauss. The inferred magnetic field strength in the emitting region is $`12\mu `$G. Given the magnetic field strength estimates in the interstellar medium ($`26\mu `$G; see Beck et al. 2003), and the expected compression of the field in the termination shock by about a factor of 3 (Kennel & Coroniti 1984), our estimates for the magnetic field in the emitting region of the pulsar wind nebula are in approximate accord.
The X-ray luminosity and the spectrum of the emitted radiation can be estimated using a simple, one-zone model for the emission nebula powered by the pulsar wind as developed by Chevalier (2000). To determine the characteristic properties of the emission a comparison of the cooling frequency, $`\nu _c`$, for which the electrons can radiate their energy in the pulsar trail, with the observing frequency, $`\nu _X`$ is necessary. The cooling frequency can be expressed as $`\nu _c=\frac{e}{2\pi m_ecB^3}(\frac{6\pi m_ec}{\sigma _T\tau _{syn}})^2`$. Substitution of the inferred magnetic strength and the electron cooling timescale due to synchrotron radiation yields $`\nu _c=1.8\times 10^{17}`$ Hz. Since $`\nu _X>\nu _c`$, the electrons are able to radiate their energy in the pulsar trail and the photon index, $`\mathrm{\Gamma }`$, is related to the power law index, $`p`$, of the electron energy distribution, $`N(\gamma )\gamma ^p`$, in the form $`\mathrm{\Gamma }=(p+2)/2`$. Based on the theoretical work on highly relativistic shocks (Bednarz & Ostrowski 1998, Lemoine & Pelletier 2003), we adopt $`p=2.2`$, yielding $`\mathrm{\Gamma }=2.1`$. To be consistent with our interpretation of emission taking place in the fast cooling region, the observed value of $`\mathrm{\Gamma }`$ should exceed 2, which is consistent with the observed value of $`\mathrm{\Gamma }=2_{0.4}^{+0.4}`$.
The luminosity of the radiating electrons in the nebula can be calculated from the luminosity per unit frequency given by
$$L_\nu =\frac{1}{2}(\frac{p2}{p1})^{p1}(\frac{6e^2}{4\pi ^2m_ec^3})^{(p2)/4}ฯต_e^{p1}ฯต_B^{(p2)/4}\gamma _w^{p2}$$
$$R_s^{(p2)/2}\dot{E}^{(p+2)/4}\nu ^{p/2}.$$
(1)
where $`R_s`$ is the distance of the shock from the pulsar expressed as $`R_s=(\dot{E}/2\pi \rho v_p^2c)^{1/2}10^{16}\dot{E}_{33}^{1/2}n^{1/2}v_{p,100}^1\mathrm{cm}`$. Here $`v_{p,100}`$ is the velocity of the pulsar in units of 100 km s<sup>-1</sup>, $`\dot{E}_{33}`$ is the spin down power of the pulsar in units of $`10^{33}`$ ergs s<sup>-1</sup>, and $`n`$ is the number density of the interstellar medium in units of 1 cm<sup>-3</sup>. The spin down power of $`3.89\times 10^{33}`$ ergs s<sup>-1</sup> and a density of 1 cm<sup>-3</sup> yields a shock radius of $`R_s6\times 10^{15}`$ cm. Assuming energy equipartition between the electron and proton fractional energy densities so that $`ฯต_e0.5`$, and a fractional energy density of the magnetic field $`ฯต_B0.01`$ (see Cheng, Taam, & Wang 2004), the corresponding luminosity given as $`\nu L_\nu `$ is $`1.6\times 10^{30}`$ ergs s<sup>-1</sup> or $`4\times 10^4\dot{E}`$. In view of the observational uncertainties, this is consistent with the observed values of $`8.3_{3.4}^{+8.7}\times 10^{29}`$ ergs s<sup>-1</sup> ($`0.510`$ keV) and $`9.1_{3.6}^{+8.6}\times 10^{29}`$ ergs s<sup>-1</sup> ($`0.12.4`$ keV).
Future theoretical investigations should be carried out to confronting the observed X-ray lightcurves and spectra having even better photon statistics than the ones which we have obtained in this first XMM-Newton observations of PSR B1929$`+`$10. The results we obtained from the X-ray trail along with the possible discovery of its radio counterpart are exciting. Follow up radio observations of the trail region at different wavelengths are currently scheduled for fall 2006 in order to further constrain its existence and, if detected, provide polarization information and the spectral index in the radio regime from it.
We acknowledg he use of data obtained from XMM-Newton observations which is an ESA science mission with instruments and contributions directly funded by ESA Member States and NASA. WB would like to thank Hermann Brunner for his help in the X-ray image reconstruction. We are grateful to Olaf Maron for the use of his database of pulsar radio fluxes, to E. Fรผrst for help with the 11cm Effelsberg radio survey and to Armando Manzali for his contribution to the the phase resolved spectral analysis which is part of the thesis work. KSC and JJJ are partially supported by a RGC grant of Hong Kong Government. AS was supported by KBN grant PBZ-KBN-054/P03/2001. ADL acknowledges a fellowship by the Italian Space Agency, ASI. PAC and ADL acknowledge financial contribution from contract ASI-INAF I/023/05/0. Optical observations were made with ESO Telescopes at the La Silla Observatories under program 077.D-0794(A). This work was supported in part by the Theoretical Institute for Advanced Research in Astrophysics (TIARA) operated under Academia Sinica and the National Science Council Excellence Projects program in Taiwan administered through grant number NSC 94-2752-M-007-001. |
warning/0506/hep-ph0506180.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The aim of this contribution is to compare whatever the best we have at hand for evaluation of the initial state radiation (ISR) effect in the process $`e^{}e^+\mu ^{}\mu ^+\gamma `$ using KKMC and PHOKHARA Monte Carlo programs. The above investigation will be partly extended to the process $`e^{}e^+\pi ^{}\pi ^+\gamma `$.
## 2 ISR in muon pair production
Both KKMC and PHOKHARA programs are full scale MC event generators, which can provide for any experimentally observable distribution. We concentrate however, on the distribution of the squared mass spectrum $`Q^2=s^{}`$ of the muon pair, because this distribution is relevant for the radiative return measurements of $`R(s)`$, and also because this particular distribution we may compare with the classical semi-analytical calculations. Here we shall also exploit the analytical formulas of ref. (see also ), which implement analytical second order ISR calculation of ref. and third order leading-logarithmic (LL) ISR calculation of refs. . The ISR formula of ref. is provided by the KKsem facility of KKMC. In the actual KKsem implementation we use version of the formula where numerically negligible (at least at LEP energies, see ref. ) second order NNLL terms are neglected.
It is important to stress from the very beginning that authors of PHOKHARA and KKMC use different terminology to describe Born level matrix element and higher order matrix element. I shall not try to unify terminology or fully explain the differences, referring the reader to original works, like refs. and . Let me explain only very briefly the main differences. The KKMC authors define Born as $`e^+e^{}f\overline{f}`$ without any photon emission and the radiative return is necessarily the first order process with respect to such a Born level. The leading-logarithmic (LL) corrections are of order $`\alpha ^nL^n`$, where $`n=1,2,3\mathrm{}\mathrm{}`$ is the standard perturbative order, while mass logarithm $`L=\mathrm{ln}(s/m_e^2)`$ is coming either from the virtual photon correction or the phase space integration over the real photon angle down to zero value. The NLL and NNLL corrections are of order $`\alpha ^nL^{n1}`$ and $`\alpha ^nL^{n2}`$ correspondingly. Concerning mass terms, they are routinely neglected in KKMC for the electron (except those which integrate up to a finite correction) while an effort is made to keep all of them for the final state fermions, at least at the Born and the first order level. KKMC implements several variants of the QED matrix elements, which feature different level of higher order and mass term truncation. PHOKHARA authors employ the leading-order (LO) as a name for the process in which one (and only one) photon is emitted in the final state. They name as the next-to-leading-order (NLO) their matrix element with the one-loop corrections and the second real photon. This terminology may seem more adequate to discuss radiative return. However, when trying to match the two terminologies one has to pay attentions to the available phase space of the first real photon. Depending on whether the minimum emission angle is imposed or not, one gets full factor $`L`$ or not, even at the LO. This affects strongly the relative magnitude of higher order corrections with respect to LO or LL. In this study we generally exclude from the considerations โnon-photonicโ corrections due to emission of additional lepton pair and vacuum polarization.
In Fig.1 we compare results of KKMC and of PHOKHARA using the best available ISR matrix element in both programs at $`\sqrt{s}=1.01942`$GeV. In KKMC we use second order matrix element with coherent exclusive exponentiation (CEEX) described in refs. . The second order CEEX matrix element has complete next-to-leading-logarithmic (NLL) contributions<sup>1</sup><sup>1</sup>1For unpolarized beams, see discussion in ref. related to eq. (128) therein. and complete next-next-to-leading-logarithmic (NNLL) contributions. The magnitude of NLL and NNLL corrections was also examined in a separate studies, see contribution of S. Yost in these proceedings. KKMC includes most of the third order LL contributions by the virtue of exponentiation<sup>2</sup><sup>2</sup>2 It is also known that exponentiation of the YFS type sums up quite substantial part of third order LL, see refs. . On the other hand, PHOKHARA implements complete second order ISR, including complete NLL and NNLL corrections (i.e. singular corrections proportional to $`\frac{\alpha }{\pi }m_e^2`$ and $`\frac{\alpha }{\pi }m_e^4`$, which integrate to finite corrections of order $`\frac{\alpha }{\pi }`$, in the limit $`m_e0`$). PHOKHARA does not resum (exponentiate) soft photon contributions to infinite order. It is worth to stress that the two MC calculation, KKMC and PHOKHARA, and semianalytical formula of KKsem of refs. represent set of three completely independent second order (using terminology of KKMC) calculation of the ISR in every aspect of calculating QED matrix element and integrating the phase space.
The main comparison of the ISR calculations is shown in Fig.1a, where the distribution $`d\sigma /dQ^2`$ from KKMC and KKsem agree very well, within 0.2%, except very low $`Q^2`$ where they diverge by about 0.3% <sup>3</sup><sup>3</sup>3 Note that similar comparison of KKMC and KKsem was done in ref. for LEP energies. At the present lower energy $`\sqrt{s}=1.01942`$GeV subleading terms are, however, more important., while Fig.1a shows certain addition cross-check. The reason for this discrepancy is not clear. Neglected NNLL in KKsem are a viable candidate, but to confirm this hypothesis one would need more tests. In the same plot we see that PHOKHARA agrees well with KKsem at low $`Q^2`$ (aligning with KKsem) and differs by about 0.25% in the central region (we need higher statistics from PHOKHARA to confirm this number) from both KKMC and KKSEM and drops sharply at soft limit, high $`Q^2`$, because of lack of soft photon resummation. In order to understand quantitatively the effect of lack of exponentiation in PHOKHARA we compare its result in Fig.1b with a variant of KKsem in which we switch off exponentiation, i.e. all terms beyond second order are truncated. The smooth curve in Fig.1b representing result of this truncation fits very well PHOKHARA result. In particular, looking into this result, one may think that the deviation of PHOKHARA by 0.25% in the central region is related to its neglect of the third order LL. This conjecture needs more test to be confirmed.
We summarize on the results of Fig.1 that KKMC with the second order CEEX matrix element, PHOKHARA with its second order matrix element and KKsem implementing second order analytical calculation agree very well, within the expected range and the pattern of the discrepancies seems to be understood.
In KKMC there is another more primitive QED matrix element denoted as EEX, see ref. for its full description, which follows closely the classical Yennie-Frautschi-Suura (YFS) exponentiation scheme and its implementation is limited to first order plus second order LL. In the second order EEX matrix element (contrary to CEEX) the NLL corrections are incomplete. (On the other hand EEX third order LL is complete, while in CEEX it is incomplete.). For technical and historical reasons, see discussion below, EEX type matrix element is used for the production of low energy hadronic final states, for example for pion pair. It is, therefore, important to check how good it is compared to KKMC with more complete coherent exclusive exponentiation (CEEX) matrix element. This is done for the muon final state in Figs. 2. (CEEX is not yet available for $`\pi `$-pairs). In Fig. 2a we see results from KKMC CEEX and several variants of EEX. We are actually plotting $`d\sigma /dQ^2`$, dividing all results by KKsem of ref. , the same as in previous Fig.1. The curves marked EEX72 represent exponentiated EEX matrix element based on complete first order, while EEX73 and EEX74 include also complete second and third LL, while CEEX203<sup>4</sup><sup>4</sup>4Indices 203, 74 etc. follow numbering of MC weights in ref. . is the same as in Fig. 1. As we see, at low $`Q^2`$, that is for the hard photon emission, results of EEX matrix elements depart from other more complete results by up to 3%! In the $`\rho `$ region it is different from the KKsem, CEEX KKMC by about 1%. The above result is also consistent with what we have seen in Fig. 1. EEX is therefore not well suited for the use in the high precision measurements of $`R(s)`$ using radiative return below $`Q^2`$=1GeV. This result is not very much surprising, as EEX of KKMC has incomplete second order NLL. The observed effect at the low $`Q^2`$ is a little bit bigger then what we expected. We have therefore done certain additional tests. We have split EEX results into three components, $`\stackrel{~}{\beta }_i,i=0,1,2`$, compared each of them with analytical result of table I of ref. , also at $`\sqrt{s}=10`$GeV. We do not show results of these tests here, but the overall pattern of discrepancies seems to be consistent with NLL class of corrections. This additional test indicates also that the main source of the problem is an approximate double real emission matrix element in EEX and not the incomplete virtual corrections. In particular we have included in these tests the complete NLL contribution in $`\stackrel{~}{\beta }_0`$ and $`\stackrel{~}{\beta }_0`$. This did not help! The whole discrepancy seems to result from the use of the LL-approximate matrix element for the double real photon emission in EEX. The above observation is consistent with the older tests in ref. at LEP energies.
### 2.1 Muon pair, ISR+FSR
Let us not include FSR in the game, again for the muon pair final state. In fact, at low $`Q^2`$ the rate of muon pair in radiative return is higher than of $`\pi `$ pairs, hence $`d\sigma /dQ^2`$ of muons can be used as a reference distribution for measuring $`R(s)`$. It is therefore worth to test FSR in KKMC and to check result for $`d\sigma /dQ^2`$ once again. In Fig. 2b we show result from KKMC for second order CEEX matrix element in which we include ISR, FSR and its interference. We compare MC result with the semianalytical result of KKsem in which with the same ISR radiative function of ref. . The FSR distribution of ref. features incomplete NLL in KKsem, so it is definitely inferior with respect to ISR counterpart โ the complete list of the FSR radiative corrections in KKsem can be found in Table II in ref. . This above semianalytical formula also misses the interference of ISR and FSR, which in first order is zero in the inclusive $`d\sigma /dQ^2`$ so this omission does not harm. In Fig. 2a we see the ratio of the corresponding results from KKMC and KKsem. (NB. PHOKHARA is able to provide result with FSR for muon pairs in the LO, and it would be interesting to include it in the comparison.) This result is rather preliminary and has to be checked. In any case, the agreement better than 1% found all over the $`dQ^2`$ range is quite satisfactory as a starting point for further investigation<sup>5</sup><sup>5</sup>5In one bin we see trace of large weight fluctuation which is probably due to rounding errors. This result was obtained using weighted events. For the MC run with weight-one events this effect would disappear. Such numerical instabilities need further investigation.
## 3 ISR for $`\pi ^+\pi ^{}`$ pair production
Let us now switch to low $`Q^2`$ $`\pi `$-pair state produced at the radiative return process. In Fig. 3 we compare KKMC with the EEX style matrix element on one hand with PHOKHARA second order (marked as PHOKHARA2) on the other hand. The EEX matrix element is the default one in KKMC, with first order exponentiation, the completed second and third order LL (EEX74). We limit ourselves to ISR only. The results of Figs. 3a-b are obtained without any cutoffs. In Fig. 3a we show the actual distributions, including also the distribution for the muon pairs. We see that for $`Q^2<0.33`$ the muon-pair cross section is bigger than that of $`\pi `$-pair. The ratio PHOKHARA/KKMC is not so well understood as the the analogous results for the muon pair shown in the previous section. The discrepancy at high $`Q^2`$ we attribute to lack of exponentiation in PHOKHARA<sup>6</sup><sup>6</sup>6This does not hinder practical applications of PHOKHARA for radiative return measurements, which concentrate at lower $`Q^2`$. while another larger discrepancy at low $`Q^2`$ is most likely due to incompleteness of second order NLL in EEX matrix element on KKMC, and it corresponds to deviation which was already seen in Fig. 2a. In Fig. 3c-d we show the analogous results for relatively mild cut on photon momentum, where photon is defined as a โmissing four momentumโ calculated knowing pion momenta and beam momenta. We ask for the momentum of such a โcollective unseen photonโ to be directed below $`15^{}`$ from the beam and to have at least $`10MeV`$ of energy. For each $`\pi `$ we require that it is situated in wide angles, eg. separated by more than $`40^{}`$ from each beam. Results in Fig. 3c-d look quite similar, except that the discrepancy between PHOKHARA and EEX KKMC is bigger (we need better statistics from PHOKHARA to see it more clearly). This can be attributed to the fact that the leading logarithm $`L`$ due to real emission is diminished by the cut on the photon angle with respect to beams.
### 3.1 How to extend CEEX ISR to hadronic final states?
In the following we show that the superior CEEX ISR matrix element can be extended to hadronic final states at low $`Q^2`$, like pion pair. This can be done provided we have some decent modeling of the hadronic final state in terms of the corresponding formfactor. In view of its practical importance, let us elaborate on this point.
In CEEX Born amplitude for $`ee\mu \mu `$ is defined as a four-spinor tensor
$$\begin{array}{cc}& ๐
\left({}_{\lambda }{}^{p};X\right)=๐
\left({}_{\lambda _a}{}^{p_a}{}_{\lambda _b}{}^{p_b}{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d};X\right)=๐
[{}_{\lambda _b}{}^{p_b}{}_{\lambda _a}{}^{p_a}][{}_{\lambda _c}{}^{p_c}{}_{\lambda _d}{}^{p_d}](X)=๐
_{[ba][cd]}(X)=\hfill \\ & =ie^2\underset{B=\gamma ,Z}{}\mathrm{\Pi }_B^{\mu \nu }(X)(G_{e,\mu }^B)_{[ba]}(G_{f,\nu }^B)_{[cd]}H_B=\underset{B=\gamma ,Z}{}๐
_{[bc][cd]}^B(X),\hfill \\ & (G_{e,\mu }^B)_{[ba]}\overline{v}(p_b,\lambda _b)G_{e,\mu }^Bu(p_a,\lambda _a),(G_{f,\mu }^B)_{[cd]}\overline{u}(p_c,\lambda _c)G_{f,\mu }^Bv(p_d,\lambda _d),\hfill \\ & G_{e,\mu }^B=\gamma _\mu \underset{\lambda =\pm }{}\omega _\lambda g_\lambda ^{B,e},G_{f,\mu }^B=\gamma _\mu \underset{\lambda =\pm }{}\omega _\lambda g_\lambda ^{B,f},\omega _\lambda =\frac{1}{2}(1+\lambda \gamma _5),\hfill \\ & \mathrm{\Pi }_B^{\mu \nu }(X)=\frac{g^{\mu \nu }}{X^2M_{B}^{}{}_{}{}^{2}+i\mathrm{\Gamma }_BX^2/M_B},\hfill \end{array}$$
(1)
and it enters as a basic building block in every spin amplitude in the CEEX scheme, with arbitrary number of photons. See eq. (43) in ref. for notation. The above Born is calculated using Chisholm identity and replaced with the bi-spinor objects of the Kleiss-Stirling method.
In case of hadronic final state the structure of the Born amplitude is
$$๐
_{[ba]}^\mu (X)J_\mu (X,q_i),J_\mu (X,q_i)X^\mu =0,$$
(2)
where $`q_i`$ are momenta of the final state hadrons, $`X=q_i`$, and
$$\begin{array}{cc}& ๐
_{[ba]}^\nu (X)=ie^2\underset{B=\gamma ,Z}{}H_B(G_{e,\mu }^B)_{[ba]}\mathrm{\Pi }_B^{\mu \nu }(X)=\underset{B=\gamma ,Z}{}๐
_{[bc]}^{B\mu }(X).\hfill \end{array}$$
(3)
In the rest frame of $`X`$ one has $`J^\mu =(0,\stackrel{}{J})`$ and we may split $`J`$ into difference of the two massless four-vectors $`J^\mu =J_+^\mu J_{}^\mu `$. In the arbitrary reference frame the above prescription extends as follows<sup>7</sup><sup>7</sup>7The author would like to thank J. Kuehn for suggesting him this solution.
$$J_\pm ^\mu =\frac{1}{2\sqrt{X^2}}(\sqrt{J^2}X^\mu \pm \sqrt{X^2}J^\mu ).$$
(4)
Each of the two corresponding components in $`๐
`$ can be expressed in terms of the of the Kleiss-Stirling bi-spinors $`s_\pm (p_1,p_2)`$, see eqs. (A4-A6) in ref. . This can be done using completeness relation for $`\overline{v}(p_b,\lambda _b)\gamma _\nu J_\pm ^\nu u(p_a,\lambda _a)`$ taking advantage of $`J_\pm ^2=0`$.
Note that the above CEEX implementation requires that we parametrize the production amplitude of the each final hadronic state one by one in terms of the formfactors in a completely exclusive manner. However, this is necessary anyway for good phenomenological description of these often resonant low energy hadronic states. We conclude that CEEX can be used for ISR for low energy hadron production. The question is only how much programming it will be and who will do it.
### 3.2 Conclusions
The case of radiative return with the muon pair final state is well understood, especially for the ISR where three different second order calculations agree very well. The case with FSR requires more tests. Matrix element EEX of KKMC with the incomplete second order NLL is not good in the radiative return at $`Q^2<1GeV`$ for precision requirement better than 1%. Method of porting CEEX matrix element of KKMC to pion pair final state is outlined.
Acknowledgments
This work is partly supported by TARI Contract No. RII3-CT-2004-506078. The author would like to thank J. Kuehn and A. Denig for useful discussion. Warm hospitality at TTP Karlsruhe University and INFN Frascati, where part of this work was done is also acknowledged. |
warning/0506/math-ph0506065.html | ar5iv | text | # Square lattice Ising model susceptibility: connection matrices and singular behavior of ๐โฝยณโพ and ๐โฝโดโพ
## 1 Introduction
Since a pioneering, and quite monumental, paper on the two-dimensional Ising models, it has been known that the magnetic susceptibility of square lattice Ising model, can be written as an infinite sum of $`(n1)`$-dimensional integrals contributions:
$`\chi (T)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\chi ^{(n)}(T)`$ (1)
The odd (respectively even) $`n`$ correspond to the high (respectively low) temperature domain. These $`(n1)`$-dimensional integrals are known to be holonomic, since they are integrals of holonomic (actually algebraic) integrands. Besides the known $`\chi ^{(1)}`$ and $`\chi ^{(2)}`$ terms, which can be expressed in terms of simple algebraic or hypergeometric functions, it is only recently that the Fuchsian differential equations satisfied by the $`\chi ^{(3)}`$ and $`\chi ^{(4)}`$ terms have been found . These two exact differential equations of quite large orders (seven and ten) can be used to find answers to a set of problems traditionally known to be subtle, and difficult, for functions with confluent singularities, like the fine-tuning of the singular behaviors for all the singularities (dominant singular behavior, sub-dominant, etc.), accurate calculations of the asymptotic behavior of the coefficients, etc.
Recall that the third, and fourth, contribution to the magnetic susceptibility $`\chi ^{(3)}`$, and $`\chi ^{(4)}`$, are given by multi-integrals and each is, thus, a particular solution of the corresponding differential equation. These differential equations exhibit a finite set of regular singular points that may (or may not) appear in the physical solutions $`\chi ^{(3)}`$ and $`\chi ^{(4)}`$. Besides the physical singularities and the non physical singularities $`s=\pm i`$ (where $`s=\mathrm{sinh}(2K)`$, $`K`$ being the usual Ising model coupling constant, $`K=\beta J`$), it is commonly believed that the $`\chi ^{(n)}`$โs have, at least, other non physical singularities given by B. Nickel . The dominant singular behaviors at all these (non physical) singularities ($`\chi ^{(3)}`$ and $`\chi ^{(4)}`$) have also been given by B. Nickel. The differential equations of the $`\chi ^{(n)}`$โs, which โencodeโ all the information on the solutions and their singular behavior, in fact, allow us to obtain not only the dominant, but also all the subdominant singular behavior, hardly detectable from straight series analysis. It is thus of interest to get (or confirm) these singular behaviors from the exact Fuchsian differential equations that we have actually obtained for $`\chi ^{(3)}`$ and $`\chi ^{(4)}`$ and, especially, the singular behavior at the two new quadratic singularities, $`\mathrm{\hspace{0.17em}1}+3w+4w^2=\mathrm{\hspace{0.17em}0}`$, (where $`w=s/(1+s^2)/2`$) found for $`\chi ^{(3)}`$ .
The physical solution $`\chi ^{(3)}`$ is defined by a double integral on two angles and is known as a series obtained by expansion (then integration) of the double integral at $`w=0`$ (or $`s=0`$). It is certainly not simple to obtain the $`\chi ^{(3)}`$ expansion around (say) the ferromagnetic critical point $`w=1/4`$, due to a singular logarithmic behavior. However, one can overcome this difficulty since, with a differential equation, it is straightforward to obtain the formal series solutions at each regular singular point (i.e., a local basis of series solutions). By connecting the formal solutions around $`w=0`$ and the formal series solutions around another regular singular point like $`w=1/4`$, one will be able to express the particular solution $`\chi ^{(3)}`$ (and also all the other formal solutions) as a linear combination of solutions valid at $`w=1/4`$. The seven local solutions at $`w=0`$ will, then, be given by the product of a $`7\times 7`$ matrix with the vector having the seven local solutions at $`w=1/4`$ as entries. In other words, succeeding in obtaining these connection matrices amounts to building a common (global) basis of solutions valid for all the regular singular points. Furthermore, with these connection matrices, we obtain, in fact, the analytic continuation in the whole complex plane of the variable $`w`$, of $`\chi ^{(3)}`$ and $`\chi ^{(4)}`$, which are known as integral representations.
Note that, remarkably, the Fuchsian differential equation for $`\chi ^{(3)}`$ has simple rational, and algebraic, solutions. These rational or algebraic solutions, known in closed form, can be understood globally. One can easily expand such globally defined solutions around any singular point of the ODE, and follow these solutions through any โjumpโ from one regular singularity to another one, and, therefore, from one well-suited basis to another well-suited basis. For a function not known in closed form, like the โphysicalโ solution $`\chi ^{(3)}`$, the decomposition on each well-suited local basis associated with every singular point of the ODE, is far from clear. The correspondence between these various (well-suited) local bases associated with each singular point of the ODE, is typically a global problem and, thus, a quite difficult one. One clearly needs to build effective methods to find such connection matrices in the case of Fuchsian differential equations of order seven, or ten ($`\chi ^{(3)}`$ and $`\chi ^{(4)}`$), or of much higher orders ($`\chi ^{(5)}`$, $`\chi ^{(6)}`$, etc.). With a method of matching of series, we will show that the connection matrices matching these various well-suited bases of series-solutions can be obtained explicitly. The entries of these matrices can be calculated with as many digits as we want. We will show that we can actually find the exact expressions of these entries as simple algebraic expressions of (in the case of the Fuchsian ODEโs of $`\chi ^{(3)}`$ and $`\chi ^{(4)}`$) powers of $`\pi `$, $`\mathrm{ln}(2)`$, $`\mathrm{ln}(3)`$ and various algebraic numbers or integers, together with more โtranscendentalโ numbers like the โferromagnetic constantโ $`I_3^+`$ introduced in equation (7.12) of :
$`I_3^+`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _1^{\mathrm{}}}{\displaystyle _1^{\mathrm{}}}{\displaystyle _1^{\mathrm{}}}๐y_1๐y_2๐y_3\left({\displaystyle \frac{y_2^21}{(y_1^21)(y_3^21)}}\right)^{1/2}Y^2`$
$`=`$ $`0.0008144625656625044393912171285627219978\mathrm{}`$
$`Y={\displaystyle \frac{y_1y_3}{(y_1+y_2)(y_2+y_3)(y_1+y_2+y_3)}}`$
Focusing on $`\chi ^{(3)}`$, and since this physical solution is known as a series expansion at $`w=0`$ (low or high temperature expansions), we will give all the connection matrices between this $`w=0`$ regular singular point and all the other regular singularities of the differential equation including the two new complex regular singularities which are roots of $`\mathrm{\hspace{0.17em}1}+3w+4w^2=0`$. We will comment on the occurrence of the โferromagnetic constantโ $`I_3^+`$ in the various blocks of the connection matrices. The decomposition of $`\chi ^{(3)}`$ in the well-suited basis for each regular singular point allows us to find all the singular behavior of the physical solution. From these results, we will deduce the asymptotic behavior of the coefficients of the series expansion of $`\chi ^{(3)}`$. These last problems are interesting, per se, for series expansions analysis of lattice statistical mechanics, since they correspond to subtle analysis of confluent singularities. Actually, we will see that even the last asymptotic evaluation problem is a (global) connection problem since the physical solution like $`\chi ^{(3)}`$ does not correspond to the obvious dominant singular behavior one might have imagined from the indicial equation.
Focusing on the two new singularities, the roots of $`\mathrm{\hspace{0.17em}1}+3w+4w^2=0`$, we will show that the physical solution $`\chi ^{(3)}`$ is not singular at these points. The factor of the logarithmic term, in the decomposition of $`\chi ^{(3)}`$ at these singular points, is known exactly and vanishes identically.
Note that a fundamental concept to understand (the symmetries, the solutions of) these exact Fuchsian differential equations is the so-called differential Galois group . Differential Galois groups have been calculated for simple enough second order, or even third order, ODEโs (see for instance ). However, finding the differential Galois group of such higher order Fuchsian differential equations (order seven for $`\chi ^{(3)}`$, order ten for $`\chi ^{(4)}`$) with eight regular singular points (for $`\chi ^{(3)}`$) is not an easy task and requires the computation of all the monodromy matrices associated with each (non apparent) regular singular point, considered in the same basis<sup>2</sup><sup>2</sup>2These monodromy matrices are the generators of the monodromy group which identifies with the differential Galois group when there are no irregular singularities, and, thus, no Stokes matrices ..
We will give the exact expression of all the monodromy matrices expressed in the same ($`w=0`$) basis of solutions, these eight matrices being the generators of the differential Galois group, which will be given in a forthcoming publication .
This method can be generalized, mutatis mutandis, to the Fuchsian differential equation of $`\chi ^{(4)}`$. Here, we give the connection matrix between $`w=0`$ and, both, the ferromagnetic, and anti-ferromagnetic, critical points. The singular behavior is straightforwardly obtained with the asymptotic behavior of the series coefficients of the physical solution $`\chi ^{(4)}`$. The monodromy matrices, expressed in the same basis of solutions are also obtained.
The paper is organized as follows. We recall, in Section 2, some results on the Fuchsian differential equation satisfied by $`\chi ^{(3)}`$, and give a new factorization for the corresponding order seven differential operator, yielding the emergence of an order two, and an order three, differential operator (denoted $`Z_2`$ and $`Y_3`$ below). We give, in Section 3, the connection matrices matching the (series) solutions around the regular singular point $`w=0`$ and around all the other regular singular points. With these connection matrices we deduce the singularity behavior and the asymptotics on the physical solution of this ODE (Section 4). In Section 5, we deduce the exact expressions of the monodromy matrices expressed in the same basis. Section 6 generalizes these results to the Fuchsian differential equation satisfied by $`\chi ^{(4)}`$. Some physics implications of our results at scaling are discussed in Section 7. Our conclusion is given in Section 8.
## 2 The order seven operator $`L_7`$
Let us first recall, with the same notations as in , the seven linearly independent solutions given in for the order seven differential operator $`L_7`$, associated with<sup>3</sup><sup>3</sup>3$`\stackrel{~}{\chi }^{(n)}`$ is defined as $`\chi ^{(n)}=(1s^4)^{1/4}/s\stackrel{~}{\chi }^{(n)}`$, for $`n`$ odd. $`\stackrel{~}{\chi }^{(3)}`$.
One finds two remarkable rational, and algebraic, solutions of the order seven differential equation associated with $`\stackrel{~}{\chi }^{(3)}`$, namely:
$`๐ฎ(L_1)={\displaystyle \frac{w}{14w}},๐ฎ(N_1)={\displaystyle \frac{w^2}{\left(14w\right)\sqrt{116w^2}}}`$ (3)
associated with the two order $`1`$ differential operators given in :
$`L_1={\displaystyle \frac{d}{dw}}{\displaystyle \frac{1}{w(14w)}},N_1={\displaystyle \frac{d}{dw}}{\displaystyle \frac{2(1+2w)}{w(116w^2)}}`$ (4)
There is a solution behaving like $`w^3`$, that we denote $`S_3`$:
$`S_3=w^3+3w^4+22w^5+74w^6+417w^7+1465w^8`$ (5)
$`+7479w^9+26839w^{10}+\mathrm{}`$
and three solutions with logarithmic terms given by equation (17) in . Note the singled-out series expansion starting with $`w^9`$, corresponding to the physical solution $`\stackrel{~}{\chi }^{(3)}`$:
$`S_9={\displaystyle \frac{\stackrel{~}{\chi }^{(3)}(w)}{8}}=w^9+36w^{11}+4w^{12}+884w^{13}+196w^{14}+\mathrm{}`$ (6)
The choice of this set of linearly independent solutions (and of these series) is, in fact, arbitrary since any linear combination of solutions is also a solution of the differential equation. Three of the above solutions are however singled out: the solutions $`๐ฎ(L_1)`$ and $`๐ฎ(N_1)`$ which are global (since they have closed expression), and the series $`S_9`$ associated with the highest critical exponent in the indicial equation ($`w^9+\mathrm{}`$), which has a unique (well-defined) expression and happens to correspond to the โphysicalโ solution $`\stackrel{~}{\chi }^{(3)}`$. Linear combinations, like $`S_3\alpha S_9`$, are, at first sight, on the same footing.
Nevertheless, introducing such a specific linear combination, B. Nickel<sup>4</sup><sup>4</sup>4We thank B. Nickel for kindly communicating this result. has been able to show that the resulting series for the particular value $`\alpha =16`$ is, also, the solution of a linear differential equation of lower order, namely order four. With this result, the factorization scheme of $`L_7`$ becomes<sup>2</sup><sup>2</sup>2The order four differential operator found by B. Nickel corresponds to $`B_2T_1L_1=B_2O_1N_1=X_1Z_2N_1`$. :
$`L_7`$ $`=`$ $`M_1Y_3Z_2N_1=B_3X_1Z_2N_1`$ (7)
$`=`$ $`B_3B_2O_1N_1=B_3B_2T_1L_1`$
where the indices correspond to the order of the differential operators ($`B_3,Y_3`$ are order three, $`B_2,Z_2`$ order two, โฆ). The differential operators $`L_7`$, $`M_1`$ and $`T_1`$ have been given in . We give in Appendix A, the differential operators $`X_1`$, $`Z_2`$ and $`Y_3`$. With these differential operators, all the factorizations (7) can be found by left and right division.
From these factorizations of $`L_7`$, one can see that the general solution of the corresponding differential equation is the direct sum of the solution of $`L_1`$ and of the general solution of the differential operator $`L_6=Y_3Z_2N_1`$. The operator $`L_7`$ has the following decomposition:
$`L_7=L_6L_1.`$ (8)
We thus consider, from now on, the differential operator $`L_6`$.
The formal solutions of $`L_6`$ (at the singular point $`w=0`$) show the occurrence of three Frobenius series and three solutions carrying logarithmic terms. With the factorizations (7), it is interesting to see which operator brings with it a singular behavior for a given regular singular point. Table 1 shows the critical exponents at each regular singular point for both differential operators $`Z_2N_1`$ and $`Y_3Z_2N_1`$. In the third and sixth column the number of independent solutions with logarithmic terms is shown.
At the singular points $`w=1`$, $`w=1/2`$, and at the two roots $`w_1`$, $`w_2`$ of $`1+3w+4w^2=0`$, we remark that the solution carrying a logarithmic term is in fact a solution of $`Z_2N_1`$. Therefore, the three solutions of the differential operator $`Y_3Z_2N_1`$, emerging from $`Y_3`$, are analytical at the non physical singular points $`w=1`$, $`w=1/2`$, and at the quadratic roots of $`1+3w+4w^2=0`$. At the singular point $`w=1/4`$, we also note that the differential operator $`Z_2N_1`$ is responsible of the $`(14w)^1`$ behavior. We will then expect the โferromagnetic constantโ $`I_3^+`$ to be localized in the blocks of the connection matrix corresponding to the solutions of the order three differential operator $`Z_2N_1`$ at the point $`w=1/4`$.
$`w`$-singularity $`Z_2N_1`$ $`N`$ $`P`$ $`Y_3Z_2N_1`$ $`N`$ $`P`$ $`0`$ $`2,1,1`$ $`1`$ $`1`$ $`3,2,2,1,1,1`$ $`3`$ $`2`$ $`1/4`$ $`1,0,1/2`$ $`0`$ $`0`$ $`2,1,0,0,0,1/2`$ $`2`$ $`2`$ $`1/4`$ $`1,1,3/2`$ $`1`$ $`1`$ $`0,0,0,1,1,3/2`$ $`3`$ $`2`$ $`\mathrm{}`$ $`1,0,0`$ $`1`$ $`1`$ $`2,1,1,1,0,0`$ $`3`$ $`2`$ $`1/2`$ $`3,1,0`$ $`1`$ $`1`$ $`4,3,3,2,1,0`$ $`1`$ $`1`$ $`1`$ $`3,1,0`$ $`1`$ $`1`$ $`4,3,3,2,1,0`$ $`1`$ $`1`$ $`\frac{3\pm i\sqrt{7}}{8}`$ $`1,1,0`$ $`1`$ $`1`$ $`4,3,2,1,1,0`$ $`1`$ $`1`$
Table 1: Critical exponents for each regular singular point for the differential operators $`Z_2N_1`$ and $`Y_3Z_2N_1`$. The columns $`N`$ show the number of solutions with logarithmic terms. The columns $`P`$ show the maximum power of the logarithm occurring in the solutions.
As far as explicit calculations are concerned, a well-suited basis necessary for explicitely writing connection matrices exists and can be described. Considering the order six operator $`L_6=Y_3Z_2N_1`$, we construct the local solutions, sequentially, as the global solution of $`N_1`$ then the two solutions coming from $`Z_2N_1`$, to which we add the three further solutions coming from $`Y_3Z_2N_1`$. We will use below this well-suited bases.
## 3 Connection matrices for $`\stackrel{~}{\chi }^{(3)}`$
Using a very simple method, let us show, in the case where one has an exact Fuchsian differential equation, that one can actually very simply, and very efficiently, obtain the connection matrices between two sets of series-solutions valid at two different points. The method consists in equating, at some matching points, the two sets of series corresponding, respectively, to expansions around $`w=0`$ and, for instance, $`w=1/4`$. The matching point should be in the radius of convergence of both series. The singular points (i.e., $`w=0`$ and $`w=1/4`$) should be neighbors, having no other singularity in between. Recall that the differential equation for $`\stackrel{~}{\chi }^{(3)}`$ has eight regular singular points, the point at infinity, five on the real axis and two ($`w_1`$ and $`w_2`$) on the upper and lower half plane each. At a given singular point $`w_s`$, the solutions are obtained as series in the variable $`x`$, where $`x=w`$ (resp. $`x=1/w`$) for the point $`w_s=0`$ (resp. $`w_s=\mathrm{}`$) and $`x=1w/w_s`$ for the other regular singular points. We take the definition $`\mathrm{ln}(x)=\mathrm{ln}(x)+i\pi `$ for negative values of $`x`$ which corresponds to matching points in the lower (resp. upper) half-plane for $`w>0`$ (resp. $`w<0`$).
The computation of the connection matrix should be more efficient when two โneighboringโ singularities are, as far as possible, far away from the other singularities and, especially, when the test points chosen half-way are, as far as possible, far from the other singularities, in order not to be โpollutedโ by the other singularities. We remark that one can calculate, in this way, just โneighboringโ singularities: connection matrices of two singularities $`w_1`$, $`w_r`$ that are not โneighborsโ should be deduced using some path of โneighboringโ connection matrices:
$`C(w_1,w_r)=C(w_1,w_2)C(w_2,w_3)\mathrm{}C(w_{r1},w_r)`$ (9)
This is the prescription we take for the singular points on the real axis and the singularity $`w_1`$ lying in the upper half-plane. For the singularity $`w_2`$ lying in the lower half-plane, the connection matrix is calculated from:
$`C(0,w_2)=C^{}(0,1/4)C^{}(1/4,w_1)=C^{}(0,w_1)`$ (10)
where denotes the complex conjugate.
Let us remark that changing the variable $`w`$ we are working with, to the more traditional $`s=\mathrm{sinh}(2K)`$ variable, or the usual high-temperature (resp. low temperature) variable $`t=\mathrm{tanh}(K)`$, or the variable $`\tau =(1/ss)/2`$, modifies the distribution of singularities in the complex plane and their radii of convergence. However, the method can still be used. One can use that freedom in the choice of the expansion variable to actually improve the convergence of our calculations.
### 3.1 Connecting solutions
Let us first show, as an example, how we compute the connection matrix between two neighboring regular singular points ($`w=0`$ and $`w=1/4`$) for order three differential operator $`Z_2N_1`$. Around the singular point $`w=0`$, the local solutions are two Frobenius series (one being the global solution $`๐ฎ(N_1)`$) and a series with a logarithmic term. The chosen basis is then (where $`x=w`$):
$`S_1^{(0)}(x)`$ $`=`$ $`๐ฎ(N_1)(x),S_2^{(0)}(x)=[0,1,5,26,106,484,\mathrm{}],`$ (11)
$`S_3^{(0)}(x)`$ $`=`$ $`S_2^{(0)}(x)\mathrm{ln}(x)+S_{30}^{(0)}(x)`$ (12)
with:
$`S_{30}^{(0)}(x)`$ $`=`$ $`[0,0,0,6,26,529/3,2149/3,\mathrm{}]`$ (13)
where $`[a_0,a_1,a_2,\mathrm{},]`$ denotes the series $`a_0+a_1x+a_2x^2+\mathrm{}`$ There are three independent series $`S_1^{(0)}`$, $`S_2^{(0)}`$ and $`S_{30}^{(0)}`$, since the operator $`Z_2N_1`$ is of order three. Similarly, around $`w=1/4`$, the local solutions read (with $`x=14w`$ and, where again, $`S_1^{(1/4)}`$ is the global solution corresponding to operator $`N_1`$):
$`S_1^{(1/4)}(x)`$ $`=`$ $`๐ฎ(N_1)(x),`$ (14)
$`S_2^{(1/4)}(x)`$ $`=`$ $`{\displaystyle \frac{1}{x}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{5}{96}}x{\displaystyle \frac{3}{64}}x^2{\displaystyle \frac{1801}{55296}}x^3+\mathrm{}`$ (15)
$`S_3^{(1/4)}(x)`$ $`=`$ $`S_2^{(1/4)}(x)\mathrm{ln}(x)+S_{30}^{(1/4)}(x)`$ (16)
with:
$`S_{30}^{(1/4)}(x)`$ $`=`$ $`[3/8,367/5760,193/6720,244483/6635520,\mathrm{}]`$ (17)
The series $`S_i^{(0)}`$ are defined around $`w=0`$, and are convergent in a radius of $`1/4`$, which corresponds to the nearest regular singular point (i.e., $`w=1/4`$). Similarly, the solutions $`S_i^{(1/4)}`$ are convergent in the disk centered at $`w=1/4`$ with same radius (i.e., $`1/4`$). Between the points $`w=0`$ and $`w=1/4`$, there is a region where both sets of solutions ($`S_i^{(0)}`$ and $`S_i^{(1/4)}`$) are convergent. This region corresponds to the common area between two disks centered respectively at $`w=0`$, and $`w=1/4`$, with the same radius $`1/4`$.
Connecting the local series-solutions at the regular singular points $`w=0`$, and $`w=1/4`$, amounts to finding the $`3\times 3`$ matrix $`C(0,1/4)`$ such that
$`S^{(0)}=C(0,1/4)S^{(1/4)}`$ (18)
where $`S^{(0)}`$ (resp. $`S^{(1/4)}`$) denotes the vector with entries $`S_i^{(0)}`$ (resp. $`S_i^{(1/4)}`$). The solutions $`S_i^{(0)}`$ and $`S_i^{(1/4)}`$ are evaluated at three arbitrary points around a point $`x_c`$ belonging to both convergence disks of the series-solutions $`S_i^{(0)}`$ and $`S_i^{(1/4)}`$.
Equation (18) is thus a linear system of nine unknowns. The entries of the connection matrix $`C(0,1/4)`$ are obtained in floating point form with a large number of digits. These entries are โrecognizedโ in symbolic form and matrix $`C(0,1/4)`$ then reads:
$`C(0,1/4)=\left[\begin{array}{ccc}1& 0& 0\\ \multicolumn{3}{c}{}\\ 1& \frac{9\sqrt{3}}{64\pi }\left(\frac{2}{3}\mathrm{ln}(24)\right)& \frac{9\sqrt{3}}{64\pi }\\ \multicolumn{3}{c}{}\\ 0& \frac{3\pi \sqrt{3}}{32}& 0\end{array}\right]`$ (22)
The entries of this matrix are combinations of radicals, of powers of $`\pi `$ and logarithms of integers. Note that there is no straightforward manner to recognize numerical values such as the ones displayed above. However, it is possible, in a โtricky wayโ, to get rid of the logarithms of integers in the entries, and obtain as many zero entries as possible. This is shown, in the following, for this very example.
The series, in the set of local solutions $`S_i^{(1/4)}`$, are solutions of the differential equation (ODE) corresponding to the third order differential operator $`Z_2N_1`$ at the regular singular point $`w=1/4`$. It is obvious that any linear combination of these series is also a solution of the differential equation. Consider the following combination instead of the third component in (16):
$`S_3^{(1/4)}(x)(\mathrm{ln}(x/24)+2/3)S_2^{(1/4)}(x)+S_{30}^{(1/4)}(x)`$ (23)
By writing the argument of the logarithm as $`x/24`$, there will be no logarithm in the connection matrix. Furthermore, by adding the second component of the basis to the third component with a factor of $`2/3`$, the entry $`(2,2)`$ of the connection matrix will be canceled. The connection matrix then reads:
$`C(0,1/4)=\left[\begin{array}{ccc}1& 0& 0\\ \multicolumn{3}{c}{}\\ 1& 0& \frac{9}{64}\frac{\sqrt{3}}{\pi }\\ \multicolumn{3}{c}{}\\ 0& \frac{3\pi \sqrt{3}}{32}& 0\end{array}\right]`$ (27)
These tricks, based on well chosen linear combinations of the solutions, allow us to obtain as many zeroes as possible, and to get rid of the logarithms. They will be used in order to compute the connection matrix for $`L_6`$ between the point $`w=0`$ and, respectively, $`w=1/4`$, $`w=1/4`$ and $`w=\mathrm{}`$.
The chosen well-suited basis of solutions, at each regular singular point calls for some comment. The factorization of the differential operator $`L_6`$ being $`Y_3Z_2N_1`$, our method of producing the solutions, sequentially, allows one to determine from which differential operator a given solution emerges. Near the points $`w=0`$, $`w=\pm 1/4`$, and $`w=\mathrm{}`$, the third order differential operator $`Y_3`$ brings three solutions (see Table 1), one Frobenius series, one solution with a $`\mathrm{log}`$ term, and one solution with a $`\mathrm{log}^2`$ term, denoted respectively $`\stackrel{~}{S}_4`$, $`\stackrel{~}{S}_5`$ and $`\stackrel{~}{S}_6`$. The solutions of the differential operator $`Y_3`$ itself are of elliptic integral type (see Appendix B). These elliptic integrals behave around $`w=\pm 1/4`$ (resp. $`w=\mathrm{}`$) like $`g(t)\mathrm{ln}(t/16)+f(t)`$, with $`t=116w^2`$ (resp. $`t=1/16w^2`$), $`g(t)`$ and $`f(t)`$ being series with rational coefficients. One may then assume that the logarithmic term that appears in the solutions of $`L_6`$, inherited from $`Y_3`$, will be of the form $`\mathrm{ln}((116w^2)/16)`$, near $`w=\pm 1/4`$, and of the form $`\mathrm{ln}(1/256/w^2)`$, near $`w=\mathrm{}`$. The general form of combination for the fourth to sixth components of the well-suited basis will be:
$`\stackrel{~}{S}_4`$ $``$ $`\stackrel{~}{S}_4`$
$`\stackrel{~}{S}_5`$ $``$ $`\stackrel{~}{S}_5+\left(a_1\mathrm{ln}(c)\right)\stackrel{~}{S}_4`$ (28)
$`\stackrel{~}{S}_6`$ $``$ $`\stackrel{~}{S}_6+2(a_1\mathrm{ln}(c))\stackrel{~}{S}_5+(\mathrm{ln}(c)^22a_1\mathrm{ln}(c)+a_2)\stackrel{~}{S}_4`$
where $`c=1,\mathrm{\hspace{0.17em}8},\mathrm{\hspace{0.17em}16}`$ for the basis at, respectively, $`w=0`$, $`w=\pm 1/4`$ and $`w=\mathrm{}`$. The values of the parameters $`a_1`$ and $`a_2`$ depend on each basis.
Note that the argument in $`\mathrm{ln}(x/24)`$ in the series solutions of the differential operator $`Z_2N_1`$ at $`w=1/4`$ will be $`\mathrm{ln}(x/4)`$ and $`\mathrm{ln}(x/24)`$ at respectively $`w=\mathrm{}`$ and $`w=1`$. Similarly to $`Y_3`$, these arguments may come from the explicit solutions of $`Z_2`$.
### 3.2 Connection matrix between $`w=0`$ and $`w=1/4`$
The first three local solutions at $`w=0`$ are given by (11), (12), (13), and the fourth, fifth and sixth solutions read
$`S_4^{(0)}(x)`$ $`=`$ $`[0,1,9,34,178,692,\mathrm{}],`$
$`S_5^{(0)}(x)`$ $`=`$ $`S_4^{(0)}(x)\mathrm{ln}(x)+S_{50}^{(0)}(x)S_4^{(0)}(x)/4,`$
$`S_6^{(0)}(x)`$ $`=`$ $`S_4^{(0)}(x)\mathrm{ln}^2(x)+2\left(S_{50}^{(0)}(x)S_4^{(0)}(x)/4\right)\mathrm{ln}(x)`$
$`+S_{60}^{(0)}(x)S_{50}^{(0)}(x)/2+25S_4^{(0)}(x)/16`$
with:
$`S_{50}^{(0)}(x)`$ $`=`$ $`[0,0,0,2,34,241/3,\mathrm{}],`$
$`S_{60}^{(0)}(x)`$ $`=`$ $`[0,0,0,0,19/3,7693/72,575593/1800,\mathrm{}].`$
At the singular point $`w=1/4`$, we make use of the combination (3.1) which amounts to taking $`x/8`$ as argument of the logarithms in the fourth, fifth and sixth component. The parameters $`a_1`$ and $`a_2`$ in (3.1) are respectively $`23/6`$ and $`41/9`$. The first three local series at $`x=14w`$ are given in (14), (15), (17), (23), and the fourth, fifth and sixth read
$`S_4^{(1/4)}(x)`$ $`=`$ $`[1,1/8,3/16,29/512,\mathrm{}],`$ (29)
$`S_5^{(1/4)}(x)`$ $`=`$ $`(\mathrm{ln}(x/8)+23/6)S_4^{(1/4)}(x)+S_{50}^{(1/4)}(x),`$
$`S_6^{(1/4)}(x)`$ $`=`$ $`\left(\mathrm{ln}^2(x/8)+{\displaystyle \frac{23}{3}}\mathrm{ln}(x/8)+41/9\right)S_4^{(1/4)}(x)`$
$`+2\left(\mathrm{ln}(x/8)+23/6\right)S_{50}^{(1/4)}(x)+S_{60}^{(1/4)}(x)`$
with:
$`S_{50}^{(1/4)}(x)`$ $`=`$ $`[0,457/480,2231/1680,128969/184320,\mathrm{}]`$
$`S_{60}^{(1/4)}(x)`$ $`=`$ $`[0,967/100,4312219/470400,595578701/116121600,\mathrm{}]`$
Connecting both solutions amounts to solving a linear system of 36 unknowns (the entries of the connection matrix). We have been able to recognize these entries which are obtained in floating point form with a large number of digits. The connection matrix $`C(0,1/4)`$ for the order six differential operator $`L_6`$ reads:
$`C(0,\mathrm{\hspace{0.17em}1}/4)=`$ (30)
$`\left[\begin{array}{cccccc}1& 0& 0& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 1& 0& \frac{9\sqrt{3}}{64\pi }& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 0& \frac{3\pi \sqrt{3}}{32}& 0& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 5& \frac{1}{3}2I_3^+& \frac{3\sqrt{3}}{64\pi }& 0& 0& \frac{1}{16\pi ^2}\\ & & \multicolumn{4}{c}{}\\ \frac{5}{4}& \frac{3\pi \sqrt{3}}{32}& \frac{45\sqrt{3}}{256\pi }& 0& \frac{1}{32}& 0\\ & & \multicolumn{4}{c}{}\\ \frac{29}{16}\frac{2\pi ^2}{3}& \frac{15\pi \sqrt{3}}{64}& \frac{225\sqrt{3}}{1024\pi }\frac{3\pi \sqrt{3}}{64}& \frac{\pi ^2}{64}& 0& 0\end{array}\right]`$ (37)
Some comments on how these entries have been โrecognizedโ will be given below. Let us remark that, once the entries of the connection matrix have been obtained, a further change of basis can be made to get it as โsimpleโ as possible.
### 3.3 Connection matrices between $`w=0`$ and the other regular singular points
The chosen basis of solutions and the connection matrices between $`w=0`$ (high or low temperature) and, respectively, the anti-ferromagnetic point $`w=1/4`$ and the point $`w=\mathrm{}`$ (corresponding to $`s=\pm i`$) are given in Appendix C.
The chosen basis, used for the regular singular points $`w=1,1/2`$ and $`1+3w+4w^2=0`$, are given in Appendix D together with the corresponding connection matrices with the point $`w=0`$. Many entries are โrecognizedโ and, in particular, those required to find the singular behavior of the physical solution. They correspond to the third column of matrices given in Appendix D.
The connection matrix between each pair of neighboring singular points is computed with the well defined procedure described above. The connection matrix between $`w=0`$, and a non neighbor singular point, is computed using (9). For instance, $`C(0,1)`$ is computed from $`C(0,1/4)`$ and $`C(1/4,1)`$ as $`C(0,1)=C(0,1/4)C(1/4,1)`$ which says that the solutions defined at $`w=1/4`$ connected to the solutions defined at $`w=0`$, are also the solutions that are connected to the solutions defined at $`w=1`$.
To be more confident of this prescription, let us underline that the connection matrices $`C(0,1)`$ and $`C(0,1/2)`$, deduced from (9), will be used below to confirm known dominant singular behavior of $`\stackrel{~}{\chi }^{(3)}`$ and find the subdominant behavior.
### 3.4 Comments and remarks
The connection matrices between $`w=0`$ and the other singular points are structured in blocks. The latter, due to the factorization of the differential operators and to the sequential building of the solutions, are easily recognized. The block $`(1,2,3)\times (1,2,3)`$ is associated with the third order differential operator $`Z_2N_1`$. The block $`(4,5,6)\times (4,5,6)`$ represents the connection between the solutions (at both $`w=0`$ and the other singular points being considered) of $`L_6`$ that are not solutions of $`Z_2N_1`$. The โferromagnetic constantโ $`I_3^+`$ appears in the connection matrix between $`w=0`$ and $`w=1/4`$, as mentioned earlier, in the block $`(1,2,3)\times (1,2,3)`$ at the column corresponding to the $`S_2^{(1/4)}`$ (see (15)) solution of the third order differential operator $`Z_2N_1`$.
To compute the connection matrix, we have used the differential operator $`L_6`$ which has a unique factorization. If, instead, we consider the differential operator $`L_7`$, the next solution (around $`w=0`$), that comes from $`M_1`$, will be the series (6) and will correspond to $`\stackrel{~}{\chi }^{(3)}`$. This seventh solution is expressed as a linear combination of the already existing components and of the solution of the differential operator $`L_1`$. We can then choose to add the latter as the seventh solution. The connection matrix will have a $`1`$ at the entry $`(7,7)`$ and zero elsewhere on the seventh line (and column), since the solution of the differential operator $`L_1`$ is global. By considering another factorization of $`L_7`$, we will get the same structure with an obvious relabelling.
Let us make a few computational remarks on the calculation of these connection matrices. At the matching of the series-solutions for which $`1500`$ coefficients<sup>3</sup><sup>3</sup>3For some checks, $`3000`$ terms have been generated. are generated from homogeneous and non-homogeneous recurrences, the entries of the matrix are computed with $`\mathrm{\hspace{0.17em}800}`$ digits for all the singular points. The numbers that come in floating form are โrecognizedโ as powers of $`\pi `$, radicals and rational numbers, and are in agreement up to $`400`$ digits<sup>4</sup><sup>4</sup>4Let us note that the โferromagnetic constantโ $`I_3^+`$ has been obtained up to more than 400 exact digits. for the connection between the solutions at $`w=0`$ and $`w=\pm 1/4`$, and up to $`100`$ digits for the connection involving other singular points like, $`w=1`$. This fact is related to the convergence rate of the series at the (midway) chosen matching points. For instance, between $`w=0`$ and $`w=1/4`$, the matching points near $`w=1/8`$ are such that both series (at $`w=0`$, and $`w=1/4`$), which have the same radius of convergence, will be faithfully reproduced with the number of terms used in the series. The matching of the solutions between $`w=1/4`$, and $`w=1`$, will then require more terms to fulfill the same accuracy than in the $`(w=0)`$-$`(w=1/4)`$ situation. This is due to the fact that, at $`w=1`$, the convergence radius being $`3/4`$, the matching points, which should be in the common region of both disks, are closer to $`w=1/4`$ than to $`w=1`$. As a general rule, the matching points are chosen around the middle of the segment in the common area between the convergence disks of the two regular singular points for which the connection matrix is computed.
The difficulty in finding โnon-localโ connection matrices is rooted in the recognition of the entries. We have given the connection matrix between $`w=0`$ and $`w=1/4`$ with entries fully recognized (apart from $`I_3^+`$) to show that the method actually works and is efficient. For the matrices concerning the connection between $`w=0`$ and the other singular points, we have concentrated our effort on the entries that will show up in the physical solution. We should note that there is no reason to expect the other (not yet recognized) entries to be โsimplyโ combinations of $`\pi `$โs, $`\mathrm{log}`$โs and radicals. These entries are probably valuations of holonomic functions. This was clearly seen in numerous examples we tackled of various differential equations (of order two and three) with known solutions of hypergeometric type. The recognition process used the fact that we actually found the explicit solutions of differential operator $`Y_3`$ and, thus, knew how the numerical logarithms can be tackled. These were โabsorbedโ in the basis. We know, on the other hand, that the problem is roomed with hypergeometric functions. We then expect some $`\pi `$โs to be present. For the entries consisting of simple product expression, recognizing the number amounts to performing simple arithmetic operations. Note that considering the inverse of the connection matrix, some entries also show up as simple rationals. The combination where $`\pi `$โs, radicals and rationals appear additively comes from looking to, for instance, at the determinant of the matrices, or block matrices, which happen to be easily recognizable (in fact rational or quadratic numbers for the roots of $`1+3w+4w^2=0`$).
Another remark is the following. We first obtained the connection matrix (30) in some general basis. The matrix had more non zero entries compared to (30) involving powers of $`\pi `$, radicals and also $`\mathrm{ln}(3)`$ and powers of $`\mathrm{ln}(2)`$. The well-suited basis we chose has โevacuatedโ all these logโs in the entries of the matrix, lessening the recognition-process effort. But, of course, all these logs will reappear in the final result such as the singular behavior of the physical solution as next sections will show.
## 4 The physical solution $`\stackrel{~}{\chi }^{(3)}`$ and its singular behavior
The calculations of connection matrices are obtained straightforwardly from the well-defined numerical process described in Section 3. Having $`N`$ singularities, one needs $`N1`$ such connection matrices in order to find the correspondence between all these well-suited bases of series-solutions.
Let us focus on some particular entries of these various connection matrices, namely the entries corresponding to the decomposition of $`\stackrel{~}{\chi }^{(3)}`$ in terms of the various well-suited bases associated with each singularity. We have used the fact that the physical solution (corresponding to $`\stackrel{~}{\chi }^{(3)}`$) decomposes as the solution of differential operator $`L_1`$, $`๐ฎ(L_1)`$ (which is $`\stackrel{~}{\chi }^{(1)}/2`$) and the physical solution of the operator $`L_6`$ denoted $`\mathrm{\Phi }_6(w)`$ :
$`\stackrel{~}{\chi }^{(3)}(w)={\displaystyle \frac{1}{6}}\stackrel{~}{\chi }^{(1)}+\mathrm{\Phi }_6(w)`$
Furthermore, our well-suited basis of solutions at the singular point $`w=0`$, does not contain, as a component, the physical solution $`\mathrm{\Phi }_6(w)`$ which is given in terms of the previously considered components as:
$`\mathrm{\Phi }_6(w)={\displaystyle \frac{4}{3}}S_1^{(0)}{\displaystyle \frac{1}{12}}S_2^{(0)}{\displaystyle \frac{1}{4}}S_4^{(0)}`$ (38)
This physical solution can now be easily obtained from the connection matrices between $`w=0`$ and any regular singular point, that we denote $`w=w_s`$ (with $`x=w`$, $`x=1/w`$ for respectively $`w=0`$ and $`w=\mathrm{}`$ and $`x=1w/w_s`$, otherwise) as:
$`\mathrm{\Phi }_6(x)={\displaystyle \underset{j=1}{\overset{6}{}}}\left({\displaystyle \frac{4}{3}}C(0,w_s)_{1j}{\displaystyle \frac{1}{12}}C(0,w_s)_{2j}{\displaystyle \frac{1}{4}}C(0,w_s)_{4j}\right)S_j^{(w_s)}`$
For instance, at the ferromagnetic critical point, this physical solution $`\mathrm{\Phi }_6(x)`$ can easily be deduced from (30), and written as:
$`\mathrm{\Phi }_6(x)={\displaystyle \frac{1}{4}}\left({\displaystyle \frac{1}{3}}2I_3^+\right)S_2^{(1/4)}{\displaystyle \frac{1}{64\pi ^2}}S_6^{(1/4)}`$
$`S_2^{(1/4)}`$ and $`S_6^{(1/4)}`$ are known from their series expansion (15), (29). This equation, giving the full expansion of $`\stackrel{~}{\chi }^{(3)}`$ at $`w=1/4`$, can hardly be obtained directly from the integrals defining $`\stackrel{~}{\chi }^{(3)}(w)`$. One has similar expansions for all the other singular points.
### 4.1 Singular behavior of $`\stackrel{~}{\chi }^{(3)}`$
Knowing the behavior of solutions $`S_j^{(w_s)}`$ near each regular singular point, it is straightforward to get the singular behavior at those points for the physical solution $`\mathrm{\Phi }_6`$ (and thus $`\stackrel{~}{\chi }^{(3)}`$).
Considering the critical behavior of $`\stackrel{~}{\chi }^{(3)}`$ near the ferromagnetic critical point $`w=\mathrm{\hspace{0.17em}1}/4`$, and denoting $`x=\mathrm{\hspace{0.17em}1}4w`$, the singular part of the โphysicalโ solution $`\stackrel{~}{\chi }^{(3)}`$ reads:
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},1/4)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{I_3^+}{x}}{\displaystyle \frac{1}{64\pi ^2}}S_4^{(1/4)}\mathrm{ln}^2(x)`$
$`+{\displaystyle \frac{1}{32\pi ^2}}\left((3\mathrm{ln}(2){\displaystyle \frac{23}{6}})S_4^{(1/4)}S_{50}^{(1/4)}\right)\mathrm{ln}(x)`$
where $`I_3^+`$ is actually the โferromagnetic constantโ (1), and $`S_i^{(1/4)}`$ the series defined in the well-suited basis (29) at $`w=1/4`$. The results agree with previous results of B. Nickel, but the correction terms are new<sup>2</sup><sup>2</sup>2These results have also been found by B. Nickel (private communication)., in particular the term $`\mathrm{\hspace{0.17em}3}\mathrm{ln}(2)/32/\pi ^2`$ in (4.1). In terms of the $`\tau =(1/ss)/2`$ variable introduced in , the singular part (4.1) reads:
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},\tau 0){\displaystyle \frac{I_3^+}{\tau ^2}}{\displaystyle \frac{\mathrm{ln}^2(\tau )}{16\pi ^2}}+\left(\mathrm{ln}(2){\displaystyle \frac{23}{24}}\right){\displaystyle \frac{\mathrm{ln}(\tau )}{4\pi ^2}}+\mathrm{}`$
Near the antiferromagnetic critical point $`w=1/4`$, $`\stackrel{~}{\chi }^{(3)}`$ behaves as:
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},1/4)`$ $`=`$ $`{\displaystyle \frac{1}{32\pi ^2}}S_4^{(1/4)}\mathrm{ln}^2(x)`$
$`{\displaystyle \frac{1}{16\pi ^2}}\left(3(2\mathrm{ln}(2))S_4^{(1/4)}+S_{50}^{(1/4)}\right)\mathrm{ln}(x)`$
At the non-physical singularities $`w=1`$ and $`w=1/2`$ the physical solution behaves, respectively, like:
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},1)={\displaystyle \frac{\sqrt{3}}{27\pi }}S_2^{(1)}\mathrm{ln}(x)`$ (41)
and
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},1/2)={\displaystyle \frac{8\sqrt{3}}{27\pi }}S_2^{(1/2)}\mathrm{ln}(x)`$ (42)
confirming Nickelโs calculations given in .
At the point $`w=\mathrm{}`$, corresponding to the non physical singularities $`s=\pm i`$, the singular behavior reads:
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},\mathrm{})={\displaystyle \frac{1}{16\pi ^2}}S_4^{(\mathrm{})}\mathrm{ln}^2(x)`$ (43)
$`{\displaystyle \frac{1}{8\pi ^2}}\left((42\pi i)S_2^{(\mathrm{})}(5+4\mathrm{ln}(2)+i{\displaystyle \frac{\pi }{2}})S_4^{(\mathrm{})}+S_{50}^{(\mathrm{})}\right)\mathrm{ln}(x)`$
At the new singularities found in , namely the roots of $`\mathrm{\hspace{0.17em}1}+3w+4w^2=0`$, which are regular singular points of the differential equation, the singular part of the physical solution reads, at first sight:
$`\stackrel{~}{\chi }^{(3)}(\mathrm{singular},w_1)={\displaystyle \frac{1}{12}}\left(a_{23}+3a_{43}\right)S_2^{(w_1)}\mathrm{ln}(x)`$
The entries $`a_{23}`$ and $`a_{43}`$ (see the connection matrix for these points in Appendix D) are however such that $`a_{23}+3a_{43}=0`$. The physical solution is thus, not singular, at the newly found quadratic singularities, confirming our conclusion given in from series analysis.
### 4.2 Asymptotic series analysis
As the physical solution $`\stackrel{~}{\chi }^{(3)}`$ is given as a series around $`w=0`$, the coefficients of the latter are controlled by the nearest singular points (i.e. $`w=\pm 1/4`$). Since the singular parts at the ferromagnetic and anti-ferromagnetic critical points (4.1), (4.1) are obtained, it is straightforward to deduce the behavior of the coefficients of series (6) for large values of $`n`$. Standard study of the asymptotic behavior of the coefficients via their linear recursion relation can be used (see ). For our purpose, we use the following identity for $`\mathrm{ln}^2(1x)`$ (where $`x`$ stands for $`x=4w`$):
$`\mathrm{ln}^2(1x)={\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}b(n)x^n,\text{where:}`$
$`b(n)={\displaystyle \underset{i=1}{\overset{n1}{}}}{\displaystyle \frac{1}{i(ni)}}={\displaystyle \frac{2}{n}}\left(\mathrm{\Psi }(n)+\gamma \right)`$ (44)
where $`\gamma =0.57721566\mathrm{}`$ denotes Eulerโs constant, and $`\mathrm{\Psi }`$ denotes the logarithmic derivative of the $`\mathrm{\Gamma }`$ function. Recalling the asymptotic expansion of $`\mathrm{\Psi }(n)`$ up to $`1/n^2`$ for large values of $`n`$, one obtains:
$`b(n){\displaystyle \frac{2}{n}}\left(\gamma +\mathrm{ln}(n){\displaystyle \frac{1}{2n}}{\displaystyle \frac{1}{12n^2}}+\mathrm{}\right)`$
With the same manipulations of $`\mathrm{ln}^2(1+x)`$, and inserting in (4.1), (4.1), one obtains the asymptotic form of coefficients of $`\stackrel{~}{\chi }^{(3)}/8w^9`$ as:
$`2^{15}{\displaystyle \frac{c(n)}{4^n}}`$ $``$ $`{\displaystyle \frac{I_3^+}{2}}{\displaystyle \frac{1}{16\pi ^2}}\left({\displaystyle \frac{1}{2}}+(1)^n\right)\left({\displaystyle \frac{\mathrm{ln}(n)}{n}}+{\displaystyle \frac{b_1}{n}}{\displaystyle \frac{1}{2n^2}}\right)`$
$`+{\displaystyle \frac{1}{16\pi ^2}}\left({\displaystyle \frac{23}{12}}+6(1)^n\right){\displaystyle \frac{1}{n}}+\mathrm{}`$
where $`b_1=\gamma +3\mathrm{ln}(2)`$.
It is this parity effect in the asymptotic behavior of the coefficients that we saw, numerically, (see equations (33) in ) where we obtained, around $`n500`$, $`c(n)\mathrm{\hspace{0.17em}13.5}\times 4^n`$ for $`n`$ even and $`c(n)\mathrm{\hspace{0.17em}11}\times 4^n`$ for $`n`$ odd. For very large values of $`n`$, the asymptotic value of the coefficient $`c(n)/4^n`$ is thus $`2^{14}I_3^+13.34415467\mathrm{}`$.
## 5 Monodromy matrices for $`\stackrel{~}{\chi }^{(3)}`$
### 5.1 Sketching the differential Galois group of $`L_7`$
As a consequence of the direct sum (8), the differential Galois group of $`L_7`$ reduces (up to a product by $`๐`$) to the differential Galois group of $`L_6`$. From the factorization of $`L_6`$, one can immediately deduce that the differential Galois group of $`L_6`$ is the semi-direct product of the differential Galois group of $`Y_3`$, of the differential Galois group of $`Z_2`$ and of the differential Galois group of $`N_1`$ (namely $`๐`$).
In some โwell-suited global basisโ of solutions, the form of the $`\mathrm{\hspace{0.17em}6}\times 6`$ matrices representing the differential Galois group of $`L_6`$, reads:
$`\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right],\mathrm{with}๐=\left[\begin{array}{cc}b& 0\\ \multicolumn{2}{c}{}\\ ๐ก& ๐ \end{array}\right]`$ (49)
where the $`\mathrm{\hspace{0.17em}2}\times 2`$ matrix $`๐ `$, and $`\mathrm{\hspace{0.17em}3}\times 3`$ matrix $`๐`$ correspond, respectively, to the differential Galois group of $`Z_2`$ and $`Y_3`$. The $`\mathrm{\hspace{0.17em}3}\times 3`$ matrix $`๐`$ is associated with the differential Galois group of $`Z_2N_1`$, and the $`\mathrm{\hspace{0.17em}3}\times 3`$ matrix $`๐`$ corresponds to the fact that we have a semi-direct product of the differential Galois group of $`Y_3`$ and $`Z_2N_1`$ in $`L_6=L_3Z_2N_1`$.
Many papers (for instance ) describe how to calculate the differential Galois groups of order 2 and order 3 differential operators. The differential Galois group of $`L_7`$ will be deduced in a forthcoming publication .
To go beyond this sketchy description of the differential Galois group, one needs to calculate specific elements like the monodromy matrices expressed in a common basis.
### 5.2 Monodromy matrices rewritten in the $`w=0`$ basis
Having the connection matrices between $`w=0`$ and each singularity, the local monodromy matrices expressed in their own well-suited basis of (series) solutions, can be rewritten in a unique global basis valid for all singularities. This will allow us, in a second step, to calculate their products and thus generate the differential Galois group. Let us define the $`2\times 2`$ and $`3\times 3`$ matrices
$`A=\left[\begin{array}{cc}1& 0\\ \multicolumn{2}{c}{}\\ \mathrm{\Omega }& 1\end{array}\right],B=\left[\begin{array}{ccc}1& 0& 0\\ \multicolumn{3}{c}{}\\ \mathrm{\Omega }& 1& 0\\ \multicolumn{3}{c}{}\\ \mathrm{\Omega }^2& 2\mathrm{\Omega }& 1\end{array}\right]`$ (55)
where $`\mathrm{\Omega }`$ denotes $`\mathrm{\hspace{0.17em}2}i\pi `$ and corresponds to the translation of the logarithm when performing a complete rotation around the regular singular point: $`\mathrm{ln}(w)\mathrm{ln}(w)+\mathrm{\Omega }`$.
The expression of the local monodromy matrix around each regular singular point $`w_s`$ in its own well-suited basis of (series) solutions reads:
$`l(w_s)=\left[\begin{array}{ccc}ฯต& 0& 0\\ \multicolumn{3}{c}{}\\ 0& C& 0\\ \multicolumn{3}{c}{}\\ 0& 0& D\end{array}\right]`$ (59)
where, $`ฯต`$ and the $`\mathrm{\hspace{0.17em}2}\times 2`$ blocks $`C`$, and $`\mathrm{\hspace{0.17em}3}\times 3`$ blocks $`D`$, are such that:
$`w=0,w=\mathrm{}`$ $``$ $`ฯต=+1,C=A,D=B`$
$`w=1/4,`$ $``$ $`ฯต=1,C=A,D=B`$
$`w=1/4,`$ $``$ $`ฯต=1,C=Id,D=B`$
$`w=1,1/2,3/8\pm i\sqrt{7}/8,`$ $``$ $`ฯต=+1,C=A,D=Id`$
The monodromy matrix around any singularity $`w=w_s`$ expressed in terms of the $`(w=0)`$ well-suited basis, and denoted $`M_{w=0}(w_s)`$, reads:
$`M_{w=0}(w_s)=C(0,w_s)l(w_s)(\mathrm{\Omega })C^1(0,w_s).`$ (60)
In order to keep track of the $`\pi `$ corresponding to the translation of the logarithm in the local monodromy matrix $`l(w_s)(\mathrm{\Omega })`$, and the $`\pi `$โs occurring in the expression of the entries of the (quite involved) connection matrix $`C(0,w_s)`$, we will denote the latter by $`\alpha =\mathrm{\hspace{0.17em}2}i\pi `$.
Let us focus on the singular point $`w=1`$. Its monodromy matrix, expressed in terms of the $`w=0`$ well-suited basis, is given by (60) with $`w_s=1`$, and where the connection matrix $`C(0,1)`$, matching the $`(w=1)`$ well-suited basis together with the $`(w=0)`$ well-suited basis, is a โquite involvedโ matrix given in Appendix D, with entries depending on $`\pi `$โs and on a set of 15 constants, not yet recognized in closed form. The monodromy $`M_{w=0}(1)`$ can finally be written as a function of only $`\alpha `$ and $`\mathrm{\Omega }`$:
$`8\alpha ^2M_{w=0}(1)(\alpha ,\mathrm{\Omega })=`$ (61)
$`\left[\begin{array}{cccccc}8\alpha ^2& 0& 0& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 48\alpha \mathrm{\Omega }& 8\alpha ^2& 48\mathrm{\Omega }& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 0& 0& 8\alpha ^2& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 1008\alpha \mathrm{\Omega }& 0& 1008\mathrm{\Omega }& 8\alpha ^2& 0& 0\\ & & \multicolumn{4}{c}{}\\ 12\alpha \left(5+16\alpha \right)\mathrm{\Omega }& 0& 12\left(5+16\alpha \right)\mathrm{\Omega }& 0& 8\alpha ^2& 0\\ & & \multicolumn{4}{c}{}\\ \alpha \left(75+44\alpha ^2\right)\mathrm{\Omega }& 0& \left(75+44\alpha ^2\right)\mathrm{\Omega }& 0& 0& 8\alpha ^2\end{array}\right]`$ (68)
Let us give one more example corresponding to the new quadratic singularities $`\mathrm{\hspace{0.17em}1}+3w+4w^2=\mathrm{\hspace{0.17em}0}`$. The monodromy matrix around one of the quadratic singularities $`w=w_1`$, expressed in terms of the $`(w=0)`$ well-suited basis, after the conjugation (60), reads:
$`8\alpha ^2M_{w=0}(w_1)(\alpha ,\mathrm{\Omega })=\left[\begin{array}{cc}A& 0\\ \multicolumn{2}{c}{}\\ B& C\end{array}\right]`$ (71)
with:
$`\left[\begin{array}{c}A\\ \\ B\end{array}\right]=\left[\begin{array}{ccc}8\alpha ^2& 0& 0\\ \multicolumn{3}{c}{}\\ 48\alpha \mathrm{\Omega }& 8\alpha \left(\alpha +6\mathrm{\Omega }\right)& 144\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 16\mathrm{\Omega }\alpha ^2& 16\mathrm{\Omega }\alpha ^2& 8\alpha \left(\alpha +6\mathrm{\Omega }\right)\\ \multicolumn{3}{c}{}\\ 16\alpha \mathrm{\Omega }& 16\alpha \mathrm{\Omega }& 48\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 4\alpha \left(4\alpha 15\right)\mathrm{\Omega }& 4\alpha \left(4\alpha 15\right)\mathrm{\Omega }& 12\left(4\alpha 15\right)\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ \alpha a& \alpha a& 3a\end{array}\right]`$ (80)
with $`a=\left(40\alpha +12\alpha ^2+75\right)\mathrm{\Omega }`$ and $`\left[\begin{array}{c}C\end{array}\right]=\mathrm{\hspace{0.17em}\hspace{0.17em}8}\alpha ^2\mathrm{๐๐}(\mathrm{๐}\times \mathrm{๐})`$.
One can actually verify that the monodromy matrix around the other quadratic singularity $`w=w_2`$ ($`w_2`$ is complex conjugate of $`w_1`$), expressed in terms of the $`(w=0)`$-well suited basis, actually identifies with (71) where $`\alpha `$ has been changed into $`\alpha `$.
We have totally similar results for all the other (regular) singularities. The expression of the other monodromy matrices $`M_{w=0}(w_s)`$, around the other (regular) singular points $`w=w_s`$, are displayed in Appendix E.
We saw that the connection matrices depend on $`I_3^+`$ and on โstill not yet recognizedโ (probably transcendantal) numbers, like $`x_{42}`$ and $`y_{41}`$ (for the connection matrix between $`w=0`$ and $`w=\mathrm{}`$). Rewriting a monodromy matrix in a unique (global) basis like the $`w=0`$ basis, amounts to performing conjugation, like (60), of simple (local) monodromy matrices depending only on $`\mathrm{\Omega }`$, by these quite involved connection matrices. As a consequence, one does expect, at first sight, these monodromy matrices, rewritten in the unique $`w=0`$ basis, to be dependent on the still unknown numbers. For instance one certainly expects the monodromy matrix around $`w=1/4`$ (see Appendix E) to be expressed in terms of the transcendental number $`I_3^+`$, or the monodromy matrix (61) to depend on 15 parameters. It is worth noting that all these matrices $`M(w_s)`$, expressed in the same $`(w=0)`$ well-suited basis, turn out to be quite simple matrices where the entries are actually rational expressions, with integer coefficients of $`\alpha `$ and $`\mathrm{\Omega }`$. Section (5.3) gives some hints on why this is so.
The introduction of the two parameters $`\alpha `$ and $`\mathrm{\Omega }`$ is a nice โtrickโ to track the $`\pi `$โs coming from the connection matrices versus the $`\pi `$โs coming from the local monodromy matrices. However, one should keep in mind that $`\alpha `$ is not independent of $`\mathrm{\Omega }`$: the โtrueโ monodromy matrices are such that $`\alpha =\mathrm{\Omega }`$ ($`\mathrm{\Omega }`$ being equal to $`\mathrm{\hspace{0.17em}2}i\pi `$). Let us denote these โtrueโ monodromy matrices by $`M_i`$, $`i=\mathrm{\hspace{0.17em}1},\mathrm{},\mathrm{\hspace{0.17em}8}`$:
$`M_1=M_{w=0}(\mathrm{})(\mathrm{\Omega },\mathrm{\Omega }),M_2=M_{w=0}(1)(\mathrm{\Omega },\mathrm{\Omega }),`$ (81)
$`M_3=M_{w=0}(1/4)(\mathrm{\Omega },\mathrm{\Omega }),M_4=M_{w=0}(w_1)(\mathrm{\Omega },\mathrm{\Omega }),`$
$`M_5=M_{w=0}(1/2)(\mathrm{\Omega },\mathrm{\Omega }),M_6=M_{w=0}(1/4)(\mathrm{\Omega },\mathrm{\Omega }),`$
$`M_7=M_{w=0}(0)(\mathrm{\Omega },\mathrm{\Omega }),M_8=M_{w=0}(w_2)(\mathrm{\Omega },\mathrm{\Omega })`$
The matrices $`M_2`$, $`M_4`$, $`M_5`$, $`M_8`$, and respectively the matrices $`M_1`$ and $`M_7`$, share the same Jordan block form. The Jordan block forms for $`M_3`$ and $`M_6`$ read respectively:
$`\left[\begin{array}{cccccc}1& 0& 0& 0& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 1& 1& 0& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 1& 1& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 0& 1& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 0& 0& 1& 1\\ \multicolumn{6}{c}{}\\ 0& 0& 0& 0& 0& 1\end{array}\right],\left[\begin{array}{cccccc}1& 0& 0& 0& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 1& 1& 0& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 1& 1& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 0& 1& 0& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 0& 0& 1& 0\\ \multicolumn{6}{c}{}\\ 0& 0& 0& 0& 0& 1\end{array}\right]`$ (94)
These matrices $`M_i`$ are the generators of a $`\mathrm{\hspace{0.17em}6}\times 6`$ matrix representation of the differential Galois group of the Fuchsian differential equation corresponding to $`L_6`$. Any element of the differential Galois group is of the form:
$`M_{P(1)}^{n_1}M_{P(2)}^{n_2}M_{P(3)}^{n_3}M_{P(4)}^{n_4}M_{P(5)}^{n_5}M_{P(6)}^{n_6}M_{P(7)}^{n_7}M_{P(8)}^{n_8}`$ (95)
where $`P`$ denotes an arbitrary permutation of eight elements and where the $`n_i`$โs are positive or negative integers. This looks, at first sight, like an infinite discrete group, but the closure of this infinite set of matrices can be quite large continuous groups like semi-direct products of $`SL(2,๐)`$ with $`SL(3,๐)`$, โฆ
Our โglobalโ (800 digits, 1500 terms) calculations yield quite involved exact connection matrices. With such large and involved computer calculations there is always a risk of a subtle mistake or misprint. At this stage, and in order to be โeven more confidentโ in our results, let us recall that the monodromy matrices must satisfy one matrix relation which will be an extremely severe non-trivial check on the validity of these eight matrices $`M_i`$, or more precisely their $`(\alpha ,\mathrm{\Omega })`$ extensions. Actually it is known (see for instance Proposition 2.1.5 in ), that the monodromy group<sup>3</sup><sup>3</sup>3Which identifies in our Fuchsian case to the differential Galois group. of a linear differential equation (with $`r`$ regular singular points) is generated by a set of matrices $`\gamma _1,\gamma _2,\mathrm{},\gamma _r`$ that satisfy $`\gamma _1\gamma _2\mathrm{}\gamma _r=\mathrm{๐๐}`$, where $`\mathrm{๐๐}`$ denotes the identity matrix. The constraint that โsomeโ product of all these matrices should be equal to the identity matrix, looks quite simple, but is, in fact, โunderminedโ by subtleties of complex analysis on how connection matrices between non neighboring singular points should be computed. The fact that the prescription (9,10) has given no contradictory results on the $`\stackrel{~}{\chi }^{(3)}`$ singular behavior may be an argument that our $`M_i`$โs are not โtoo farโ from these โelementaryโ $`\gamma _i`$โs. In other words, one of the products (95) must be equal to the identity matrix for some set of $`n_i`$โs and for some permutation $`P`$. With the particular choice (81) of ordering of the eight singularities, this product, actually reads:
$`M_1M_2M_3M_4M_5M_6M_7M_8=\mathrm{๐๐}`$ (96)
Of course, from this relation, one also has seven other relations deduced by cyclic permutations. It is important to note that these relations (96) are not verified by extensions like (61), (71) depending on two independant parameters $`\alpha `$ and $`\mathrm{\Omega }`$, of the monodromy matrices $`M_i`$. If one imposes relations (96) for the $`(\alpha ,\mathrm{\Omega })`$ extensions of the $`M_i`$โs, one will find that, necessarily, $`\alpha `$ has to be equal to $`\mathrm{\Omega }`$, but (of course<sup>4</sup><sup>4</sup>4A matrix identity like (96) yields a set of polynomial (with integer coefficients) relations on $`\mathrm{\Omega }=\mathrm{\hspace{0.17em}\hspace{0.17em}2}i\pi `$. The number $`\pi `$ being transcendental it is not a solution of a polynomial with integer coefficients. These polynomial relations have, thus, to be polynomial identities valid for any $`\mathrm{\Omega }`$.) one will find that these matrix identities are verified for any value of $`\mathrm{\Omega }`$, not necessarily equal to $`\mathrm{\hspace{0.17em}2}i\pi `$.
### 5.3 Comments
The entries of the connection matrices have been seen to be expressed as various polynomials, or algebraic combinations of power of $`\pi `$, $`\mathrm{ln}(2)`$, $`\mathrm{ln}(N)`$ ($`N`$ integer), algebraic numbers, etc., and more โinvolvedโ transcendental numbers like (1). On the other hand, the monodromy matrices $`M_{w=0}(w_s)`$, expressed in the same $`(w=0)`$ well-suited bases, have entries which are rational expressions with integer coefficients of $`\alpha `$ and $`\mathrm{\Omega }`$. To get some hint as to how this occurs, let us consider, for instance, the regular singular point $`w=1`$. The local monodromy matrix is almost the unity matrix (only one solution with log) with elements:
$`l(1)_{ij}=\delta _{ij}+\mathrm{\Omega }\delta _{i3}\delta _{j2}`$ (97)
The product (60) giving the global monodromy matrix will be given by
$`M_{w=0}(1)_{ij}=\delta _{ij}+\mathrm{\Omega }C(0,1)_{i3}C^1(0,1)_{2j}`$ (98)
where one can see that only the third column of $`C(0,1)`$ and the second row of its inverse will contribute. These entries have been โrecognizedโ (see Appendix D).
Let us assume that there is another solution with a log term (this is not so, see Table 1). An entry (for instance $`l(1)_{65}`$) of the local monodromy matrix changes from zero to $`\mathrm{\Omega }`$. In this case equation (98) becomes:
$`M_{w=0}(1)_{ij}=\delta _{ij}+\mathrm{\Omega }C(0,1)_{i3}C^1(0,1)_{2j}+\mathrm{\Omega }C(0,1)_{i6}C^1(0,1)_{5j}`$
The entries $`C(0,1)_{i6}`$ and $`C^1(0,1)_{5j}`$ will appear in the global monodromy matrix. In fact, changing the entry $`l(1)_{65}`$ from zero to $`\mathrm{\Omega }`$ means that a formal solution will exhibit logโs, and this will correspond to the entries $`C(0,1)_{i6}`$. As a pratical rule, we found that such entries (corresponding to solutions with logโs) can be easily โrecognizedโ in contrast with the entries corresponding to Frobenius series which will be canceled by the zero entries of $`l(1)`$. The entries corresponding to Frobenius series are probably valuations of holonomic functions.
Let us now assume (for the actual situation) that the whole column $`C(0,1)_{i3}`$ has unknown entries. Recalling the fact that the product of the monodromy matrices, expressed in the same basis, should be equal to the identity matrix (this is what we found for our eight matrices $`M_i`$, see (96)), one then expects the โnot yet guessed constantsโ (i.e., the column $`C(0,1)_{i3}`$) to be given by a non linear system of equations. This is indeed what occurs for this example, and we recover that way the entries given for this case in Appendix D.
A last remark is the following. Right now, we have considered all the matrices (connection and therefore monodromy matrices expressed in a unique basis) with respect to the ($`w=0`$) well-suited basis of solutions. This is motivated by the physical solution $`\stackrel{~}{\chi }^{(3)}`$ which is known as series around $`w=0`$. In fact, we can switch to another $`w=\stackrel{~}{w}`$ well-suited basis of solutions. This amounts to considering the connection $`C(\stackrel{~}{w},w_s)=C^1(0,\stackrel{~}{w})C(0,w_s)`$. For instance, we have actually performed the same calculations for the $`(w=1/4)`$ basis of series solutions. We have calculated all the connection matrices from the $`(w=1/4)`$ basis to the other singular point basis series solutions, and deduced the exact expressions of the corresponding monodromy matrices now expressed in the same $`(w=1/4)`$ basis of series solutions. It is worth noting that we get, this time, for the monodromy $`M_{w=1/4}(w_s)`$ around singular point $`w_s`$ and expressed in the $`(w=1/4)`$ basis, a matrix whose entries depend rationally on $`\alpha `$, $`\mathrm{\Omega }`$, but, this time, also (except for the monodromy matrix at $`w=1`$) on the โferromagnetic constantโ $`I_3^+`$. One verifies that the product of these monodromy matrices in the same order as (96), is actually equal to the identity matrix when $`\alpha =\mathrm{\Omega }`$, the matrix identity being valid for any value of $`\alpha =\mathrm{\Omega }`$ (equal or not to $`\mathrm{\hspace{0.17em}2}i\pi `$), and for any value of $`I_3^+`$ (equal or not to its actual value given in (1)).
We have similar results for the monodromy matrices around singular point $`w_s`$, expressed in the $`(w=\mathrm{})`$ basis, but, now, the monodromy matrices $`M_{w=\mathrm{}}(w_s)`$ depend on $`\alpha `$, $`\mathrm{\Omega }`$, and, this time, on the (not yet recognized) constants $`y_{41}`$ and $`x_{42}`$. Again, the product of these monodromy matrices in the same order as (96), is actually equal to the identity matrix when $`\alpha =\mathrm{\Omega }`$, the matricial identity being valid for any value of $`\alpha =\mathrm{\Omega }`$ (equal or not to $`\mathrm{\hspace{0.17em}2}i\pi `$) and for any values of $`y_{41}`$ and $`x_{42}`$ (equal, or not, to their actual values given in Appendix C).
## 6 Mutatis mutandis: Connection matrices and singular behavior for $`\stackrel{~}{\chi }^{(4)}`$
### 6.1 Connection matrices
The Fuchsian differential equation for<sup>2</sup><sup>2</sup>2$`\stackrel{~}{\chi }^{(n)}`$ is defined as $`\chi ^{(n)}=(1s^4)^{1/4}\stackrel{~}{\chi }^{(n)}`$, for $`n`$ even. $`\stackrel{~}{\chi }^{(4)}`$, the four-particle contribution to the susceptibility, is given in . The order ten differential operator $`_{10}`$ associated with this differential equation has 36 (equivalent up to isomorphisms) factorizations (see Appendix F in ). Consider, for instance, two of these factorizations:
$`_{10}`$ $`=`$ $`N_8M_2L_{25}L_{12}L_3L_0`$ (99)
$`=`$ $`M_1L_{24}L_{13}L_{17}L_{11}N_0`$
The notations are the same as those in , the $`M`$ operators are of order four, the $`N`$ and $`L`$ operators are respectively of order two and one. The two factorizations above mean that $`_{10}`$ is a direct sum of an order eight differential operator, $`_8`$ $`=M_2L_{25}L_{12}L_3L_0`$ and of the order two differential operator $`N_0`$ (which, see , has remarkably $`\stackrel{~}{\chi }^{(2)}`$ as solution):
$`_{10}=_8N_0`$ (100)
As was the case for $`\stackrel{~}{\chi }^{(3)}`$, it is thus sufficient to consider the differential operator $`_8`$ for which a general form of $`\mathrm{\hspace{0.17em}8}\times 8`$ matrices, representing $`๐ขal(_8)`$, the differential Galois group of $`_8`$, is deduced:
$`\left[\begin{array}{cc}\multicolumn{2}{c}{}\\ ๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]`$ (103)
$`๐`$, $`๐`$ and $`๐`$ are $`4\times 4`$ matrices, the latter being lower triangular. Recall that $`_8`$ has four known global solutions (see and below).
Similarly to the calculation on $`\stackrel{~}{\chi }^{(3)}`$, we can, for instance, calculate connection matrices associated with the correspondence between the series near $`x=16w^2=\mathrm{\hspace{0.17em}0}`$ (high temperature) with the series near $`x=16w^2=\mathrm{\hspace{0.17em}1}`$ (ferromagnetic and antiferromagnetic critical point), and find how the โphysical solutionโ $`\stackrel{~}{\chi }^{(4)}`$ can be decomposed on the various well-suited bases around each singular point (physical or non-physical) of the order ten Fuchsian differential equation.
We use the factorization (99) to construct the basis of solutions, sequentially, as the four solutions corresponding to the differential operator $`L_{25}L_{12}L_3L_0`$ that we call respectively $`S_1`$, $`S_2`$, $`S_3`$ and $`S_4`$. To these solutions, we add the four solutions coming from $`_8`$ and inherited from the differential operator $`M_2`$, that we call $`S_5`$, $`\mathrm{}`$, $`S_8`$. Here, again, an optimal choice of basis is made in order to have as many zeroes as possible in the connection matrix with as โsimpleโ entries as possible. The basis of solutions at $`x=0`$ and $`x=1`$ (with respectively $`t=x`$ and $`t=1x`$) have similar forms and read:
$`S_1(t)`$ $`=`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}1},S_2(t)=\mathrm{eq}.(33)\mathrm{in}\text{[10]},`$
$`S_3(t)`$ $`=`$ $`\mathrm{eq}.(32)\mathrm{in}\text{[10]},S_4(t)=\mathrm{eq}.(43)\mathrm{in}\text{[10]},`$
$`S_5(t)`$ $`=`$ $`\mathrm{see}\mathrm{below},S_6(t)=S_5(\mathrm{ln}(t/16)+a_1)+S_{60}`$
$`S_7(t)`$ $`=`$ $`(\mathrm{ln}(t/16)^2+2a_1\mathrm{ln}(t/16)+a_2)S_5`$
$`+2S_{60}\left(\mathrm{ln}(t/16)+a_1\right)+S_{70}`$
$`S_8(t)`$ $`=`$ $`(\mathrm{ln}(t/16)^3+3a_1\mathrm{ln}(t/16)^2+3a_2\mathrm{ln}(t/16)+a_3)S_5`$
$`+3(\mathrm{ln}(t/16)^2+2a_1\mathrm{ln}(t/16)+a_2)S_{60}`$
$`+3\left(\mathrm{ln}(t/16)+a_1\right)S_{70}+S_{80}`$
where the constants $`a_1`$, $`a_2`$ and $`a_3`$ and the series read, near $`x=0`$
$`a_1`$ $`=`$ $`79/60,a_2=751/1800,a_3=10619/375,`$
$`S_5^{(0)}(t)`$ $`=`$ $`[0,0,1,45/32,425/256,945/512,\mathrm{}],`$
$`S_{60}^{(0)}(t)`$ $`=`$ $`[0,2/3,0,2353/13440,121619/322560,\mathrm{}],`$
$`S_{70}^{(0)}(t)`$ $`=`$ $`[8,119/45,0,560333/1411200,\mathrm{}],`$
$`S_{80}^{(0)}(t)`$ $`=`$ $`[0,0,0,0,127639044817/85349376000,\mathrm{}]`$
and, near $`x=1`$ :
$`a_1`$ $`=`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}35}/6,a_2=107/9,a_3=1051745657/749700`$
$`S_5^{(1)}(t)`$ $`=`$ $`[1,1/4,7/64,45/256,3385/16384,\mathrm{}],`$
$`S_{60}^{(1)}(t)`$ $`=`$ $`[0,7/120,3809/13440,42401/16120,9271027/18923520,\mathrm{}],`$
$`S_{70}^{(1)}(t)`$ $`=`$ $`[0,1099/75,741847/78400,218499331/101606400,\mathrm{}],`$
$`S_{80}^{(1)}(t)`$ $`=`$ $`[0,0,0,37462660457/592220160,\mathrm{}]`$
The connection matrix between $`x=0`$ and $`x=1`$ comes out as:
$`C(0,1)=\left[\begin{array}{cc}\mathrm{๐}& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]`$ (106)
where $`\mathbf{\hspace{0.17em}1}`$ denotes the $`\mathrm{\hspace{0.17em}4}\times 4`$ identity matrix and $`\mathbf{\hspace{0.17em}0}`$ denotes the $`\mathrm{\hspace{0.17em}4}\times 4`$ zero matrix. The $`\mathrm{\hspace{0.17em}4}\times 4`$ identity matrix corresponds to the fact the four solutions $`S_1`$, $`\mathrm{}`$, $`S_4`$ are global solutions. The two lower $`4\times 4`$ blocks read:
$`๐=\left[\begin{array}{cccc}a_{51}& a_{52}& \frac{5}{2}& a_{54}\\ \multicolumn{4}{c}{}\\ 0& \frac{2}{3}\pi & 0& \frac{1}{32}\\ \multicolumn{4}{c}{}\\ a_{71}& 0& a_{73}& 0\\ \multicolumn{4}{c}{}\\ a_{81}& \pi ^3& a_{83}& a_{84}\end{array}\right],๐=\left[\begin{array}{cccc}0& 0& 0& \frac{1}{2\pi ^3}\\ \multicolumn{4}{c}{}\\ 0& 0& \frac{1}{2\pi }& 0\\ \multicolumn{4}{c}{}\\ 0& \frac{\pi }{2}& 0& 0\\ \multicolumn{4}{c}{}\\ \frac{\pi ^3}{2}& 0& 0& 0\end{array}\right]`$ (115)
with
$`a_{71}={\displaystyle \frac{\pi ^2}{6}}{\displaystyle \frac{2422}{225}},a_{73}={\displaystyle \frac{5\pi ^2}{6}}+{\displaystyle \frac{2422}{225}},a_{84}={\displaystyle \frac{\pi ^2}{32}}{\displaystyle \frac{1211}{600}}`$
The โnot yet recognizedโ entries of this matrix read:
$`a_{51}`$ $``$ $`17.882936774520,a_{52}7.767669067696,`$
$`a_{54}`$ $``$ $`0.530951641617,a_{81}92.773462923758,`$
$`a_{83}`$ $``$ $`77.887072991056`$
Here again, the block structure of the connection matrix relies on the factorization of $`_8`$ and on the โsequentialโ building of the solutions. The block matrix $`๐`$ represents, specifically, the connection between the solutions inherited from $`M_2`$ at both points $`x=0`$ and $`x=1`$. This fourth order differential operator $`M_2`$ in $`_8`$ (corresponding to $`\stackrel{~}{\chi }^{(4)}`$) is structurally very similar (see the remark at end of Appendix B) to operator $`Y_3`$ in $`L_6`$ ($`\stackrel{~}{\chi }^{(3)}`$). Similarly to $`\stackrel{~}{\chi }^{(3)}`$ case, a ferromagnetic (and anti-ferromagnetic) constant (see (119) below) is localized at the fifth line.
We have also computed the connection matrices<sup>3</sup><sup>3</sup>3The matching points are taken in the lower half-plane of the variable $`x`$. (not given here) between the solutions at $`x=0`$ and respectively $`x=4`$ (corresponding to Nickelโs non-physical singularities) and $`x=\mathrm{}`$ (corresponding to the non-physical singularities $`s=\pm i`$). Denoting by $`M_{x=0}(0)`$, $`M_{x=0}(1)`$, $`M_{x=0}(4)`$ and $`M_{x=0}(\mathrm{})`$, the monodromy matrices expressed in the same $`x=0`$ well-suited basis obtained with similar conjugation like (60), one obtains:
$`M_{x=0}(\mathrm{})M_{x=0}(4)M_{x=0}(1)M_{x=0}(0)=\mathrm{๐๐}`$ (116)
This identity is valid irrespective of the still unknown constants.
### 6.2 Singular behavior of $`\stackrel{~}{\chi }^{(4)}`$
The particular physical solution corresponding to $`\stackrel{~}{\chi }^{(4)}=\stackrel{~}{\chi }^{(2)}/3+\mathrm{\Phi }_8`$ (see ) is given, in terms of the basis chosen at the point $`x=0`$, by:
$`\mathrm{\Phi }_8={\displaystyle \frac{1}{384}}\left(5S_1^{(0)}5S_3^{(0)}2S_5^{(0)}\right)`$ (117)
At the ferromagnetic, and anti-ferromagnetic, critical point $`x=1`$, the solution can be deduced from the above connection matrix and reads:
$`\mathrm{\Phi }_8={\displaystyle \frac{1}{384}}\left(2a_{51}5\right)S_1^{(1)}{\displaystyle \frac{a_{52}}{192}}S_2^{(1)}{\displaystyle \frac{a_{54}}{192}}S_4^{(1)}+{\displaystyle \frac{1}{384\pi ^3}}S_8^{(1)}`$
Here again, the above decomposition corresponds to an expansion at the point $`x=1`$ of the triple integral defining $`\stackrel{~}{\chi }^{(4)}`$.
From this solution, the singular part of $`\stackrel{~}{\chi }^{(4)}`$ reads (with $`t=1x`$):
$`\stackrel{~}{\chi }^{(4)}(\mathrm{singular},1)`$ $`=`$ $`{\displaystyle \frac{I_4^{}}{t}}+{\displaystyle \frac{1}{384\pi ^3}}S_5^{(1)}\mathrm{ln}^3(t)`$
$`{\displaystyle \frac{1}{32\pi ^3}}\left((\mathrm{ln}(2){\displaystyle \frac{35}{24}})S_5^{(1)}{\displaystyle \frac{35}{24}}S_{60}^{(1)}\right)\mathrm{ln}^2(t)`$
$`+{\displaystyle \frac{1}{8\pi ^3}}((\mathrm{ln}(2)^2{\displaystyle \frac{35}{12}}\mathrm{ln}(2)+{\displaystyle \frac{107}{144}})S_5^{(1)}`$
$`({\displaystyle \frac{1}{2}}\mathrm{ln}(2){\displaystyle \frac{35}{48}})S_{60}^{(1)}+{\displaystyle \frac{1}{16}}S_{70}^{(1)})\mathrm{ln}(t)`$
$`+{\displaystyle \frac{1}{48\pi }}_2F_1(1/2,1/2;2;t)\mathrm{ln}(t)`$
The constant $`I_4^{}`$ reads, in terms of the โnot yet recognizedโ numbers $`a_{52}`$, $`a_{54}`$:
$`I_4^{}={\displaystyle \frac{1}{36\pi }}+{\displaystyle \frac{a_{52}}{128}}{\displaystyle \frac{a_{54}\pi }{24}}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0.0000254485110658}\mathrm{}`$ (119)
The first term at the right-hand-side of (119) comes from $`\stackrel{~}{\chi }^{(2)}`$, as well as the last term in (6.2).
Similarly, the singular behavior of the physical solution $`\stackrel{~}{\chi }^{(4)}`$ at the other singular points can easily be obtained from the corresponding connection matrices (not given here). At the singular point $`x=4`$, the physical solution behaves like (with $`t=4x`$):
$`\stackrel{~}{\chi }^{(4)}(\mathrm{singular},4)`$ $`=`$ $`{\displaystyle \frac{it^{13/2}}{2^{10}3^25005}}\left(1+{\displaystyle \frac{5}{4}}t+{\displaystyle \frac{261}{272}}t^2+\mathrm{}\right)`$ (120)
confirming the calculations in .
The singular behavior of $`\stackrel{~}{\chi }^{(4)}`$ at the singular point $`x=\mathrm{}`$ reads (with $`t=1/x`$):
$`\stackrel{~}{\chi }^{(4)}(\mathrm{singular},\mathrm{})=20it^{1/2}(A_0+3A_1\mathrm{ln}(t)`$ (121)
$`+3((a_14\mathrm{ln}(2))S_5^{\mathrm{}}+S_{60}^{\mathrm{}})\mathrm{ln}^2(t)+S_5^{\mathrm{}}\mathrm{ln}^3(t))`$
$`+{\displaystyle \frac{(t)^{1/2}}{36\pi }}\left(1+{\displaystyle \frac{3t}{4}}_2F_1(1/2,5/2;2;t)\mathrm{ln}(t){\displaystyle \frac{9\pi t}{16}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}b_nt^n\right)`$
with
$`A_1`$ $`=`$ $`{\displaystyle \frac{2}{5}}(2K1)S_{41}^{\mathrm{}}+\left(16\mathrm{ln}^2(2)8a_1\mathrm{ln}(2)+a_2\right)S_5^{\mathrm{}}`$
$`+`$ $`2(a_14\mathrm{ln}(2))S_{60}^{\mathrm{}}+3S_{70}^{\mathrm{}}`$
$`A_0`$ $`=`$ $`2\pi ^3\left(i_{52}+i{\displaystyle \frac{24}{\pi ^2}}(2K1)\right)S_2^{\mathrm{}}/5`$
$``$ $`\left({\displaystyle \frac{48}{\pi ^2}}(2K1)+i(5+2r_{53})\right)\pi ^3S_3^{\mathrm{}}/5`$
$``$ $`\left(64\mathrm{ln}^3(2)48a_1\mathrm{ln}^2(2)+12a_2\mathrm{ln}(2)a_3\right)S_5^{\mathrm{}}`$
$`+`$ $`{\displaystyle \frac{6}{5}}(2K1)S_{40}^{\mathrm{}}+3\left(16\mathrm{ln}^2(2)8a_1\mathrm{ln}(2)+a_2\right)S_{60}^{\mathrm{}}`$
$`+`$ $`3(a_14\mathrm{ln}(2))S_{70}^{\mathrm{}}+S_{80}^{\mathrm{}}`$
$`b_n={\displaystyle \frac{\mathrm{\Gamma }(n+1/2)\mathrm{\Gamma }(n+5/2)}{\mathrm{\Gamma }(n+2)\mathrm{\Gamma }(n+1)}}\left(\mathrm{\Psi }(n+2)+\mathrm{\Psi }(n+1)\mathrm{\Psi }(n+{\displaystyle \frac{5}{2}})\mathrm{\Psi }(n+{\displaystyle \frac{1}{2}})\right)`$
where $`K=0.915965\mathrm{}`$ is Catalanโs constant and the other parameters, constants and series are: $`a_1=2/5\pi i`$, $`a_2=1\pi ^24\pi i/5`$, $`a_3=6\pi ^2/5+48193/7500+\pi (\pi ^23)i`$, $`i_{52}=0.740250494\mathrm{}`$, $`r_{53}=2.225246651\mathrm{}`$, and
$`S_2^{\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{16t+2t^2}{2(t1)}},S_3^{\mathrm{}}={\displaystyle \frac{312t+8t^2}{8(t1)^{3/2}}}`$
$`S_{40}^{\mathrm{}}`$ $`=`$ $`[2,41/2,313/48,3047/480,\mathrm{}],`$
$`S_{41}^{\mathrm{}}`$ $`=`$ $`[1,25/2,61/8,129/16,\mathrm{}],`$
$`S_5^{\mathrm{}}`$ $`=`$ $`[0,1,7/10,47/64,981/1280,\mathrm{}],`$
$`S_{60}^{\mathrm{}}`$ $`=`$ $`[0,0,161/300,2039/4800,\mathrm{}],`$
$`S_{70}^{\mathrm{}}`$ $`=`$ $`[0,0,1847/18000,2627/36000,\mathrm{}],`$
$`S_{80}^{\mathrm{}}`$ $`=`$ $`[0,0,0,14423879/7200000,\mathrm{}]`$
The last bracket in (121) comes from $`\stackrel{~}{\chi }^{(2)}`$.
Having the singular part of $`\stackrel{~}{\chi }^{(4)}`$ at the ferromagnetic and anti-ferromagnetic critical points, it is straightforward to obtain the asymptotic behavior of the series coefficients. This time, one needs the form of the coefficients in the expansion of $`\mathrm{ln}^3(1x)`$ that we find to be<sup>4</sup><sup>4</sup>4An asymptotic form can be obtained using various packages available at http://algol.inria.fr/libraries/software.html like the command โequivalentโ in gfun , see details in .
$`\mathrm{ln}^3(1x)={\displaystyle \underset{n=3}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{3}{n}}\left(\mathrm{\Psi }(n)+\gamma \right)^2+{\displaystyle \frac{\pi ^2}{2n}}{\displaystyle \frac{3}{n}}\mathrm{\Psi }(1,n)\right)x^n`$ (122)
where $`\mathrm{\Psi }(1,n)`$ is the first derivative of $`\mathrm{\Psi }(n)`$. Expanding $`\mathrm{\Psi }(n)`$ and $`\mathrm{\Psi }(1,n)`$ up to $`1/n^2`$ for large values of $`n`$, one obtains the following asymptotic behavior for the coefficients of the $`\stackrel{~}{\chi }^{(4)}`$ series:
$`c(n)I_4^{}{\displaystyle \frac{\mathrm{ln}^2(n)}{128\pi ^3n}}+{\displaystyle \frac{\mathrm{ln}(n)}{128\pi ^3n^2}}`$
$`{\displaystyle \frac{b_1\mathrm{ln}(n)}{64\pi ^3n}}{\displaystyle \frac{b_2}{2304\pi ^3n}}+{\displaystyle \frac{b_11}{128\pi ^3n^2}}+\mathrm{}`$
where:
$`b_1`$ $`=`$ $`\gamma +4\mathrm{ln}(2){\displaystyle \frac{35}{6}},`$
$`b_2`$ $`=`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}288}\mathrm{ln}^2(2)+144\gamma \mathrm{ln}(2)+18\gamma ^2210\gamma 840\mathrm{ln}(2)+45\pi ^2+214`$
## 7 $`\stackrel{~}{\chi }^{(1)}+\stackrel{~}{\chi }^{(3)}`$ versus $`\stackrel{~}{\chi }`$ at scaling
Thus far we have discussed, in Sections 4 and 6.2 the mathematical aspects of the solutions to the Fuchsian differential equations for $`\stackrel{~}{\chi }^{(3)}`$ and $`\stackrel{~}{\chi }^{(4)}`$. However, the physics implications of the solutions we have obtained call for some remarks near the physical critical points. Taking, as an example, the ferromagnetic singularity for $`\stackrel{~}{\chi }^{(3)}`$, the sum of the first two $`n`$-particle terms behave at $`\tau 0`$ as:
$`\stackrel{~}{\chi }^{(1)}+\stackrel{~}{\chi }^{(3)}`$ $``$ $`{\displaystyle \frac{1+I_3^+}{\tau ^2}}{\displaystyle \frac{\mathrm{ln}^2(\tau )}{16\pi ^2}}+\left(\mathrm{ln}(2){\displaystyle \frac{23}{24}}\right){\displaystyle \frac{\mathrm{ln}(\tau )}{4\pi ^2}}`$
$`+{\displaystyle \frac{11}{48}}+{\displaystyle \frac{3}{8}}I_3^+{\displaystyle \frac{1}{4\pi ^2}}\left(\mathrm{ln}^2(2){\displaystyle \frac{23}{12}}\mathrm{ln}(2)+{\displaystyle \frac{14}{144}}\right)+\mathrm{}`$
The exact susceptibility, as reported in , yields for the normalized susceptibility $`\stackrel{~}{\chi }`$:
$`\stackrel{~}{\chi }={\displaystyle \frac{s}{(1s^4)^{1/4}}}\chi ={\displaystyle \frac{\left(\tau +\sqrt{1+\tau ^2}\right)^{1/2}}{(1+\tau ^2)^{1/8}}}\times `$ (124)
$`\left(c_1\tau ^2F_+(\tau )+{\displaystyle \frac{\tau ^{1/4}}{\sqrt{2}}}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{q=p^2}{\overset{\mathrm{}}{}}}b_+^{(p,q)}\tau ^q\mathrm{ln}^p(\tau )\right)`$
where $`c_1=1.000815260\mathrm{}`$ is given with some 50 digits in . $`F_+(\tau )`$ and $`b_+^{(p,q)}`$ are given in . The constants $`1+I_3^+`$ and $`c_1`$ verify $`1+I_3^++I_5^+=c_1`$ with 9 digits, $`I_5^+`$, corresponding to $`\chi ^{(5)}`$, is the constant given in (and with some 30 digits in ). Thus, and as suggested in , the partial sums of the $`\chi ^{(n)}`$ would converge rapidly to the full $`\chi `$. Furthermore, adding $`\chi ^{(3)}`$ term has resulted in a series expansion that reproduces the first 24 terms of $`\chi `$ to be compared with only eight first terms for $`\chi ^{(1)}`$ series.
However, equation (124) shows a $`\tau ^{1/4}`$ divergence as an overall factor to the logarithmic singularities. This structure, absent in (7), could suggest, in the most pessimistic scenario, that the $`n`$-particle sequence is perhaps useless in understanding scaling corrections and that one should be cautious in accepting the conclusions of studies of higher field derivatives of the susceptibility, based on similar $`n`$-particle representations . The same situation occurs for the low temperature regime when we compare the first two $`n`$-particle terms ($`\stackrel{~}{\chi }^{(2)}`$ and $`\stackrel{~}{\chi }^{(4)}`$) with the full $`\stackrel{~}{\chi }`$ at scaling <sup>1</sup><sup>1</sup>1For the leading amplitude, $`\stackrel{~}{\chi }^{(2)}`$ and $`\stackrel{~}{\chi }^{(4)}`$ give $`1/12\pi +I_4^{}1.0009593\mathrm{}/12\pi `$ which is very close to $`1.0009603\mathrm{}/12\pi `$ for the full $`\stackrel{~}{\chi }`$ . .
This observation raises several profound issues, which we do not address here. One is how the logarithmic terms in the entire sum add up to make the $`\tau ^{1/4}`$ divergence be factored out. If one assumes that the other $`\stackrel{~}{\chi }^{(2n+1)}`$ terms share the same singularity structure as $`\stackrel{~}{\chi }^3`$, in particular the occurrence (in variable $`\tau `$ or $`s`$) of only integer critical exponents at the ferromagnetic critical point, the $`\tau ^{1/4}`$ divergence, as an overall factor, implies the following correspondence :
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{N(n)}{}}}\alpha _{n,m}S_{n,m}(\tau )\mathrm{ln}^m(\tau )\tau ^{1/4}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{q=p^2}{\overset{\mathrm{}}{}}}b_+^{(p,q)}\tau ^q\mathrm{ln}^p(\tau )`$
with $`S_{n,m}(\tau )`$ analytical at $`\tau =0`$ and $`\alpha _{n,m}`$ numerical coefficients. $`N(n)`$ is the maximum power of logarithmic terms occurring in the solution around the ferromagnetic point of the differential equation of $`\stackrel{~}{\chi }^{2n+1}`$. This correspondence requires probably a very particular structure in the successive differential equations. Obtaining the differential equation for $`\stackrel{~}{\chi }^{(5)}`$ (or for $`\stackrel{~}{\chi }^{(6)}`$), and obtaining much larger series for the full susceptibility $`\chi `$, will certainly help to guess such a structure and understand the susceptibility of the two-dimensional Ising model which continues to be a treasure-trove of profound insights into both the mathematics and physics of integrable systems.
Let us note that the phenomenon we have discussed may be more widespread than that observed here. If so, a whole new chapter could be opened on field-theoretical expansions. The challenging problem one faces here is to link linear and non linear descriptions of a physical problem, namely the description in terms of an infinite number of holonomic (linear) expressions for a physical quantity of a non linear nature. Actually the latter is โPainlevรฉ likeโ since its series expansion can be obtained from a program of polynomial growth which uses exclusively a quadratic finite difference double recursion generalizing the Painlevรฉ equations . The difficulty to link holonomic versus non-linear descriptions of physical problems is typically the kind of problems one faces with the Feynman diagram approach of particle physics, but the susceptibility of the Ising model is, obviously, the simplest non trivial example to address such an important issue.
## 8 Conclusion
We have introduced a simple and very efficient method to calculate numerically, with an arbitrary number of digits, the connection matrices between the independent solutions, defined at two singular points, of differential equations of quite high orders. We have considered the order seven, and ten, Fuchsian ODEโs corresponding to the three and four particle contribution to the magnetic susceptibility of the Ising model. The entries of the connection matrix between two regular singular points have been obtained in floating point form and most of them have been recognized, particularly those that show up in the singular behavior of the physical solutions. They are expressed as polynomial, or algebraic, combinations of $`\pi `$, $`\mathrm{ln}(2)`$, $`\mathrm{}`$, radicals, and more involved numbers (not yet recognized) such as the โferromagnetic constantโ (1). The method allows us to obtain the series expansions of the physical solutions $`\stackrel{~}{\chi }^{(3)}`$ (and $`\stackrel{~}{\chi }^{(4)}`$) around any other regular singular point, besides the already known series around $`w=0`$. We obtained, in this way, near each singular point all the dominant, and subdominant, singular behaviors of the physical solutions. Such subdominant singular behavior is certainly hard to obtain from series analysis. At the newly found quadratic singularities of the differential equation, we showed that the physical solution $`\stackrel{~}{\chi }^{(3)}`$ itself is not singular. Also note, at $`w=1/4`$, that the behavior in $`(14w)^{3/2}`$ corresponding to the largest critical exponent for the ODE is actually absent in the physical solution. Note the remarkable fact that the factorization of differential operator $`L_7`$ (and $`_{10}`$) associated with $`\stackrel{~}{\chi }^{(3)}`$ (respectively $`\stackrel{~}{\chi }^{(4)}`$) shows clearly the differential operator responsible of the non-physical singularities given in and the newly found quadratic numbers . In both cases ($`\stackrel{~}{\chi }^{(3)}`$ and $`\stackrel{~}{\chi }^{(4)}`$), these non-physical singularities are carried by the differential operator $`Z_2N_1`$ (respectively $`L_{25}L_{12}L_3L_0`$) occurring at the right of $`L_7`$ (respectively $`_{10}`$).
The physical solutions $`\stackrel{~}{\chi }^{(3)}`$ (and $`\stackrel{~}{\chi }^{(4)}`$) being known as series around $`w=0`$, the growth behavior of the corresponding series coefficients should be controlled by the singular behavior at the nearest singular points which are the ferromagnetic and anti-ferromagnetic critical points in both cases ($`w=\pm 1/4`$ and $`x=1`$). This growth is easily found from the expansion around the ferromagnetic and anti-ferromagnetic points.
The connection matrices we have obtained allow us to relate the solutions around any given singular point to a common (non-local) basis of solutions. In this respect, we have obtained the exact expression of all the monodromy matrices, expressed in the same basis, and we have seen that they are simple matrices with rational function entries. In a forthcoming publication , we will give the whole structure of the differential Galois group for the two previous Fuchsian differential equations.
As far as the physics implications of the solutions are concerned, we have compared the corrections to scaling at the ferromagnetic point given by the first two terms ($`\chi ^{(1)}`$ and $`\chi ^{(3)}`$) with the full $`\chi `$. Qualitative difference is found raising profound issues on the $`n`$-particle representation of the susceptibility. The same observation occurs for the antiferromagnetic point, and also for the low temperature regime.
Acknowledgments We thank Jacques-Arthur Weil for many valuable comments on differential Galois group and connection matrices. We would like to thank B. Nickel for his inspired comment on solution $`S_3`$, many exchanges of informations and for pointing some misprints in our series-solutions, in the earlier version of the manuscript. We would like to thank A. J. Guttmann, I. Jensen, and W. Orrick for a set of useful comments on the singularity behavior of physical solutions. One of us (JMM) would like to thank B.M. McCoy for many extensive discussions on the problem of the holonomic description of non-linear problems. We would like to thank an anonymous referee for raising important points on the physics implications of our results that we discussed in Section 7. (S. B) and (S. H) acknowledge partial support from PNR3.
## 9 Note added in the Proofs
After completion of the revised version of our manuscript we were told that, as consequence of the work of B. M. McCoy, C. A. Tracy and T.T. Wu, the two transcendental numbers $`I_3^+`$ and $`I_4^{}`$ can actually be written in terms of polylogarithms, namely the Clausen function $`Cl_2`$ and of the Riemann zeta function, as follows :
$`I_3^+={\displaystyle \frac{1}{2\pi ^2}}\left({\displaystyle \frac{\pi ^2}{3}}+23\sqrt{3}Cl_2({\displaystyle \frac{\pi }{3}})\right),Cl_2(\theta )={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}(n\theta )}{n^2}}`$
$`I_4^{}={\displaystyle \frac{1}{16\pi ^3}}\left({\displaystyle \frac{4\pi ^2}{9}}{\displaystyle \frac{1}{6}}{\displaystyle \frac{7}{2}}\zeta (3)\right)`$
The derivation of these results has never been published but these results appeared in a conference proceedings . We have actually checked that $`I_3^+`$ and $`I_4^{}`$ we got from the calculations displayed in our paper as floating numbers with respectively 421 digits and 431 digits accuracy are actually in agreement with the previous two formula. These two results provide a clear answer to the question of how โcomplicated and transcendentalโ some of our constants occurring in the entries of the connection matrices can be. These extremely interesting results are not totally surprising when one recalls the deep link between zeta functions, polylogarithms and hypergeometric series .
## 10 Appendix A
We give, in this Appendix, the explicit expressions of the differential operators $`X_1`$ and $`Z_2`$ and $`Y_3`$. The order one differential operator reads
$`X_1={\displaystyle \frac{d}{dw}}+{\displaystyle \frac{p_0}{p_1}}`$ (125)
with:
$`p_1`$ $`=`$ $`\left(1+w\right)\left(4w1\right)\left(1+2w\right)\left(4w+1\right)\left(1+3w+4w^2\right)`$
$`\left(13w18w^2+104w^3+96w^4\right)`$
$`\left(17w4w^247w^3+36w^4+280w^5+160w^6+256w^7\right)`$
$`p_0`$ $`=`$ $`w(58+909w+3284w^224711w^372352w^4+181016w^5`$
$`+1251768w^6+2852880w^7+1454592w^811455616w^9`$
$`31712256w^{10}20418560w^{11}+20840448w^{12}+34963456w^{13}`$
$`+30146560w^{14}+15728640w^{15})`$
The order two differential operator $`Z_2`$ is
$`Z_2={\displaystyle \frac{1}{p_2}}{\displaystyle \underset{n=0}{\overset{2}{}}}p_n{\displaystyle \frac{d}{dw^n}}`$ (126)
where the polynomials $`p_i`$โs, now, read:
$`p_2`$ $`=`$ $`w\left(4w1\right)^2\left(4w+1\right)\left(1+3w+4w^2\right)\left(1+w\right)\left(1+2w\right)`$
$`\left(13w18w^2+104w^3+96w^4\right)`$
$`p_1`$ $`=`$ $`(4w1)(16w111w^2108w^3+1080w^44488w^5`$
$`40368w^694272w^748384w^8+72704w^9+49152w^{10})`$
$`p_0`$ $`=`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}4}+48w276w^21520w^33192w^44224w^571552w^6`$
$`307200w^7239616w^8+98304w^9+98304w^{10}`$
The order three differential operator $`Y_3`$ is given by
$`Y_3={\displaystyle \frac{1}{p_3}}{\displaystyle \underset{n=0}{\overset{3}{}}}p_n{\displaystyle \frac{d^n}{dw^n}}`$ (127)
where the polynomials $`p_i`$โs, now, read:
$`p_3=w^2\left(w1\right)\left(1+2w\right)\left(1+3w+4w^2\right)`$ (128)
$`\left(4w1\right)^3\left(4w+1\right)^3\left(96w^4+104w^318w^23w+1\right)^3`$
$`(1+19w368w^23296w^3+17882w^4+272599w^5+160900w^6`$
$`6979208w^7+7550800w^8+203094872w^9278920192w^{10}`$
$`3959814304w^{11}2115447424w^{12}+20894729472w^{13}`$
$`+39719728128w^{14}+20516098048w^{15}+256763363328w^{16}`$
$`327065010176w^{17}8810227761152w^{18}+414933057536w^{19}`$
$`+116411936538624w^{20}+296827723186176w^{21}+317648030138368w^{22}`$
$`+179148186189824w^{23}+194933533179904w^{24}+112931870081024w^{25}`$
$`55246164328448w^{26}+11063835754496w^{27}+1511828488192w^{28})`$
$`p_2=w\left(4w1\right)^2\left(4w+1\right)^2\left(96w^4+104w^318w^23w+1\right)^2`$
$`(6+102w2018w^223962w^3+242904w^4+2575633w^5`$
$`12389010w^6178413527w^7+80727412w^8+6252221348w^9`$
$`+2456938016w^{10}178278888104w^{11}103902989696w^{12}`$
$`+3814815965856w^{13}+1524977514176w^{14}67400886678400w^{15}`$
$`74115827788032w^{16}+797710351468032w^{17}+2324376661856256w^{18}`$
$`1561280104050688w^{19}16314064973299712w^{20}`$
$`27005775986622464w^{21}40259640226480128w^{22}`$
$`+35764751009841152w^{23}+1007304244270727168w^{24}`$
$`+1460771505523654656w^{25}13359756413056843776w^{26}`$
$`63988213537189134336w^{27}116684614339309600768w^{28}`$
$`75710498024932245504w^{29}+57121462326803824640w^{30}`$
$`+132479693600191414272w^{31}+111232702128767107072w^{32}`$
$`+106152703871500156928w^{33}+83508376521540632576w^{34}`$
$`+10084606300752183296w^{35}9404395631251816448w^{36}`$
$`+2682738003029262336w^{37}+297237575406452736w^{38})`$
$`p_1=\mathrm{\hspace{0.17em}2}\left(4w1\right)\left(4w+1\right)\left(96w^4+104w^318w^23w+1\right)`$
$`(325w+1013w^2+7893w^3353904w^41562671w^5`$
$`+43285825w^6+192457911w^72690351207w^815077420736w^9`$
$`+94510776436w^{10}+707838800508w^{11}2327528107216w^{12}`$
$`23421365465744w^{13}+45755890012000w^{14}+568028144875200w^{15}`$
$`824814656530816w^{16}10390722028797440w^{17}`$
$`+12438134957505536w^{18}+145637031330319360w^{19}`$
$`127616737495506944w^{20}1708173874007113728w^{21}`$
$`52355400373420032w^{22}+15741676181476802560w^{23}`$
$`+24085046332129804288w^{24}57977682482294161408w^{25}`$
$`168033877030234750976w^{26}56941336876602621952w^{27}`$
$`426707803148891717632w^{28}200805832817071095808w^{29}`$
$`+8716841486700848873472w^{30}6642009916749838811136w^{31}`$
$`192590979400145399971840w^{32}564260086660360537374720w^{33}`$
$`585770764250229243904000w^{34}+235172208485444226121728w^{35}`$
$`+1203159617695281059987456w^{36}+1323272087085206269329408w^{37}`$
$`+997072075164663150542848w^{38}+789138181323007857786880w^{39}`$
$`+388137877034203055390720w^{40}+4946627729914186432512w^{41}`$
$`26947297377570617556992w^{42}+10614515947351012540416w^{43}`$
$`+998718253365681192960w^{44})`$
$`p_0=\mathrm{\hspace{0.17em}2}w(348+2768w+248784w^2358217w^350461860w^4`$
$`+16394998w^5+5283255372w^6+3911764831w^7329364073508w^8`$
$`572985025996w^9+13847002317264w^{10}+38091073842520w^{11}`$
$`437846238222272w^{12}1682624909395232w^{13}`$
$`+10892230218721408w^{14}+52959188332189824w^{15}`$
$`214291413015639808w^{16}1200734422407578112w^{17}`$
$`+3319489124092462080w^{18}+20066023020568346624w^{19}`$
$`38248948302383529984w^{20}254480826931185762304w^{21}`$
$`+261281404771497082880w^{22}+2480194764802183397376w^{23}`$
$`+148352203759030894592w^{24}19049822668612433870848w^{25}`$
$`29328532357149024583680w^{26}+103410036785394615320576w^{27}`$
$`+391034390334579595542528w^{28}+11096790708133489016832w^{29}`$
$`1530120948962096058466304w^{30}2868669407093825701150720w^{31}`$
$`6126661019209831555268608w^{32}+2808943911875675603075072w^{33}`$
$`+40458568379798955017371648w^{34}169712327643359793079386112w^{35}`$
$`1092943871171162347998806016w^{36}1781375524629107822238367744w^{37}`$
$`+250471471742289487729786880w^{38}+4679788548889591917580386304w^{39}`$
$`+7101176295364126941625974784w^{40}+5918768536906007398653624320w^{41}`$
$`+4083406571846803705271681024w^{42}+2567747434748530216944009216w^{43}`$
$`+846246487598480459424595968w^{44}49595159800068478383161344w^{45}`$
$`37040268890013610134208512w^{46}+21784239691989525951676416w^{47}`$
$`+1753178556765355785584640w^{48})`$
## 11 Appendix B: Solutions of the differential operator $`Y_3`$
Considering the critical exponents at the regular singular points, as well as the formal solutions of differential operator $`Y_3`$, one can make the following remarks. The roots of the polynomial of degree 28 in polynomial $`p_3`$ (see (128)) are apparent singularities. The roots of the polynomial of degree four in one of the factors of the same polynomial $`p_3`$ are not apparent singularities. While the formal solutions near $`w=0`$, $`w=\pm 1/4`$, and $`w=\mathrm{}`$, have one Frobenius solution and two logarithmic solutions, the formal solutions near the other regular singular points are free of logarithmic solutions. The critical exponents at $`w=1`$, $`w=1/2`$, roots of $`1+3w+4w^2=0`$, and roots of $`13w18w^2+104w^3+96w^4=0`$, are respectively $`(1,0,1)`$, $`(1,0,1)`$, $`(1,0,1)`$ and $`(1,1,2)`$. This leads us to look for the solutions of the third order differential operator $`Y_3`$ as a linear combination of powers of elliptic integrals with a common factor โtaking careโ of the non logarithmic singularity behavior of the singular points.
Defining
$`K(x)=_2F_1(1/2,1/2;1;x),E(x)=_2F_1(1/2,1/2;1;x)`$
and
$`s(w)`$ $`=`$ $`w^2\left(116w^2\right)^3\left(1+2w\right)\left(1w\right)\left(1+3w+4w^2\right)`$
$`\left(13w18w^2+104w^3+96w^4\right)`$
one obtains the three independent solutions of the differential operator $`Y_3`$ as:
$`S_1(Y_3)`$ $`=`$ $`{\displaystyle \frac{1}{s(w)}}(P_1K^2(16w^2)+P_2E^2(16w^2)`$
$`+P_3K(16w^2)E(16w^2))`$
$`S_2(Y_3)`$ $`=`$ $`{\displaystyle \frac{1}{s(w)}}(P_4K^2(1/16w^2)16w^2P_2E^2(1/16w^2)`$
$`+P_5K(1/16w^2)E(1/16w^2))`$
$`S_3(Y_3)`$ $`=`$ $`{\displaystyle \frac{1}{s(w)}}((P_1+P_2+P_3)K^2(116w^2)+P_2E^2(116w^2)`$
$`(2P_2+P_3)K(116w^2)E(116w^2))`$
with
$`P_4`$ $`=`$ $`{\displaystyle \frac{P_1}{16w^2}}{\displaystyle \frac{(116w^2)^2}{16w^2}}P_2{\displaystyle \frac{116w^2}{16w^2}}P_3,`$
$`P_5`$ $`=`$ $`2(116w^2)P_2P_3`$
where the three polynomials $`P_1`$, $`P_2`$ and $`P_3`$ read:
$`P_1`$ $`=`$ $`(1+4w)(15w69w^2+537w^3+2964w^44100w^5`$
$`46816w^674688w^7+230656w^8+647680w^9+475136w^{10}`$
$`8192w^{11}+720896w^{12})`$
$`P_2`$ $`=`$ $`1+5w+25w^29w^32408w^417460w^519696w^6`$
$`+28800w^73328w^862464w^936864w^{10}`$
$`P_3`$ $`=`$ $`2(13w65w^2+143w^3+3888w^4+15144w^510624w^6`$
$`172416w^7241536w^8+111616w^9+282624w^{10}`$
$`+180224w^{11}+98304w^{12})`$
Remark: Let us note the very close similarity between the differential operator $`Y_3`$, occurring at the left of differential operator $`L_6`$ (see (7)) for $`\stackrel{~}{\chi }^{(3)}`$, and the differential operator $`M_2`$ (see (99)) occurring at the left of differential operator $`_8`$ for $`\stackrel{~}{\chi }^{(4)}`$. For this order four differential operator $`M_2`$, we have been able, using the same ansatz, to obtain in closed form three of the four solutions, also expressed as a linear combination of products of elliptic integrals. Note that, setting $`\lambda =16w^2`$, one can detect in the solutions of $`Y_3`$ (and also in the three solutions of $`M_2`$ we have found) the structure of $`\mathrm{\Sigma }_3`$ permutation group , $`\lambda `$, $`1/\lambda `$, $`1\lambda `$, $`\mathrm{\hspace{0.17em}1}1/\lambda `$, etc.
## 12 Appendix C: Connection matrices between $`w=0`$ and $`w=1/4`$, $`w=\mathrm{}`$
### 12.1 Connection matrix between $`w=0`$ and $`w=1/4`$
The basis of solutions at the anti-ferromagnetic critical point $`w=1/4`$ are chosen as follows (with $`x=1+4w`$)
$`S_1^{(1/4)}(x)`$ $`=`$ $`๐ฎ(N_1)(x),`$
$`S_2^{(1/4)}(x)`$ $`=`$ $`[1,0,1/10,87/700,313/1680,\mathrm{}],`$
$`S_3^{(1/4)}(x)`$ $`=`$ $`[0,1,17/10,23/25,1/30,\mathrm{}],`$
$`S_4^{(1/4)}(x)`$ $`=`$ $`[1,5/2,3/8,5/16,83/512,\mathrm{}],`$
$`S_5^{(1/4)}(x)`$ $`=`$ $`S_4^{(1/4)}(x)\left(\mathrm{ln}(x/8)+6\right)+S_{50}^{(1/4)}(x)`$
$`S_6^{(1/4)}(x)`$ $`=`$ $`S_4^{(1/4)}(x)\left(\mathrm{ln}^2(x/8)+12\mathrm{ln}(x/8)+23264/315\right)`$
$`+2S_{50}^{(1/4)}(x)\left(\mathrm{ln}(x/8)+6\right)+S_{60}^{(1/4)}(x)`$
with:
$`S_{50}^{(1/4)}(x)`$ $`=`$ $`[0,97/6,553/240,2339/672,1678457/645120,\mathrm{}],`$
$`S_{60}^{(1/4)}(x)`$ $`=`$ $`[0,0,0,85997/18000,8450503/1814400,\mathrm{}].`$
Here again, an optimal choice of the components is made in order to remove logarithms and have as many zeroes as possible in the entries of the matrix. The same method of matching the series-solutions at a half-way point between $`w=0`$ and $`w=1/4`$, gives
$`C(0,1/4)=`$ (129)
$`\left[\begin{array}{cccccc}1& 0& 0& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 2& r_{22}& r_{23}& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 2\pi i& r_{32}+r_{22}\pi i& r_{33}+r_{23}\pi i& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 6& \frac{1}{\pi }i_{52}& \frac{1}{\pi }i_{53}& 0& 0& \frac{1}{8\pi ^2}\\ & & \multicolumn{4}{c}{}\\ \frac{5}{2}+6\pi i& a_{52}& a_{53}& 0& \frac{1}{16}& \frac{1}{8\pi }i\\ & & \multicolumn{4}{c}{}\\ \frac{23}{8}\frac{17\pi ^2}{3}+5\pi i& a_{62}& a_{63}& \frac{\pi ^2}{32}& \frac{\pi }{8}i& \frac{1}{8}\end{array}\right]`$ (136)
with:
$`r_{22}r_{33}r_{23}r_{32}=\mathrm{\hspace{0.17em}25}/12288`$
$`a_{52}=3r_{32}{\displaystyle \frac{5}{4}}r_{22}+i_{52}i,a_{53}=3r_{33}{\displaystyle \frac{5}{4}}r_{23}+i_{53}i`$
$`a_{62}=\left({\displaystyle \frac{25}{16}}{\displaystyle \frac{2\pi ^2}{3}}{\displaystyle \frac{5\pi }{2}}i\right)r_{22}\left({\displaystyle \frac{5}{2}}+6\pi i\right)r_{32}i_{52}\pi `$
$`a_{63}=\left({\displaystyle \frac{25}{16}}{\displaystyle \frac{2\pi ^2}{3}}{\displaystyle \frac{5\pi }{2}}i\right)r_{23}\left({\displaystyle \frac{5}{2}}+6\pi i\right)r_{33}i_{53}\pi `$
and where<sup>2</sup><sup>2</sup>2The numbers $`r_{ij}`$โs and $`i_{ij}`$โs are peculiar to each connection matrix.:
$`r_{22}0.059050961331,r_{23}0.018643190255,`$
$`r_{32}0.1631382423131,i_{52}1.839621665835,`$
$`i_{53}0.015467563102`$
### 12.2 Connection matrix between $`w=0`$ and $`w=\mathrm{}`$
The basis of solutions at the singular point $`w=\mathrm{}`$ are chosen as follows (with $`x=1/w`$):
$`S_1^{(\mathrm{})}(x)`$ $`=`$ $`๐ฎ(N_1),`$
$`S_2^{(\mathrm{})}(x)`$ $`=`$ $`[1,1,7/16,1/16,7/256,\mathrm{}],`$
$`S_3^{(\mathrm{})}(x)`$ $`=`$ $`\left(\mathrm{ln}(x/4)2/3\right)S_2^{(\mathrm{})}(x)+S_{30}^{(\mathrm{})}(x),`$
$`S_4^{(\mathrm{})}(x)`$ $`=`$ $`[0,1,0,1/32,9/512,\mathrm{}],`$ (137)
$`S_5^{(\mathrm{})}(x)`$ $`=`$ $`\left(\mathrm{ln}(x/16)+a_1\right)S_4^{(\mathrm{})}(x)+S_{50}^{(\mathrm{})}(x),`$
$`S_6^{(\mathrm{})}(x)`$ $`=`$ $`\left(\mathrm{ln}^2(x/16)+2a_1\mathrm{ln}(x/16)+a_2\right)S_4^{(\mathrm{})}(x)`$
$`+2\left(\mathrm{ln}(x/16)+a_1\right)S_{50}^{(\mathrm{})}(x)+S_{60}^{(\mathrm{})}(x)`$
with:
$`a_1=5{\displaystyle \frac{\pi }{2}}i,a_2={\displaystyle \frac{\pi ^2}{4}}+{\displaystyle \frac{379}{11}}+5\pi i,`$
$`S_{30}^{(\mathrm{})}(x)`$ $`=`$ $`[2/3,1/6,1/24,1/96,7/768,\mathrm{}],`$
$`S_{50}^{(\mathrm{})}(x)`$ $`=`$ $`[0,0,3/2,3/64,107/512,23113/491520,\mathrm{}],`$
$`S_{60}^{(\mathrm{})}(x)`$ $`=`$ $`[0,0,0,93/44,80891/13516,105811/4055040,\mathrm{}].`$
The connection matrix reads
$`C(0,\mathrm{})=`$ (138)
$`\left[\begin{array}{cccccc}1& 0& 0& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 1& \frac{1}{16}& \frac{3}{16\pi }i& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ \pi i& 0& \frac{1}{16}& 0& 0& 0\\ & & \multicolumn{4}{c}{}\\ 11+y_{41}i& x_{42}\frac{1}{\pi }i& \frac{2}{\pi ^2}\frac{15}{16\pi }i& 0& 0& \frac{1}{4\pi ^2}\\ & & \multicolumn{4}{c}{}\\ a_{51}& a_{52}& \frac{9}{16}\frac{49}{64\pi }i& 0& \frac{1}{16}& \frac{1}{8\pi }i\\ & & \multicolumn{4}{c}{}\\ a_{61}& a_{62}& \frac{11}{32}+\frac{5\pi }{16}i\frac{75}{256\pi }i& \frac{\pi ^2}{64}& \frac{\pi }{16}i& \frac{1}{16}\end{array}\right]`$ (145)
where:
$`x_{42}1.534248223197,y_{41}22.932479960454,`$
$`a_{51}={\displaystyle \frac{5}{4}}+{\displaystyle \frac{\pi }{2}}y_{41}+7\pi i,a_{52}={\displaystyle \frac{11}{64}}{\displaystyle \frac{\pi }{2}}x_{42}i{\displaystyle \frac{\pi }{32}}i,`$
$`a_{61}={\displaystyle \frac{29}{16}}+{\displaystyle \frac{16\pi ^2}{3}}{\displaystyle \frac{\pi ^2}{4}}iy_{41}+{\displaystyle \frac{5\pi }{2}}i,a_{62}={\displaystyle \frac{25}{256}}{\displaystyle \frac{7\pi ^2}{192}}{\displaystyle \frac{\pi ^2}{4}}x_{42}.`$
## 13 Appendix D
### 13.1 Basis of solutions for $`w=1`$, $`w=1/2`$ and $`\mathrm{\hspace{0.17em}\hspace{0.17em}1}+3w+4w^2=0`$.
The basis near $`w=1`$ is (with $`x=1w`$):
$`S_1^{(1)}(x)`$ $`=`$ $`๐ฎ(N_1)(x),`$
$`S_2^{(1)}(x)`$ $`=`$ $`[0,0,0,1,65/24,383/72,\mathrm{}],`$
$`S_3^{(1)}(x)`$ $`=`$ $`S_2^{(1)}(x)\left(\mathrm{ln}(x/24)+2666/75\right)+S_{30}^{(1)}(x),`$
$`S_4^{(1)}(x)`$ $`=`$ $`[0,1,0,0,0,213149176769/914630737500,\mathrm{}],`$
$`S_5^{(1)}(x)`$ $`=`$ $`[0,0,1,0,0,806017240807/426827677500,\mathrm{}],`$
$`S_6^{(1)}(x)`$ $`=`$ $`[0,0,0,0,1,555108887/158084325,\mathrm{}],`$
with:
$`S_{30}^{(1)}(x)=[0,96/5,628/25,0,812657/18000,\mathrm{}].`$
The basis near $`w=1/2`$ reads (with $`x=1+2w`$)
$`S_1^{(1/2)}(x)`$ $`=`$ $`๐ฎ(N_1)(x),`$
$`S_2^{(1/2)}(x)`$ $`=`$ $`[0,0,0,1,8/3,46/9,247/27,\mathrm{}],`$
$`S_3^{(1/2)}(x)`$ $`=`$ $`S_2^{(1/2)}(x)\mathrm{ln}(x)+S_{30}^{(1/2)}(x),`$
$`S_4^{(1/2)}(x)`$ $`=`$ $`[0,1,0,0,0,55489/60345,\mathrm{}],`$
$`S_5^{(1/2)}(x)`$ $`=`$ $`[0,0,1,0,0,159977/80460,\mathrm{}],`$
$`S_6^{(1/2)}(x)`$ $`=`$ $`[0,0,0,0,1,1492/447,\mathrm{}]`$
where:
$`S_{30}^{(1/2)}(x)=[0,3/4,7/8,0,95/144,\mathrm{}].`$
The basis near $`w_1=3/8+i\sqrt{7}/8`$ root of $`1+3w+4w^2`$ is (with $`x=1w/w_1`$)
$`S_1^{(w_1)}(x)`$ $`=`$ $`๐ฎ(N_1)(x),`$
$`S_2^{(w_1)}(x)`$ $`=`$ $`[0,1,49/6461/(64\sqrt{7})i,655/1024747/(1024\sqrt{7})i,\mathrm{}],`$
$`S_3^{(w_1)}(x)`$ $`=`$ $`S_2^{(w_1)}(x)\mathrm{ln}(x)+S_{30}^{(w_1)}(x),S_4^{(w_1)}(x)=[0,0,1,0,0,\mathrm{}],`$
$`S_5^{(w_1)}(x)`$ $`=`$ $`[0,0,0,1,0,\mathrm{}],S_6^{(w_1)}(x)=[0,0,0,0,1,\mathrm{}]`$
with:
$`S_{30}^{(w_1)}(x)=[0,0,657/896+61/(128\sqrt{7})i,41203/43008+1991/(6144\sqrt{7})i,\mathrm{}].`$
### 13.2 Connection matrices for $`w=1`$, $`w=1/2`$ and $`1+3w+4w^2=0`$
For the singular point $`w=1`$, the connection matrix with $`w=0`$ reads
$`C(0,\mathrm{\hspace{0.17em}1})=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]\mathrm{where}\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]\mathrm{and}\left[\begin{array}{c}๐\end{array}\right]\mathrm{read}`$ (151)
$`\left[\begin{array}{ccc}1& 0& 0\\ \multicolumn{3}{c}{}\\ 4+i_{21}i& \frac{\sqrt{3}}{144}i& \frac{\sqrt{3}}{144\pi }\\ \multicolumn{3}{c}{}\\ 2\pi i& \frac{\pi \sqrt{3}}{216}& 0\\ \multicolumn{3}{c}{}\\ 4\frac{4}{\pi }i_{51}\frac{5}{\pi }i_{21}+i_{41}i& r_{42}+\frac{4}{\pi }r_{52}i+\frac{\sqrt{3}}{48}i& \frac{7\sqrt{3}}{48\pi }\\ \multicolumn{3}{c}{}\\ 5\frac{2}{\pi }i_{61}+\left(\frac{2\pi }{3}+\frac{25}{8\pi }\right)i_{21}+i_{51}i& r_{52}+\frac{2}{\pi }r_{62}i\frac{25\sqrt{3}}{1728}i& \frac{5\sqrt{3}}{576\pi }+\frac{\sqrt{3}}{18}i\\ \multicolumn{3}{c}{}\\ \frac{13}{2}+\frac{\pi ^2}{3}+i_{61}i& r_{62}\frac{\pi ^2\sqrt{3}}{432}i\frac{25\sqrt{3}}{2304}i& \frac{11\pi \sqrt{3}}{432}\frac{25\sqrt{3}}{2304\pi }\end{array}\right]`$ (158)
$`\left[\begin{array}{ccc}r_{44}+i_{44}i& r_{45}+i_{45}i& r_{46}+i_{46}i\\ \multicolumn{3}{c}{}\\ \frac{\pi }{4}i_{44}+i_{54}i& \frac{\pi }{4}i_{45}+i_{55}i& \frac{\pi }{4}i_{46}+i_{56}i\\ \multicolumn{3}{c}{}\\ \frac{\pi }{2}i_{54}& \frac{\pi }{2}i_{55}& \frac{\pi }{2}i_{56}\end{array}\right]`$ (162)
where:
$`i_{21}1.838093775180,i_{41}4.136525226980,i_{51}8.13898927603`$
$`i_{61}20.74366088704,r_{42}2.542631644752,r_{52}0.01184208897`$
$`r_{62}4.87108777344,r_{44}1.622875171987,r_{45}1.954781507112`$
$`r_{46}3.51387499953,i_{44}0.158271118920,i_{54}2.13873967059`$
$`i_{45}0.041310289307,i_{55}2.46759854730,i_{46}0.02873064396`$
$`i_{56}4.392293882282,`$
These numbers are such that:
$`i_{46}i_{55}r_{44}+r_{45}i_{56}i_{44}r_{46}i_{55}i_{44}i_{46}r_{45}i_{54}`$
$`+i_{45}r_{46}i_{54}i_{45}i_{56}r_{44}={\displaystyle \frac{468398}{18984375\pi ^2}}`$
The connection matrix between $`w=0`$ and the singular point $`w=1/2`$, reads
$`C(0,1/2)=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]\mathrm{where}\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]\mathrm{and}\left[\begin{array}{c}๐\end{array}\right]\mathrm{read}`$ (168)
$`\left[\begin{array}{ccc}1& 0& 0\\ \multicolumn{3}{c}{}\\ 2\frac{3\mathrm{ln}(2)}{\pi }i& r_{22}& \frac{\sqrt{3}}{9\pi }\\ \multicolumn{3}{c}{}\\ \pi i+3\mathrm{ln}(2)& \frac{\pi \sqrt{3}}{54}+\pi r_{22}i& \frac{\sqrt{3}}{9}i\\ \multicolumn{3}{c}{}\\ r_{41}+i_{41}i& r_{42}+i_{42}i& \frac{11\sqrt{3}}{9\pi }\\ \multicolumn{3}{c}{}\\ r_{51}+i_{51}i& r_{52}+i_{52}i& \frac{5\sqrt{3}}{36\pi }+\frac{7\sqrt{3}}{9}i\\ \multicolumn{3}{c}{}\\ r_{61}+i_{61}i& r_{62}+i_{62}i& \frac{25\sqrt{3}}{144\pi }\frac{13\pi \sqrt{3}}{27}+\frac{5\sqrt{3}}{18}i\end{array}\right],`$ (175)
$`\left[\begin{array}{ccc}r_{44}+i_{44}i& r_{45}+i_{45}i& r_{46}+i_{46}i\\ \multicolumn{3}{c}{}\\ \frac{3\pi }{4}i_{44}+i_{54}i& \frac{3\pi }{4}i_{45}+i_{55}i& \frac{3\pi }{4}i_{46}+i_{56}i\\ \multicolumn{3}{c}{}\\ r_{64}\frac{\pi ^2}{2}i_{44}i& r_{65}\frac{\pi ^2}{2}i_{45}i& r_{66}\frac{\pi ^2}{2}i_{46}i\end{array}\right]`$ (179)
where:
$`r_{22}0.02539959775,r_{41}6.805351589429,r_{51}7.203810787172,`$
$`r_{61}8.75798651623,i_{41}5.23529215352,i_{51}12.14972643902,`$
$`i_{61}7.505979318469,r_{42}0.512271205543,r_{52}0.75497554989,`$
$`r_{62}=2.232400538972,i_{42}0.462196540081,i_{52}0.143220115658,`$
$`i_{62}0.18195427623,r_{44}0.1681290553,r_{45}0.00270658055,`$
$`r_{46}0.00323043290,i_{44}0.14301292413,i_{45}0.690508507395,`$
$`i_{46}1.26354926677,i_{54}0.34844554701,r_{64}0.812327323812,`$
$`i_{55}0.50108648504,r_{65}2.347957990666,i_{56}1.132041888142,`$
$`r_{66}5.35056326640,`$
The connection matrix between $`w=0`$ and the singular point $`w_1=3/8+i\sqrt{7}/8`$ root of $`\mathrm{\hspace{0.17em}1}+3w+4w^2=0`$, reads
$`C(0,w_1)=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]\text{where}\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]\text{and}\left[\begin{array}{c}๐\end{array}\right]\text{read}`$ (185)
$`\left[\begin{array}{ccc}1& 0& 0\\ \multicolumn{3}{c}{}\\ r_{21}\frac{3}{2\pi }r_{31}i& r_{22}\frac{3}{2\pi }r_{32}i+\frac{275\sqrt{7}}{16384}i& a\\ \multicolumn{3}{c}{}\\ r_{31}+\frac{2\pi }{3}r_{21}i+\frac{2\pi }{3}i& r_{32}+\frac{2\pi }{3}r_{22}i\frac{623\pi }{24576}i& \frac{2\pi }{3}ia\\ \multicolumn{3}{c}{}\\ r_{41}+i_{41}i& r_{42}+i_{42}i& \frac{1}{3}a\\ \multicolumn{3}{c}{}\\ r_{51}+i_{51}i& r_{52}+i_{52}i& (\frac{5}{4}+\frac{2\pi }{3}i)a\\ \multicolumn{3}{c}{}\\ r_{61}+i_{61}i& r_{62}+i_{62}i& (\frac{25}{16}\pi ^2\frac{5\pi }{3}i)a\end{array}\right]`$ (192)
$`\left[\begin{array}{ccc}r_{44}+i_{44}i& r_{45}+i_{45}i& r_{46}+i_{46}i\\ \multicolumn{3}{c}{}\\ r_{54}+i_{54}i& r_{55}+i_{55}i& r_{56}+i_{56}i\\ \multicolumn{3}{c}{}\\ r_{64}+i_{64}i& r_{65}+i_{65}i& r_{66}+i_{66}i\end{array}\right]`$ (196)
where:
$`a={\displaystyle \frac{825\sqrt{7}1869i}{16384\pi }},`$
$`r_{21}0.30983963151,r_{31}1.38629436111,r_{22}0.07996746793,`$
$`r_{32}0.044743829620,r_{41}4.70316610599,i_{41}5.10203220992,`$
$`r_{42}0.028522637766,i_{42}0.03731267544,r_{51}1.404170417754,`$
$`i_{51}10.77185269595,r_{52}0.25654299002,i_{52}0.03695328252,`$
$`r_{61}6.98898250954,i_{61}17.585497074,r_{62}0.18342705750,`$
$`i_{62}1.339914984659,r_{44}0.00394832042,i_{44}0.043931830095,`$
$`r_{45}0.02716280332,i_{45}0.0900753899,r_{46}0.070134204478,`$
$`i_{46}0.050869745772,r_{54}0.2122947699,i_{54}0.033562029788,`$
$`r_{55}0.496361798471,i_{55}0.00455966493,r_{56}0.36867647137,`$
$`i_{56}0.040697038977,r_{64}0.1279407612,i_{64}0.68382860060,`$
$`r_{65}0.14739127007,i_{65}1.64596123266,r_{66}0.189914623980,`$
$`i_{66}1.29483325656,`$
## 14 Appendix E: Monodromy matrices in the $`w=0`$-basis
The monodromy matrix around $`w=0`$ expressed in terms of its own $`(w=0)`$ well-suited basis is given in (59).
The monodromy matrix around $`w=1/2`$, expressed in terms of the $`(w=0)`$ well-suited basis, after a conjugation similar to (60), and thus using the previously given connection matrices, reads in terms of $`\alpha `$ and $`\mathrm{\Omega }`$:
$`4\alpha ^2M_{w=0}(1/2)(\alpha ,\mathrm{\Omega })=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]`$ (199)
where:
$`\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]=\left[\begin{array}{ccc}4\alpha ^2& 0& 0\\ \multicolumn{3}{c}{}\\ 48\alpha \mathrm{\Omega }& 4\alpha \left(12\mathrm{\Omega }+\alpha \right)& 96\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 24\mathrm{\Omega }\alpha ^2& 24\mathrm{\Omega }\alpha ^2& 4\left(\alpha 12\mathrm{\Omega }\right)\alpha \\ \multicolumn{3}{c}{}\\ 528\alpha \mathrm{\Omega }& 528\alpha \mathrm{\Omega }& 1056\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 12\left(14\alpha +5\right)\alpha \mathrm{\Omega }& 12\left(14\alpha +5\right)\alpha \mathrm{\Omega }& 24\left(14\alpha +5\right)\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ \alpha a\mathrm{\Omega }& \alpha a\mathrm{\Omega }& 2\mathrm{\Omega }a\end{array}\right]`$ (208)
with $`a=\left(75+52\alpha ^2+60\alpha \right)`$ and $`\left[\begin{array}{c}๐\end{array}\right]=\mathrm{\hspace{0.17em}\hspace{0.17em}4}\alpha ^2\mathrm{๐๐}(\mathrm{๐}\times \mathrm{๐})`$.
The monodromy matrix around $`w=\mathrm{\hspace{0.17em}1}/4`$, expressed in terms of the $`(w=0)`$-well suited basis reads:
$`24\alpha ^4M_{w=0}(1/4)(\alpha ,\mathrm{\Omega })=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]`$ (211)
where $`\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]`$ and $`\left[\begin{array}{c}๐\end{array}\right]`$ read respectively:
$`\left[\begin{array}{ccc}24\alpha ^4& 0& 0\\ \multicolumn{3}{c}{}\\ 48\alpha ^4& 24\alpha ^4& 144\alpha ^2\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 0& 0& 24\alpha ^4\\ \multicolumn{3}{c}{}\\ 48\left(5\alpha ^4+8\mathrm{\Omega }^2+8\mathrm{\Omega }^2\alpha ^2\right)& 32\left(4\mathrm{\Omega }\alpha ^275\mathrm{\Omega }15\alpha ^2\right)\mathrm{\Omega }& 48\left(9\alpha ^2+80\mathrm{\Omega }\right)\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 12\left(5\alpha ^2+4\mathrm{\Omega }+4\mathrm{\Omega }\alpha ^2\right)\alpha ^2& 4\left(754\alpha ^2\right)\alpha ^2\mathrm{\Omega }& 300\alpha ^2\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ \left(87+8\alpha ^2\right)\alpha ^4& 0& 3\left(4\alpha ^275\right)\alpha ^2\mathrm{\Omega }\end{array}\right],`$ (218)
and:
$`\left[\begin{array}{ccc}24\alpha ^4& 384\alpha ^2\mathrm{\Omega }& 1536\mathrm{\Omega }^2\\ \multicolumn{3}{c}{}\\ 0& 24\alpha ^4& 192\alpha ^2\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 0& 0& 24\alpha ^4\end{array}\right]`$ (222)
The monodromy matrix around $`w=1/4`$, expressed in terms of the $`(w=0)`$ well-suited basis reads:
$`12\alpha ^4M_{w=0}(1/4)(\alpha ,\mathrm{\Omega })=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]`$ (225)
where $`\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]`$ and $`\left[\begin{array}{c}๐\end{array}\right]`$ read respectively:
$`\left[\begin{array}{ccc}12\alpha ^4& 0& 0\\ \multicolumn{3}{c}{}\\ 48\alpha ^4& 12\alpha ^4& 0\\ \multicolumn{3}{c}{}\\ 24\alpha ^5& 0& 12\alpha ^4\\ \multicolumn{3}{c}{}\\ a_{41}& a_{42}& 192\mathrm{\Omega }\left(10\mathrm{\Omega }3\alpha ^2\right)\\ \multicolumn{3}{c}{}\\ a_{51}& a_{52}& 48\alpha \mathrm{\Omega }\left(5\alpha +20\mathrm{\Omega }6\alpha ^2\right)\\ \multicolumn{3}{c}{}\\ a_{61}& a_{62}& 48\alpha ^2\mathrm{\Omega }\left(5\alpha +10\mathrm{\Omega }3\alpha ^2\right)\end{array}\right],`$ (232)
with:
$`a_{41}=144\alpha ^4192\mathrm{\Omega }^2192\mathrm{\Omega }^2\alpha ^2,`$
$`a_{42}=16\mathrm{\Omega }\left(60\alpha \mathrm{\Omega }+75\mathrm{\Omega }+8\mathrm{\Omega }\alpha ^218\alpha ^3+15\alpha ^2\right),`$
$`a_{51}=12\alpha \left(5\alpha ^3+6\alpha ^4+8\mathrm{\Omega }^2+8\mathrm{\Omega }^2\alpha ^22\alpha \mathrm{\Omega }2\alpha ^3\mathrm{\Omega }\right),`$
$`a_{52}=2\alpha \mathrm{\Omega }\left(300\mathrm{\Omega }+32\mathrm{\Omega }\alpha ^2+240\alpha \mathrm{\Omega }80\alpha ^375\alpha \right),`$
$`a_{61}=\alpha ^2\left(69\alpha ^2+60\alpha ^3+34\alpha ^4+48\mathrm{\Omega }^2+48\mathrm{\Omega }^2\alpha ^224\alpha \mathrm{\Omega }24\alpha ^3\mathrm{\Omega }\right),`$
$`a_{62}=2\alpha ^2\mathrm{\Omega }\left(150\mathrm{\Omega }30\alpha ^2+16\mathrm{\Omega }\alpha ^2+120\alpha \mathrm{\Omega }44\alpha ^375\alpha \right),`$
and:
$`\left[\begin{array}{ccc}12\left(\alpha +4\mathrm{\Omega }\right)^2\alpha ^2& 192\left(\alpha +4\mathrm{\Omega }\right)\alpha \mathrm{\Omega }& 768\mathrm{\Omega }^2\\ \multicolumn{3}{c}{}\\ 24\alpha ^3\mathrm{\Omega }\left(\alpha +4\mathrm{\Omega }\right)& 12\alpha ^2\left(\alpha ^2+32\mathrm{\Omega }^2\right)& 96\left(\alpha +4\mathrm{\Omega }\right)\alpha \mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ 48\alpha ^4\mathrm{\Omega }^2& 48\left(\alpha 4\mathrm{\Omega }\right)\alpha ^3\mathrm{\Omega }& 12\left(\alpha +4\mathrm{\Omega }\right)^2\alpha ^2\end{array}\right]`$ (236)
The monodromy matrix around $`w=\mathrm{}`$, expressed in terms of the $`w=0`$-well suited basis reads:
$`24\alpha ^4M_{w=0}(\mathrm{})(\alpha ,\mathrm{\Omega })=\left[\begin{array}{cc}๐& \mathrm{๐}\\ \multicolumn{2}{c}{}\\ ๐& ๐\end{array}\right]`$ (239)
where $`\left[\begin{array}{c}๐\\ \\ ๐\end{array}\right]`$ and $`\left[\begin{array}{c}๐\end{array}\right]`$ read respectively:
$`\left[\begin{array}{ccc}24\alpha ^4& 0& 0\\ \multicolumn{3}{c}{}\\ 288\alpha ^3\mathrm{\Omega }& 24\alpha ^3\left(\alpha +6\mathrm{\Omega }\right)& 864\mathrm{\Omega }\alpha ^2\\ \multicolumn{3}{c}{}\\ 48\alpha ^4\mathrm{\Omega }& 24\alpha ^4\mathrm{\Omega }& 24\alpha ^3\left(\alpha +6\mathrm{\Omega }\right)\\ \multicolumn{3}{c}{}\\ a_{41}& a_{42}& 96\left(21\alpha ^2+160\mathrm{\Omega }\right)\mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ a_{51}& a_{52}& 120\left(6\alpha ^2\alpha +32\mathrm{\Omega }\right)\alpha \mathrm{\Omega }\\ \multicolumn{3}{c}{}\\ a_{61}& a_{62}& 6\left(20\alpha 22536\alpha ^2+160\mathrm{\Omega }\right)\alpha ^2\mathrm{\Omega }\end{array}\right]`$ (246)
with:
$`a_{41}=96\left(\alpha ^3+16\mathrm{\Omega }\alpha ^216\mathrm{\Omega }\right)\mathrm{\Omega }`$
$`a_{42}=16\left(33\alpha ^360\alpha ^2+240\alpha \mathrm{\Omega }+8\mathrm{\Omega }\alpha ^2600\mathrm{\Omega }\right)\mathrm{\Omega },`$
$`a_{51}=24\left(2\alpha ^315\alpha ^24\alpha +16\mathrm{\Omega }\alpha ^216\mathrm{\Omega }\right)\alpha \mathrm{\Omega },`$
$`a_{52}=4\left(40\alpha ^345\alpha ^2150\alpha +240\alpha \mathrm{\Omega }+8\mathrm{\Omega }\alpha ^2600\mathrm{\Omega }\right)\alpha \mathrm{\Omega },`$
$`a_{61}=6\left(20\alpha ^283\alpha +4\alpha ^3+16\mathrm{\Omega }\alpha ^216\mathrm{\Omega }\right)\alpha ^2\mathrm{\Omega },`$
$`a_{62}=\left(525\alpha 44\alpha ^3+240\alpha \mathrm{\Omega }+8\mathrm{\Omega }\alpha ^2600\mathrm{\Omega }\right)\alpha ^2\mathrm{\Omega },`$
and:
$`\left[\begin{array}{ccc}24\left(\alpha +4\mathrm{\Omega }\right)^2\alpha ^2& 768\mathrm{\Omega }\left(\alpha +4\mathrm{\Omega }\right)\alpha & 6144\mathrm{\Omega }^2\\ \multicolumn{3}{c}{}\\ 24\alpha ^3\mathrm{\Omega }\left(\alpha +4\mathrm{\Omega }\right)& 24\alpha ^2\left(\alpha ^2+32\mathrm{\Omega }^2\right)& 384\alpha \mathrm{\Omega }\left(\alpha +4\mathrm{\Omega }\right)\\ \multicolumn{3}{c}{}\\ 24\alpha ^4\mathrm{\Omega }^2& 48\alpha ^3\mathrm{\Omega }\left(\alpha +4\mathrm{\Omega }\right)& 24\alpha ^2\left(\alpha +4\mathrm{\Omega }\right)^2\end{array}\right]`$ (250) |
warning/0506/astro-ph0506046.html | ar5iv | text | # Limits on Turbulent H I Fluctuations Towards PSR B0329+54 On Scales Between 0.0025 and 12.5AU
## 1 Introduction
During the last decade, studies of H $`\mathrm{I}`$ absorption lines have revealed angular or temporal variations corresponding to spatial scales on the order of tens of AU. In some directions structure in the H $`\mathrm{I}`$ has been observed and interpreted as individual H $`\mathrm{I}`$ clouds passing across the line-of-sight with densities of $`10^4`$$`10^5`$$`\mathrm{cm}^3`$. Based on these measurements a significant fraction (10โ15%) of the cold H $`\mathrm{I}`$ gas must be in these clouds (Frail et al. 1994). However, it is difficult to reconcile the high H $`\mathrm{I}`$ densities at AU-scales implied by these measurements with other information about the interstellar medium (ISM), as they would be greatly out of pressure equilibrium and should be short lived. Heiles (1997) proposed that the observed H $`\mathrm{I}`$ components are formed from sheets and filaments where the large column densities are produced by the appropriate viewing angle. Deshpande (2000) suggests that these data have been misinterpreted and that a single power-law describes the distribution of cold H $`\mathrm{I}`$ in the ISM. It has also been suggested that the observed changes in the H $`\mathrm{I}`$ absorption profiles towards pulsars are the result of scintillation along with a velocity gradient in a uniform H $`\mathrm{I}`$ medium (Gwinn 2001). Angular power spectra from H $`\mathrm{I}`$ emission observations reveal a power-law distribution of H $`\mathrm{I}`$ structures on parsec scales consistent with a turbulent medium (e.g., Green 1993). It seems possible that the AU-scale H $`\mathrm{I}`$ fluctuations are part of the same turbulence that is present on parsec scales, but this has not been conclusively demonstrated.
If AU-scale fluctuations are present and are part of the turbulent cascade seen at parsec scales then it may be possible to determine some of the fundamental characteristics of the H $`\mathrm{I}`$ gas. If the gas is purely hydrodynamical (HD) then the smallest size scale that shows fluctuations, called the inner scale, is larger than the โmolecularโ mean-free-path length (Tennekes & Lumley 1994, Frisch 1996). For typical situations in the ISM the H $`\mathrm{I}`$ mean-free-path length corresponds to scales from $`1100`$ AU. Collisions of the H $`\mathrm{I}`$ with the ions and electrons in magneto-hydrodynamical (MHD) turbulence could create H $`\mathrm{I}`$ fluctuations on scales smaller than $``$ AU. If we are able to measure the inner scale of the HD turbulence then it is possible to estimate the kinematic viscosity of the H $`\mathrm{I}`$ in the ISM. The kinematic viscosity is given approximately by $`\nu \lambda _{\mathrm{mfp}}\mathrm{c}_{\mathrm{th}}`$, where $`\lambda _{\mathrm{mfp}}`$ is the mean-free-path length and $`\mathrm{c}_{\mathrm{th}}`$ is the thermal sound speed. If the H $`\mathrm{I}`$ turbulence is MHD then sub-AU scale H $`\mathrm{I}`$ fluctuations would be present which would tell us that magnetic fields cannot be ignored in any aspect of the dynamics of interstellar gas.
We have made H $`\mathrm{I}`$ absorption measurements towards the pulsar B0329+54 with the Green Bank Telescope (GBT) to measure H $`\mathrm{I}`$ fluctuations on sub-AU scales. B0329+54 has a parallactic distance of $`1.03_{0.12}^{+0.13}`$$`\mathrm{kpc}`$ and a proper motion of $`95_{11}^{+12}`$$`\mathrm{km}\mathrm{sec}^1`$ (Brisken et al. 2002). The pulsarโs proper motion corresponds to scales of 0.03$`\mathrm{AU}`$ in one day and 1$`\mathrm{AU}`$ in one month for an H $`\mathrm{I}`$ cloud half way to the pulsar.
## 2 Observations
The 100$`\mathrm{m}`$ Green Bank telescope (GBT) of the National Radio Astronomy Observatory<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. (NRAO) was used for this measurement. The GBT is an unblocked aperture telescope with a spatial resolution of 9$`\text{.}^{}`$2 at 21$`\mathrm{cm}`$. The system temperature on cold sky was $`20`$ K. The detector was the NRAO spectral processor, an FFT spectrometer, configured to have 1024 channels for each linear polarization with a bandwidth of 1.25$`\mathrm{MHz}`$, producing a spectral resolution of 0.26$`\mathrm{km}\mathrm{sec}^1`$ per channel. The accumulation memory is 32-bit, providing good dynamic range. The spectral processor integration time was 14 pulse periods, approximately 10 seconds.
The pulsar B0329+54 is an ideal target for study of small-scale structure in cold H $`\mathrm{I}`$. It is very bright so a change in opacity of 0.1 can be detected in a time less than the scintillation time-scale of $`15`$ minutes. Its declination is such that it can be observed by the GBT continuously for $`20`$$`\mathrm{hr}`$ during which $`50`$ scintles are seen. The observations reported here were made in 18 separate observing sessions. Three long ($`20`$$`\mathrm{hr}`$) sessions, separated by two weeks, were made to probe H $`\mathrm{I}`$ fluctuations on sub-AU scales. Fifteen short ($`12`$$`\mathrm{hr}`$) observations spread over the subsequent fifteen months sample scales larger than an AU.
The data were calibrated using a method similar to that described by Weisberg (1978). For each integration an absorption spectrum was formed by taking the difference between the pulsar โonโ and pulsar โoffโ spectra. If $`\tau _{\mathrm{i},\mathrm{j}}`$ is the H $`\mathrm{I}`$ opacity for the i<sup>th</sup> spectral channel and the j<sup>th</sup> integration sample then
$$T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{on}})\mathrm{e}^{\tau _{\mathrm{i},\mathrm{j}}}=\frac{\mathrm{T}_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{on}},\mathrm{}_{\mathrm{on}})\mathrm{T}_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{off}},\mathrm{}_{\mathrm{on}})}{\mathrm{T}_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{off}},\mathrm{}_{\mathrm{off}})/_{\mathrm{i}=1}^{\mathrm{n}_{\mathrm{chan}}}\mathrm{T}_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{off}},\mathrm{}_{\mathrm{off}})}$$
(1)
where $`n_{\mathrm{chan}}`$ is the number of spectral channels. The symbol $`T`$ denotes intensity in units of Kelvins where, for example, $`T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{off}},\mathrm{}_{\mathrm{on}})`$ corresponds to the intensity at the frequency of the H $`\mathrm{I}`$ line emission when the pulsar is โoffโ, $`T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{off}},\mathrm{}_{\mathrm{off}})`$ corresponds to the intensity at frequencies other than the H $`\mathrm{I}`$ line emission when the pulsar is โoffโ, and $`T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{on}})`$ corresponds to the intensity at all observed frequencies when the pulsar is โonโ. The power is converted from detector counts to Kelvin by using a calibrated noise diode that was injected every pulsar cycle for 10% of the pulsar period. Each integration is then weighted and summed such that
$$T_\mathrm{i}(\mathrm{p}_{\mathrm{on}})e^{\tau _\mathrm{i}}=\frac{_{j=1}^{n_{\mathrm{int}}}T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{on}})\mathrm{e}^{\tau _{\mathrm{i},\mathrm{j}}}\frac{\mathrm{T}_{\mathrm{i},\mathrm{j}}^2(\mathrm{p}_{\mathrm{on}},\mathrm{}_{\mathrm{off}})}{\mathrm{T}_{\mathrm{sys}_{\mathrm{i},\mathrm{j}}}^2}}{_{j=1}^{n_{\mathrm{int}}}\frac{T_{\mathrm{i},\mathrm{j}}^2(\mathrm{p}_{\mathrm{on}},\mathrm{}_{\mathrm{off}})}{T_{sys_{\mathrm{i},\mathrm{j}}}^2}}$$
(2)
where $`T_{sys_{\mathrm{i},\mathrm{j}}}`$ is the system temperature, $`T_\mathrm{i}(\mathrm{p}_{\mathrm{on}})`$ is the weighted average pulsar flux, and there are $`n_{\mathrm{int}}`$ independent integration samples.
The flux from B0329+54 varies for two reasons: 1) pulse to pulse variations due to intrinsic emission variations; and 2) due to scintillation. Because B0329+54 varies in intensity during our 10$`\mathrm{s}`$ integration period and can be a significant fraction of the total system temperature, some of the emission spectrum will appear in the absorption spectrum as โghostsโ. Weisberg (1978) has shown that the measured pulsar absorption spectrum is a linear combination of $`T_\mathrm{i}(\mathrm{p}_{\mathrm{off}},\mathrm{}_{\mathrm{on}})`$ and the desired pulsar absorption spectrum. Therefore the measured H $`\mathrm{I}`$ emission spectrum is used to fit and remove the observed โghostsโ. Simultaneously a polynomial model is used to remove the average pulsar flux ($`T_i(\mathrm{p}_{\mathrm{on}})`$), which cannot be independently measured, and any structure in the baseline due to instrumental effects. The โghostโ and polynomial fit is constrained only in spectral regions where there is no H $`\mathrm{I}`$ absorption. The โghostโ spectrum and the polynomial fit is then extrapolated through the regions with H $`\mathrm{I}`$ absorption. Since the 1.25 MHz bandwidth of these observations is a significant fraction of the scintillation bandwidth<sup>2</sup><sup>2</sup>2The scintillation bandwidth is the $`e^1`$ scale over which the pulsarโs signal becomes decorrelated in frequency due to scintillation of the pulsar signal from a non-uniform distribution of electrons in the ISM. It effectively represents the average frequency domain size of an observed scintle. for B0329+54 (measured to be $`5.1\pm 0.1`$ MHz at $`1640`$ MHz by Minter (2001) ), the average pulsar flux contains scintillation induced structures that vary slowly with frequency. In order to remove the baseline structure arising from scintillation structures and instrumental effects, a 5th order polynomial is used.
For each integration the RMS noise was computed for the part of the spectrum without H $`\mathrm{I}`$ line emission for both $`T_\mathrm{j}(\mathrm{p}_{\mathrm{on}})`$ and $`T_\mathrm{j}(\mathrm{p}_{\mathrm{off}})`$. Since the H $`\mathrm{I}`$ line emission contributes significantly to the system temperature ($`20\mathrm{K}`$) the RMS noise will be larger for frequencies where there is H $`\mathrm{I}`$ line emission. Thus the frequency dependent noise is given by
$$\sigma _{T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{on}})}=\sigma _{T_\mathrm{j}(\mathrm{p}_{\mathrm{on}})}\left(1+T_\mathrm{i}^{\mathrm{HI}}/T_{sys}\right)$$
(3)
where $`T_\mathrm{i}^{\mathrm{HI}}`$ is the H $`\mathrm{I}`$ line emission strength in the ith spectral channel. A similar expression is used for $`\sigma _{T_{\mathrm{i},\mathrm{j}}(\mathrm{p}_{\mathrm{off}})}`$. These errors are then propagated for all subsequent calculations.
Systematic errors are also important in the comparison of H $`\mathrm{I}`$ absorption towards pulsars from different epochs. The systematic errors arise from the removal of the average pulsar flux, $`T_i(\mathrm{p}_{\mathrm{on}})`$, via a polynomial fit from the results of equation 2 in order to determine the $`e^\tau `$ spectra. The actual pulsar flux will certainly not follow the value of the polynomial as it is extrapolated over the regions of H $`\mathrm{I}`$ absorption that were excluded from the polynomial fitting. This difference between the actual pulsar flux and the extrapolated polynomial value results in a systematic error that must be taken into account when comparing the H $`\mathrm{I}`$ absorption spectra between different epochs.
In Figure 1 we show an example of how the systematic error is calculated. A frequency band without H $`\mathrm{I}`$ emission or absorption was observed and reduced in the same manner as the pulsar H $`\mathrm{I}`$ absorption data. The upper panel of Figure 1 shows the results of equation 2 and the polynomial fit. The polynomial fit has only been done in the shaded regions, which are the same channel numbers used for the polynomial fit for the H $`\mathrm{I}`$ absorption measurements. The bottom panel of Figure 1 shows the resulting $`e^\tau `$ spectrum. As can been seen in Figure 1, there is excess noise in the regions where the polynomial fit was extrapolated (the white areas in Figure 1). This systematic noise is shown by the single error bar above the data in the bottom panel of Figure 1. The random noise is shown by the two solid lines in the bottom panel of Figure 1. The systematic noise was determined by taking the mean deviation of the $`e^\tau `$ spectra away from $`e^\tau =1`$ within the region where the polynomial fit was extrapolated (the white areas in the bottom panel of Figure 1) and then subtracting, in quadrature, the random noise term. The systematic noise was determined by (i) subtracting the polynomial obtained by omitting the H $`\mathrm{I}`$ absorption channels (shown in Figure 1) with a polynomial obtained by using all of the channels; and (ii) taking the rms of the difference between these two fitted spectra over the omitted channel ranges.
The average systematic error found using this method amounts to $`\pm 0.005`$ in the $`e^\tau `$ spectra. Given that previous detections of H $`\mathrm{I}`$ fluctuations have been recently questioned (e.g. Johnston et al. 2003, Stanimiroviฤ et al. 2003) we choose to be conservative in our error estimates. The quoted uncertainties for our measurements thus contain both random and systematic errors, which have been added in quadrature. We use a value of $`\pm 0.005`$ in the $`e^\tau `$ spectra, as determined above, for the systematic errors in all H $`\mathrm{I}`$ spectra.
## 3 Results
Figure 2 compares the absorption spectrum towards B0329+54 for two epochs, July 14, 2002 and July 2, 2003. The solid lines above and below the difference spectrum are the $`1\sigma `$ random noise uncertainties. The systematic error from the polynomial fitting is indicated by the error bar above the data. There are no significant differences between the two spectra shown in Figure 2. In fact, no variations have been detected over the entire 16 months of observations.
The structure function for time variations in the opacity is defined as
$$D_\tau (\mathrm{\Delta }t,i)=\left(\tau (t,i)\tau (t+\mathrm{\Delta }t,i)\right)^2$$
(4)
where $`\tau (t,i)`$ is the opacity at time $`t`$ in the ith spectral channel. The structure function is actually computed by binning values of $`\left(\tau (t,i)\tau (t+\mathrm{\Delta }t,i)\right)^2`$ and then averaging the data within a given bin. The errors for the structure function are computed from the distribution of values within a given bin. Analysis of the data using a structure function has two advantages: (1) it reduces the noise for a given delay through averaging; and (2) it provides information on the power spectrum of H $`\mathrm{I}`$ absorption fluctuations. If the power spectrum of opacity variations is a power law, as is the case for turbulent fluctuations, then the structure function will also be a power law. If only noise is present in the data then the structure function will be constant at a level of $`2\sigma _{\tau (i)}^2`$. (We refer the reader to Spangler et al. (1989) for a more detailed discussion of structure functions.)
Figure 3 plots the structure function for our observations of B0329+54 for two frequency channels. The structure function data are consistent with there being no turbulence in the H $`\mathrm{I}`$ absorption down to our detection limits. Figure 4 shows our observational upper limits to the turbulent fluctuations in opacity derived from the structure functions in each spectral channel. The data cover differences between ten minutes to sixteen months corresponding to angular scales of 0.37$`\mu \mathrm{as}`$ to 23.8$`\mathrm{mas}`$ and geometric linear size scales of $`0.002512.5`$$`\mathrm{AU}`$, assuming H $`\mathrm{I}`$ gas half way to the pulsar and using the pulsarโs proper motion velocity of 95$`\mathrm{km}\mathrm{sec}^1`$ and distance of 1.03$`\mathrm{kpc}`$. (For the remainder of this paper we calculate distance scales assuming a distance halfway to the pulsar of $`515`$$`\mathrm{pc}`$.) Furthermore, comparing our spectra to those of Gordon et al. (1969), who first detected H $`\mathrm{I}`$ absorption towards B0329+54, shows no evidence for variations in H $`\mathrm{I}`$ absorption greater than $`\mathrm{\Delta }e^\tau 0.01`$ (the noise level of the Gordon et al. measurements is $`\mathrm{\Delta }e^\tau 0.01`$ while the noise for our measurements is $`\mathrm{\Delta }e^\tau 0.002`$) with a velocity resolution of 1.7$`\mathrm{km}\mathrm{sec}^1`$. This corresponds to a linear scale of $`350`$$`\mathrm{AU}`$.
Strong interstellar scintillation is observed towards B0329+54. The multi-path propagation that results from scintillation effectively creates a spatial smoothing of H $`\mathrm{I}`$ absorption structures within the observed angularly broadened size of the pulsar which might mask absorption variations that exist on very small scales. Semenkov et al. (2003) have limited the angular broadening size of B0329+54 at 1600$`\mathrm{MHz}`$ to be $`1.8`$$`\mathrm{mas}`$ corresponding to a scintillation smoothing scale of $`1.15`$$`\mathrm{AU}`$ assuming a $`\lambda ^2`$ scaling of the angular broadening size. Our observations probe scales up to $`12.5`$$`\mathrm{AU}`$ and thus scintillation cannot account for a lack of H $`\mathrm{I}`$ absorption fluctuations in the data. The scintillation smoothing provides a lower limit to the smallest scale opacity structure that we can probe. However we only have upper limits to the scattering disk size. Therefore we use the geometric linear size scale of 0.0025$`\mathrm{AU}`$ (see above) rather than the scattering disk size as the smallest size scale probed by our observations.
Gwinn (2001) suggests that interstellar scintillation coupled with gradients in the Doppler velocity of H $`\mathrm{I}`$ can produce small-scale fluctuations in H $`\mathrm{I}`$ absorption spectra towards pulsars. The GBT observations of B0329+54 can decouple density fluctuations from velocity gradients since an absorption spectrum can be measured for individual scintles. The data points in Figure 3 are divided into three groups. The cross symbols are data averaged over a single scintle, the square symbols consist of $`2`$$`\mathrm{hr}`$ averages within a given epoch, and the triangle symbols are data averaged over a single observing epoch. No turbulent fluctuations are detected on any timescales from the scintillation timescale to sixteen months in our observations.
Gwinn (2001) predicts that the opacity variations that would be observed are given by $`\sigma _\tau =\sqrt{\frac{\tau \mathrm{\Delta }v}{c_s}}`$ where $`\mathrm{\Delta }v`$ is the amount by which the velocity of the absorption feature changes and $`c_s`$ is the thermal sound speed. We have fit Gaussians to the peak absorption features for each epoch. From these fits we find no trend in the value of the line center and can limit any change in velocity of the absorbing gas to $`\mathrm{\Delta }v<0.065`$$`\mathrm{km}\mathrm{sec}^1`$. We can derive upper limits to the thermal sound speed from the H $`\mathrm{I}`$ spin temperature ($`\mathrm{T}_\mathrm{s}`$) found from comparing the H $`\mathrm{I}`$ emission spectrum ($`\mathrm{T}_{\mathrm{em}}`$) with the H $`\mathrm{I}`$ absorption spectrum. The spin temperature is found from $`\mathrm{T}_\mathrm{s}=\frac{\mathrm{T}_{\mathrm{em}}}{1\mathrm{e}^\tau }`$ and should be considered an upper limit since some warm H $`\mathrm{I}`$ contributes to $`\mathrm{T}_{\mathrm{em}}`$ but typically does not contribute to the H $`\mathrm{I}`$ absorption. The spin temperature upper limits are shown in Figure 5. The sound speed can then be estimated from $`c_s=0.093\sqrt{\mathrm{T}_\mathrm{s}(\mathrm{K})}`$$`\mathrm{km}\mathrm{sec}^1`$. An estimate of Gwinnโs $`\sigma _\tau `$ prediction is shown as the dotted line in Figure 4. The derived spin temperatures are not likely to be incorrect by more than a factor of a few. We thus consider the predictions using Gwinnโs formula for opacity variations shown in Figure 4 to be a reasonable upper limit for opacity variations induced by the combination of scintillation and velocity gradients. From Figure 4 it can be seen that our measured limits are smaller than the upper limits determined for Gwinnโs $`\sigma _\tau `$ formula. Although this suggests that Gwinnโs hypothesis could be incorrect, we cannot claim this with any certainty since only an upper limit can be determined for Gwinnโs prediction.
## 4 Discussion
Our GBT observations towards B0329+54 were made primarily to probe sub-AU H $`\mathrm{I}`$ structures. The structure functions of the change in opacity versus time are consistent with noise (i.e. the structure functions are constant values). The structure functions thus place upper limits on the opacity variations as discussed in ยง 3. So we do not detect H $`\mathrm{I}`$ opacity variations consistent with a turbulent power law distribution on scales $`<12.5`$$`\mathrm{AU}`$ greater than $`0.026`$ for the $`31`$, $`21`$, $`18`$, and $`+4`$$`\mathrm{km}\mathrm{sec}^1`$ absorption lines, $`0.12`$ for the $`11`$$`\mathrm{km}\mathrm{sec}^1`$ absorption line, and $`0.055`$ for the $`1`$$`\mathrm{km}\mathrm{sec}^1`$ absorption line. This is somewhat surprising since Frail et al. (1994) detected H $`\mathrm{I}`$ variations of $`\mathrm{\Delta }\tau 0.1`$ in all observed pulsars. However, Johnston et al. (2003) have made multi-epoch observations of H $`\mathrm{I}`$ absorption towards four southern pulsars and find no significant variations and Stanimiroviฤ et al. (2003) reach a similar conclusion from re-observations of the Frail et al. pulsars to increase the number of temporal baselines. VLBA H $`\mathrm{I}`$ absorption measurements have been summarized by Faison (2002). Only two sources show significant H $`\mathrm{I}`$ fluctuations. Small-scale H $`\mathrm{I}`$ structure measured through H $`\mathrm{I}`$ absorption is currently found in only two of the 15 sources observed in Johnston et al. (2003), Stanimiroviฤ et al. (2003), Faison (2002) and this work.
On large scales the distribution of H $`\mathrm{I}`$ is influenced by spiral density waves, supernovae, etc. as is evident by the observed shells and filaments in neutral hydrogen surveys. Dickey & Lockman (1990) have argued that no more than 10% of the total H $`\mathrm{I}`$ exists in small-scale structures ($`1`$$`\mathrm{pc}`$). That is, there are not large variations in the H $`\mathrm{I}`$ column density on small spatial scales and the concept of distinct H $`\mathrm{I}`$ clouds does not describe most of the Galactic neutral hydrogen. The distribution of H $`\mathrm{I}`$ can be analyzed by producing angular power spectra of H $`\mathrm{I}`$ emission over different spatial scales. The results are well fit by a power-law with a slope of approximately $`3`$ towards different directions in the Galaxy (e.g., Crovisier & Dickey 1983; Green 1993; Dickey et al. 2001). These results can be described by a turbulent cascade of energies (Lazarian & Pogosyan 2000). But does this turbulence extend down to very small spatial scales? In other words, what is the inner scale of the turbulence?
We are unable to answer this question with the current data on B0329+54 since we have not detected any significant turbulent H $`\mathrm{I}`$ fluctuations.
Of particular note are the results of Shishov et al. (2003) from diffractive scintillation measurements of B0329+54. They find that on scales below $`3\times 10^{16}\mathrm{cm}`$ that the diffractive scintillation can be explained with a scattering screen comprised solely of ionized gas.<sup>3</sup><sup>3</sup>3The value $`3\times 10^{16}\mathrm{cm}`$ was obtained by Shishov et al. (2003) using a velocity of $`139`$$`\mathrm{km}\mathrm{sec}^1`$ for the pulsar. Using a velocity of $`95`$$`\mathrm{km}\mathrm{sec}^1`$ as measured by Brisken et al. (2002) results in a scale of $`2\times 10^{16}\mathrm{cm}`$. However, on scales above $`3\times 10^{16}\mathrm{cm}`$ their results require that some neutral gas is also present within the scattering screen. Could this indicate that the inner scale for the neutral gas is $`2000`$ AU? If this is the inner scale then an H $`\mathrm{I}`$ density of $`1\mathrm{cm}^3`$ at a temperature of $`100\mathrm{K}`$ would give a kinematic viscosity of $`3\times 10^{21}\mathrm{cm}^2\mathrm{s}^1`$.
The research of J.S.K. at NRAO was supported by the NSF Research Experiences for Undergraduates program. We thank Crystal Brogan, Avinash Deshpande, and Sneลบana Stanimiroviฤ for stimulating discussions. We thank Jay Lockman for many discussions about the observations and calibration procedures and for commenting on the manuscript. We would like to thank the anonymous referee whose comments greatly improved the quality of this work. |
warning/0506/gr-qc0506014.html | ar5iv | text | # Appendix
## Appendix
In this Appendix, we first provide the form of the Riemann tensor for our time-dependent metric ansatz; it depends on five functions, reducing to for $`\dot{a}=\dot{b}=0`$:
$`R_{\mu \nu \alpha \beta }`$ $`=`$ $`4A\delta _{[\mu }^0\delta _{\nu ]}^r\delta _{[\alpha }^0\delta _{\beta ]}^r+4BZ_i\delta _{[\mu }^0\delta _{\nu ]}^{\theta _i}\delta _{[\alpha }^0\delta _{\beta ]}^{\theta _i}+4CZ_i\delta _{[\mu }^r\delta _{\nu ]}^{\theta _i}\delta _{[\alpha }^r\delta _{\beta ]}^{\theta _i}+2r^4\psi Z_iZ_j\delta _{[\mu }^{\theta _i}\delta _{\nu ]}^{\theta _j}\delta _{[\alpha }^{\theta _i}\delta _{\beta ]}^{\theta _j}`$ (8)
$`+`$ $`4DZ_i\left(\delta _{[\mu }^0\delta _{\nu ]}^{\theta _i}\delta _{[\alpha }^r\delta _{\beta ]}^{\theta _i}+\delta _{[\mu }^r\delta _{\nu ]}^{\theta _i}\delta _{[\alpha }^0\delta _{\beta ]}^{\theta _i}\right),`$
with
$`A`$ $`=`$ $`{\displaystyle \frac{b}{2}}\left[3a^{}b^{}+a^{\prime \prime }b+2ab^{\prime \prime }\right]+{\displaystyle \frac{1}{2a^2}}\left[{\displaystyle \frac{\dot{a}\dot{b}}{b}}+{\displaystyle \frac{a\ddot{a}2(\dot{a})^2}{a}}\right]`$ (9)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}(rab)\left[a^{}b+2ab^{}\right]`$ (10)
$`C`$ $`=`$ $`{\displaystyle \frac{a^{}r}{2a}},\text{ }D={\displaystyle \frac{\dot{a}r}{2a}},\text{ }\psi ={\displaystyle \frac{1a}{r^2}}`$ (11)
and $`Z_i`$ is the $`i^{\text{th}}`$ angular component of the metric.
We are interested in relating $`๐ข^{0r}(k)`$ and $`๐ข^{0r}(k1)`$ to exhibit their โpolynomial in $`\psi `$โ nature, outlining the steps here. We find that effectively, the Riemann tensor appears in the field equation in the combination
$$\stackrel{~}{R}_{\alpha \beta \gamma \rho }=2\psi g_{\alpha \beta }g_{\gamma \rho }+4Dr^2g_{\gamma \rho }\delta _\alpha ^r\delta _\beta ^0.$$
(12)
The first term above is $`๐ข^{0r}(k1)`$ while the second is (in odd dimensions) $`k^1G^{0r}(k)`$ when replacing the $`k^{\text{th}}`$ Riemann tensor in the field equation (5). After some algebra, this leads to
$$๐ข^{0r}(k)=2\psi ๐ข^{0r}(k1)+k^1๐ข^{0r}(k)๐ข^{0r}(k)=\frac{2k\psi }{k1}๐ข^{0r}(k1),$$
(13)
and so
$$๐ข^{0r}(k)=(\text{const.})\times \psi ^{k1}๐ข^{0r}(1);$$
(14)
the final term here is Einstein, and contains a factor of $`D`$ and a product of all the angular dependence but no $`\psi `$.
Our other task here is, for concreteness, to show explicitly in D=5 how the seeming obstructions to Birkhoff are really deSitter vacua, as shown generically in text.
For the metric ansatz:
$$ds^2=ab^2dt^2+a^1dr^2+2bfdtdr+r^2(d\psi ^2+\mathrm{sin}^2\psi (d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2))$$
(15)
we form the cosmological, Einstein, and Gauss-Bonnett actions
$$I=d^5x\sqrt{g}\left(\mathrm{\Lambda }+R+\alpha (R^24R^{\mu \nu }R_{\mu \nu }+R^{\mu \nu \alpha \beta }R_{\mu \nu \alpha \beta })\right),$$
(16)
keeping only terms linear in $`f`$ and its derivatives. The field equations at $`f=0`$ read:
$`{\displaystyle \frac{\delta I}{\delta a}}|_{f=0}`$ $`=`$ $`3(4\alpha (1a)+r^2)b^{}=0`$ (17)
$`{\displaystyle \frac{\delta I}{\delta b}}|_{f=0}`$ $`=`$ $`\mathrm{\Lambda }r^3+6r(1a)3r^2a^{}12\alpha a^{}(1a)`$ (18)
$`๐ข^{0r}`$ $``$ $`{\displaystyle \frac{\delta I}{\delta f}}|_{f=0}=3(r^2+4\alpha (1a))\dot{a}a^1.`$ (19)
Now consider the obstruction, $`4\alpha (1a)+r^2=0`$, to Birkhoff. This is just $`a=1+\frac{r^2}{4\alpha }`$, vacuum deSitter. The constant $`4\alpha `$ is determined by inserting into the $`b`$ equation:
$$\frac{\delta I}{\delta b}|_{f=0,a=1+r^2/(4\alpha )}=\frac{r^3(2\alpha \mathrm{\Lambda }3)}{2\alpha }=0.$$
(20)
The Lovelock coefficients determine the deSitter cosmological constant according to the tuning $`\alpha =3\mathrm{\Lambda }/2`$; in the absence of $`\mathrm{\Lambda }`$, there are no roots, except flat space.
We thank B. Tekin for a discussion. This work was supported by NSF grant PHY-04-01667. |
warning/0506/hep-ph0506239.html | ar5iv | text | # Doubleโlepton polarization asymmetries in the Exclusive ๐ตโ๐โขโโบโขโโป decay beyond the Standard Model
## 1 Introduction
Rare $`B`$ meson decays, induced by flavor changing neutral current (FCNC) $`bs(d)\mathrm{}^+\mathrm{}^{}`$ transitions constitute one of the most important class of decays for testing the gauge structure of the Standard Model (SM). These decays which are forbidden in the SM at tree level, occur at loop level and provide insight to check the predictions of the SM at quantum level. Moreover, these decays are also quite sensitive to the existence of new physics beyond the SM, since new particles running at loops can give contribution to these decays. The new physics manifests itself in rare decays in two different ways; one via modification of the existing Wilson coefficients in the SM, or through the introduction of some new operators with new coefficients which are absent in the SM. Some of the most important exclusive FCNC decays governed by $`bs(d)`$ transition at quark level are $`BK^{}\gamma `$ and $`B(\pi ,\rho ,K,K^{})\mathrm{}^+\mathrm{}^{}`$ decays. The decays of the kind $`BM\mathrm{}^+\mathrm{}^{}`$, where $`M`$ stands for pseudoscalar or vector mesons, enable the investigation of the experimental observables, such as, lepton pair forwardโbackward (FB) asymmetry, lepton polarizations, etc. One of the most efficient ways in looking for new physics beyond the SM is the measurement of lepton polarization in the decays. Polarization of a single lepton has been studied in $`BK^{}\mathrm{}^+\mathrm{}^{}`$ , $`BX_s\mathrm{}^+\mathrm{}^{}`$ , $`BK\mathrm{}^+\mathrm{}^{}`$ , $`B\pi (\rho )\mathrm{}^+\mathrm{}^{}`$ and $`B\mathrm{}^+\mathrm{}^{}\gamma `$ decays in detail in fitting the parameters of the SM and set constraints on new physics beyond the SM. Moreover, as has already been pointed out in , some of the single lepton polarization asymmetries might be quite small to be observed and might not provide sufficient number of observables in checking the structure of the effective Hamiltonian. By taking both lepton polarizations into account simultaneously, maximum number of independent observables are constructed. It is clear that, measurement of many more observables which would be useful in further improvement of the parameters of the SM probing new physics beyond the SM. It should be noted here that both lepton polarizations in the $`BK^{}\tau ^+\tau ^{}`$ and $`BK\mathrm{}^+\mathrm{}^{}`$ decays are studied in and , respectively. The decays of $`B`$ mesons induced by the $`bd\mathrm{}^+\mathrm{}^{}`$ transition are promising in looking for CP violation since the CKM factors $`V_{tb}V_{td}^{}`$, $`V_{ub}V_{ud}^{}`$ and $`V_{cb}V_{cd}^{}`$ in the SM are all of the same order. For this reason CP violation is much more considerable in the decays induced by $`bd`$ transition. So, study of the exclusive decays $`B_d(\pi ,\rho ,\eta )\mathrm{}^+\mathrm{}^{}`$ are quite promising for the confirmation of the CP violation and these decays have extensively been investigated in the SM and beyond .
The aim of the present work is to study the doubleโlepton polarization asymmetries in the exclusive $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay in a model independent way, including all possible forms of interactions into the effective Hamiltonian. Moreover, we study the correlation between the doubleโlepton polarization asymmetries and the branching ratio of the $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay, in order to find such regions of new Wilson coefficients in which the branching ratio (as well as singleโlepton polarization) coincides with the SM prediction while the doubleโlepton polarization asymmetries do not. It is clear that if such a region of the new Wilson coefficients exists it is an indication of the fact that new physics beyond the SM can be established by measurement of the doubleโlepton polarizations only. Note that the doubleโlepton polarizations in the $`BK\mathrm{}^+\mathrm{}^{}`$ and $`B\mathrm{}^+\mathrm{}^{}\gamma `$ decays are studied in and in detail.
The paper is organized as follows. In section 2, using a general form of the effective Hamiltonian, we obtain the matrix element in terms of the form factors of the $`B\rho `$ transition. In section 3 we derive the analytical results for the nine doubleโlepton polarization asymmetries. Last section is devoted to the numerical analysis, discussion and conclusions.
## 2 Double lepton polarization asymmetries in $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay
In this section we calculate the double lepton polarizations using a general form of the effective Hamiltonian. The $`B\rho \mathrm{}^+\mathrm{}^{}`$ process is governed by $`bd\mathrm{}^+\mathrm{}^{}`$ transition at quark level. The matrix element for the $`bd\mathrm{}^+\mathrm{}^{}`$ transition can be written in terms of the twelve model independent fourโFermi interactions in the following form:
$`_{eff}`$ $`=`$ $`{\displaystyle \frac{G_F\alpha }{\sqrt{2}\pi }}V_{td}V_{tb}^{}\{C_{SL}\overline{d}_Ri\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}b_L\overline{\mathrm{}}\gamma ^\mu \mathrm{}+C_{BR}\overline{d}_Li\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}b_R\overline{\mathrm{}}\gamma ^\mu \mathrm{}`$ (1)
$`+`$ $`C_{LL}^{tot}\overline{d}_L\gamma _\mu b_L\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L+C_{LR}^{tot}\overline{d}_L\gamma _\mu b_L\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{RL}\overline{d}_R\gamma _\mu b_R\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L`$
$`+`$ $`C_{RR}\overline{d}_R\gamma _\mu b_R\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{LRLR}\overline{d}_Lb_R\overline{\mathrm{}}_L\mathrm{}_R+C_{RLLR}\overline{d}_Rb_L\overline{\mathrm{}}_L\mathrm{}_R`$
$`+`$ $`C_{LRRL}\overline{d}_Lb_R\overline{\mathrm{}}_R\mathrm{}_L+C_{RLRL}\overline{d}_Rb_L\overline{\mathrm{}}_R\mathrm{}_L+C_T\overline{d}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma ^{\mu \nu }\mathrm{}`$
$`+`$ $`iC_{TE}ฯต^{\mu \nu \alpha \beta }\overline{d}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma _{\alpha \beta }\mathrm{}\},`$
where
$`d_L={\displaystyle \frac{1\gamma _5}{2}}d,d_R={\displaystyle \frac{1+\gamma _5}{2}}d,`$
$`C_X`$ are the coefficients of the fourโFermi interactions and $`q`$ is the momentum transfer. Among all these Wilson coefficients, several already exits in the SM. Indeed, the first two coefficients in Eq. (1), $`C_{SL}`$ and $`C_{BR}`$, are the nonlocal Fermi interactions, which correspond to $`2m_sC_7^{eff}`$ and $`2m_bC_7^{eff}`$ in the SM, respectively. The next four terms with coefficients $`C_{LL}`$, $`C_{LR}`$, $`C_{RL}`$ and $`C_{RR}`$ are the vector type interactions. Two of these vector interactions containing $`C_{LL}^{tot}`$ and $`C_{LR}^{tot}`$ do already exist in the SM in the form $`(C_9^{eff}C_{10})`$ and $`(C_9^{eff}+C_{10})`$. Therefore, $`C_{LL}^{tot}`$ and $`C_{LR}^{tot}`$ can be written as
$`C_{LL}^{tot}`$ $`=`$ $`C_9^{eff}C_{10}+C_{LL},`$
$`C_{LR}^{tot}`$ $`=`$ $`C_9^{eff}+C_{10}+C_{LR},`$
where $`C_{LL}`$ and $`C_{LR}`$ describe the contributions of the new physics. The terms with coefficients $`C_{LRLR}`$, $`C_{RLLR}`$, $`C_{LRRL}`$ and $`C_{RLRL}`$ describe the scalar type interactions. The remaining last two terms lead by the coefficients $`C_T`$ and $`C_{TE}`$, obviously, describe the tensor type interactions.
It should be noted here that, in further analysis we will assume that all new Wilson coefficients are real, as is the case in the SM, while only $`C_9^{eff}`$ contains imaginary part and it is parametrized in the following form
$`C_9^{eff}=\xi _1+\lambda _u\xi _2,`$ (2)
where
$`\lambda _u={\displaystyle \frac{V_{ub}V_{ud}^{}}{V_{tb}V_{td}^{}}},`$
and
$`\xi _1`$ $`=`$ $`4.128+0.138\omega (\widehat{s})+g(\widehat{m}_c,\widehat{s})C_0(\widehat{m}_b){\displaystyle \frac{1}{2}}g(\widehat{m}_d,\widehat{s})(C_3+C_4)`$
$``$ $`{\displaystyle \frac{1}{2}}g(\widehat{m}_b,\widehat{s})(4C_3+4C_4+3C_5+C_6)+{\displaystyle \frac{2}{9}}(3C_3+C_4+3C_5+C_6),`$
$`\xi _2`$ $`=`$ $`[g(\widehat{m}_c,\widehat{s})g(\widehat{m}_u,\widehat{s})](3C_1+C_2),`$ (3)
where $`\widehat{m}_q=m_q/m_b`$, $`\widehat{s}=q^2`$, $`C_0(\mu )=3C_1+C_2+3C_3+C_4+3C_5+C_6`$, and
$`\omega (\widehat{s})`$ $`=`$ $`{\displaystyle \frac{2}{9}}\pi ^2{\displaystyle \frac{4}{3}}Li_2(\widehat{s}){\displaystyle \frac{2}{3}}\mathrm{ln}(\widehat{s})\mathrm{ln}(1\widehat{s}){\displaystyle \frac{5+4\widehat{s}}{3(1+2\widehat{s})}}\mathrm{ln}(1\widehat{s})`$ (4)
$``$ $`{\displaystyle \frac{2\widehat{s}(1+\widehat{s})(12\widehat{s})}{3(1\widehat{s})^2(1+2\widehat{s})}}\mathrm{ln}(\widehat{s})+{\displaystyle \frac{5+9\widehat{s}6\widehat{s}^2}{3(1\widehat{s})(1+2\widehat{s})}},`$
represents the $`O(\alpha _s)`$ correction coming from one gluon exchange in the matrix element of the operator $`๐ช_9`$ , while the function $`g(\widehat{m}_q,\widehat{s})`$ represents oneโloop corrections to the fourโquark operators $`O_1`$$`O_6`$ , whose form is
$`g(\widehat{m}_q,\widehat{s})={\displaystyle \frac{8}{9}}\mathrm{ln}(\widehat{m}_q)+{\displaystyle \frac{8}{27}}+{\displaystyle \frac{4}{9}}y_q{\displaystyle \frac{2}{9}}(2+y_q)`$ (5)
$``$ $`\sqrt{\left|1y_q\right|}\left\{\theta (1y_q)\left[\mathrm{ln}\left({\displaystyle \frac{1+\sqrt{1y_q}}{1\sqrt{1y_q}}}\right)i\pi \right]+\theta (y_q1)\mathrm{arctan}\left({\displaystyle \frac{1}{\sqrt{y_q1}}}\right)\right\},`$
where $`y_q=4\widehat{m}_q^2/\widehat{s}`$.
In addition to the short distance contributions, $`BX_d\mathrm{}^+\mathrm{}^{}`$ decay also receives long distance contributions, which have their origin in the real $`\overline{u}u`$, $`\overline{d}d`$ and $`\overline{c}c`$ intermediate states, i.e., $`\rho `$, $`\omega `$ and $`J/\psi `$ family. There are four different approaches in taking long distance contributions into consideration: a) HQET based approach , b) AMM approach , c) LSW approach , and d) KS approach . In the present work we choose the AMM approach, in which these resonance contributions are parametrized using the BreitโWigner form for the resonant states. The effective coefficient $`C_9^{eff}`$ including the $`\rho `$, $`\omega `$ and $`J/\psi `$ resonances are defined as
$`C_9^{eff}=C_9(\mu )+Y_{res}(\widehat{s}),`$ (6)
where
$`Y_{res}`$ $`=`$ $`{\displaystyle \frac{3\pi }{\alpha ^2}}\{(C^{(0)}(\mu )+\lambda _u[3C_1(\mu )+C_2(\mu )]){\displaystyle \underset{V_i=\psi }{}}K_i{\displaystyle \frac{\mathrm{\Gamma }(V_i\mathrm{}^+\mathrm{}^{})M_{V_i}}{M_{V_i}^2q^2iM_{V_i}\mathrm{\Gamma }_{V_i}}}`$ (7)
$``$ $`\lambda _ug(\widehat{m}_u,\widehat{s})[3C_1(\mu )+C_2(\mu )]{\displaystyle \underset{V_i=\rho ,\omega }{}}{\displaystyle \frac{\mathrm{\Gamma }(V_i\mathrm{}^+\mathrm{}^{})M_{V_i}}{M_{V_i}^2q^2iM_{V_i}\mathrm{\Gamma }_{V_i}}}\}.`$
The phenomenological factor $`K_i`$ has the universal value for the inclusive $`BX_{s(d)}\mathrm{}^+\mathrm{}^{}`$ decay $`K_i2.3`$ , which we use in our calculations.
The decay amplitude for the exclusive $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay is obtained from the matrix element of the effective Hamiltonian in Eq. (1) over $`B`$ and $`\rho `$ meson states, which can be parametrized in terms of various form factors. It follows from (1) that, the following matrix elements
$`\rho \left|\overline{d}\gamma _\mu (1\pm \gamma _5)b\right|B,`$
$`\rho \left|\overline{d}i\sigma _{\mu \nu }q^\nu (1\pm \gamma _5)b\right|B,`$
$`\rho \left|\overline{d}(1\pm \gamma _5)b\right|B,`$
$`\rho \left|\overline{d}\sigma _{\mu \nu }b\right|B,`$
are needed in obtaining the decay amplitude of the $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay. These matrix elements are defined as follows:
$`\rho (p_\rho ,\epsilon )\left|\overline{d}\gamma _\mu (1\pm \gamma _5)b\right|B(p_B)=`$
$`ฯต_{\mu \nu \lambda \sigma }\epsilon ^\nu p_\rho ^\lambda q^\sigma {\displaystyle \frac{2V(q^2)}{m_B+m_\rho }}\pm i\epsilon _\mu ^{}(m_B+m_\rho )A_1(q^2)`$
$`i(p_B+p_\rho )_\mu (\epsilon ^{}q){\displaystyle \frac{A_2(q^2)}{m_B+m_\rho }}iq_\mu {\displaystyle \frac{2m_\rho }{q^2}}(\epsilon ^{}q)\left[A_3(q^2)A_0(q^2)\right],`$
$`\rho (p_\rho ,\epsilon )\left|\overline{d}i\sigma _{\mu \nu }q^\nu (1\pm \gamma _5)b\right|B(p_B)=`$
$`4ฯต_{\mu \nu \lambda \sigma }\epsilon ^\nu p_\rho ^\lambda q^\sigma T_1(q^2)\pm 2i\left[\epsilon _\mu ^{}(m_B^2m_\rho ^2)(p_B+p_\rho )_\mu (\epsilon ^{}q)\right]T_2(q^2)`$
$`\pm 2i(\epsilon ^{}q)\left[q_\mu (p_B+p_\rho )_\mu {\displaystyle \frac{q^2}{m_B^2m_\rho ^2}}\right]T_3(q^2),`$
$`\rho (p_\rho ,\epsilon )\left|\overline{d}\sigma _{\mu \nu }b\right|B(p_B)=`$
$`iฯต_{\mu \nu \lambda \sigma }\{2T_1(q^2)\epsilon _{}^{}{}_{}{}^{\lambda }(p_B+p_\rho )^\sigma +{\displaystyle \frac{2}{q^2}}(m_B^2m_\rho ^2)[T_1(q^2)T_2(q^2)]\epsilon _{}^{}{}_{}{}^{\lambda }q^\sigma `$
$`{\displaystyle \frac{4}{q^2}}[T_1(q^2)T_2(q^2){\displaystyle \frac{q^2}{m_B^2m_\rho ^2}}T_3(q^2)](\epsilon ^{}q)p_\rho ^\lambda q^\sigma \}.`$
where $`q=p_Bp_\rho `$ is the momentum transfer and $`\epsilon `$ is the polarization vector of $`\rho `$ meson.
Note that the matrix element
$`\rho (p_\rho ,\epsilon )\left|\overline{d}\sigma _{\mu \nu }\gamma _5b\right|B(p_B)`$
can easily be obtained from (2)by using the identity
$`\sigma _{\alpha \beta }={\displaystyle \frac{i}{2}}ฯต_{\alpha \beta \rho \sigma }\sigma ^{\rho \sigma }\gamma _5.`$
In order to ensure finiteness of (2) and (2) at $`q^2=0`$, we assume that $`A_3(q^2=0)=A_0(q^2=0)`$ and $`T_1(q^2=0)=T_2(q^2=0)`$. The matrix element $`\rho \left|\overline{d}(1\pm \gamma _5)b\right|B`$ can be calculated by contracting both sides of Eq. (2) with $`q^\mu `$ and using equation of motion. Neglecting the mass of the $`d`$ quark we get
$`\rho (p_\rho ,\epsilon )\left|\overline{d}(1\pm \gamma _5)b\right|B(p_B)={\displaystyle \frac{1}{m_b}}\left[2im_\rho (\epsilon ^{}q)A_0(q^2)\right].`$ (11)
In deriving Eq. (11) we have used the relationship
$`2m_\rho A_3(q^2)=(m_B+m_\rho )A_1(q^2)(m_Bm_\rho )A_2(q^2),`$
which follows from the equations of motion.
Using the definition of the form factors, as given above, the amplitude of the $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay can be written as
$`(B\rho \mathrm{}^+\mathrm{}^{})={\displaystyle \frac{G\alpha }{4\sqrt{2}\pi }}V_{tb}V_{td}^{}`$ (12)
$`\times \{\overline{\mathrm{}}\gamma ^\mu (1\gamma _5)\mathrm{}[2A_1ฯต_{\mu \nu \lambda \sigma }\epsilon ^\nu p_\rho ^\lambda q^\sigma iB_1\epsilon _\mu ^{}+iB_2(\epsilon ^{}q)(p_B+p_\rho )_\mu +iB_3(\epsilon ^{}q)q_\mu ]`$
$`+\overline{\mathrm{}}\gamma ^\mu (1+\gamma _5)\mathrm{}\left[2C_1ฯต_{\mu \nu \lambda \sigma }\epsilon ^\nu p_\rho ^\lambda q^\sigma iD_1\epsilon _\mu ^{}+iD_2(\epsilon ^{}q)(p_B+p_\rho )_\mu +iD_3(\epsilon ^{}q)q_\mu \right]`$
$`+\overline{\mathrm{}}(1\gamma _5)\mathrm{}\left[iB_4(\epsilon ^{}q)\right]+\overline{\mathrm{}}(1+\gamma _5)\mathrm{}\left[iB_5(\epsilon ^{}q)\right]`$
$`+4\overline{\mathrm{}}\sigma ^{\mu \nu }\mathrm{}\left(iC_Tฯต_{\mu \nu \lambda \sigma }\right)\left[2T_1\epsilon _{}^{}{}_{}{}^{\lambda }(p_B+p_\rho )^\sigma +B_6\epsilon _{}^{}{}_{}{}^{\lambda }q^\sigma B_7(\epsilon ^{}q)p_{\rho }^{}{}_{}{}^{\lambda }q^\sigma \right]`$
$`+16C_{TE}\overline{\mathrm{}}\sigma _{\mu \nu }\mathrm{}[2T_1\epsilon _{}^{}{}_{}{}^{\mu }(p_B+p_\rho )^\nu +B_6\epsilon _{}^{}{}_{}{}^{\mu }q^\nu B_7(\epsilon ^{}q)p_{\rho }^{}{}_{}{}^{\mu }q^\nu \},`$
where
$`A_1`$ $`=`$ $`(C_{LL}^{tot}+C_{RL}){\displaystyle \frac{V}{m_B+m_\rho }}2(C_{BR}+C_{SL}){\displaystyle \frac{T_1}{q^2}},`$
$`B_1`$ $`=`$ $`(C_{LL}^{tot}C_{RL})(m_B+m_\rho )A_12(C_{BR}C_{SL})(m_B^2m_\rho ^2){\displaystyle \frac{T_2}{q^2}},`$
$`B_2`$ $`=`$ $`{\displaystyle \frac{C_{LL}^{tot}C_{RL}}{m_B+m_\rho }}A_22(C_{BR}C_{SL}){\displaystyle \frac{1}{q^2}}\left[T_2+{\displaystyle \frac{q^2}{m_B^2m_\rho ^2}}T_3\right],`$
$`B_3`$ $`=`$ $`2(C_{LL}^{tot}C_{RL})m_\rho {\displaystyle \frac{A_3A_0}{q^2}}+2(C_{BR}C_{SL}){\displaystyle \frac{T_3}{q^2}},`$
$`C_1`$ $`=`$ $`A_1(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$
$`D_1`$ $`=`$ $`B_1(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$
$`D_2`$ $`=`$ $`B_2(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$
$`D_3`$ $`=`$ $`B_3(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$
$`B_4`$ $`=`$ $`2(C_{LRRL}C_{RLRL}){\displaystyle \frac{m_\rho }{m_b}}A_0,`$
$`B_5`$ $`=`$ $`2(C_{LRLR}C_{RLLR}){\displaystyle \frac{m_\rho }{m_b}}A_0,`$
$`B_6`$ $`=`$ $`2(m_B^2m_\rho ^2){\displaystyle \frac{T_1T_2}{q^2}},`$
$`B_7`$ $`=`$ $`{\displaystyle \frac{4}{q^2}}\left(T_1T_2{\displaystyle \frac{q^2}{m_B^2m_\rho ^2}}T_3\right).`$ (13)
From this expression of the decay amplitude, for the unpolarized differential decay width we get the following result:
$`{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(B\rho \mathrm{}^+\mathrm{}^{})={\displaystyle \frac{G^2\alpha ^2m_B}{2^{14}\pi ^5}}\left|V_{tb}V_{td}^{}\right|^2\lambda ^{1/2}(1,\widehat{r},\widehat{s})v\mathrm{\Delta }(\widehat{s}),`$ (14)
with
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{2m_B^2}{3\widehat{r}_\rho \widehat{s}}}\text{Re}\{6m_B\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_1D_1)(B_4^{}B_5^{})`$ (15)
$``$ $`12m_B^2\widehat{m}_{\mathrm{}}^2\widehat{s}\lambda \left[B_4B_5^{}+(B_3D_2D_3)B_1^{}(B_2+B_3D_3)D_1^{}\right]`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\widehat{s}(1\widehat{r}_\rho )\lambda (B_2D_2)(B_4^{}B_5^{})`$
$`+`$ $`12m_B^4\widehat{m}_{\mathrm{}}^2\widehat{s}(1\widehat{r}_\rho )\lambda (B_2D_2)(B_3^{}D_3^{})`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\lambda \widehat{s}^2(B_4B_5)(B_3^{}D_3^{})`$
$`+`$ $`48\widehat{m}_{\mathrm{}}^2\widehat{r}_\rho \widehat{s}\left(3B_1D_1^{}+2m_B^4\lambda A_1C_1^{}\right)`$
$`+`$ $`48m_B^5\widehat{m}_{\mathrm{}}\widehat{s}\lambda ^2(B_2+D_2)B_7^{}C_{TE}^{}`$
$``$ $`16m_B^4\widehat{r}_\rho \widehat{s}(\widehat{m}_{\mathrm{}}^2\widehat{s})\lambda \left(\left|A_1\right|^2+\left|C_1\right|^2\right)`$
$``$ $`m_B^2\widehat{s}(2\widehat{m}_{\mathrm{}}^2\widehat{s})\lambda \left(\left|B_4\right|^2+\left|B_5\right|^2\right)`$
$``$ $`48m_B^3\widehat{m}_{\mathrm{}}\widehat{s}(1\widehat{r}_\rho \widehat{s})\lambda \left[(B_1+D_1)B_7^{}C_{TE}^{}+2(B_2+D_2)B_6^{}C_{TE}^{}\right]`$
$``$ $`6m_B^4\widehat{m}_{\mathrm{}}^2\widehat{s}\lambda \left[2(2+2\widehat{r}_\rho \widehat{s})B_2D_2^{}\widehat{s}\left|(B_3D_3)\right|^2\right]`$
$`+`$ $`96m_B\widehat{m}_{\mathrm{}}\widehat{s}(\lambda +12\widehat{r}_\rho \widehat{s})(B_1+D_1)B_6^{}C_{TE}^{}`$
$`+`$ $`8m_B^2\widehat{s}^2\left[v^2\right|C_T|^2+4(32v^2)\left|C_{TE}|^2\right]\left[4(\lambda +12\widehat{r}_\rho \widehat{s})\right|B_6|^2`$
$``$ $`4m_B^2\lambda (1\widehat{r}_\rho \widehat{s})B_6B_7^{}+m_B^4\lambda ^2\left|B_7|^2\right]`$
$``$ $`4m_B^2\lambda \left[\widehat{m}_{\mathrm{}}^2(22\widehat{r}_\rho +\widehat{s})+\widehat{s}(1\widehat{r}_\rho \widehat{s})\right](B_1B_2^{}+D_1D_2^{})`$
$`+`$ $`\widehat{s}\left[6\widehat{r}_\rho \widehat{s}(3+v^2)+\lambda (3v^2)\right]\left(\left|B_1\right|^2+\left|D_1\right|^2\right)`$
$``$ $`2m_B^4\lambda \left\{\widehat{m}_{\mathrm{}}^2[\lambda 3(1\widehat{r}_\rho )^2]\lambda \widehat{s}\right\}\left(\left|B_2\right|^2+\left|D_2\right|^2\right)`$
$`+`$ $`128m_B^2\{4\widehat{m}_{\mathrm{}}^2[20\widehat{r}_\rho \lambda 12\widehat{r}_\rho (1\widehat{r}_\rho )^2\lambda \widehat{s}]`$
$`+`$ $`\widehat{s}[4\widehat{r}_\rho \lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )^2+\lambda \widehat{s}]\left\}\right|C_T\left|^2\right|t_1|^2`$
$`+`$ $`512m_B^2\{\widehat{s}[4\widehat{r}_\rho \lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )^2+\lambda \widehat{s}]`$
$`+`$ $`8\widehat{m}_{\mathrm{}}^2[12\widehat{r}_\rho (1\widehat{r}_\rho )^2+\lambda (\widehat{s}8\widehat{r}_\rho )]\left\}\right|C_{TE}\left|^2\right|t_1|^2`$
$``$ $`64m_B^2\widehat{s}^2\left[v^2\right|C_T|^2+4(32v^2)\left|C_{TE}|^2\right]\{2[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )]B_6t_1^{}`$
$``$ $`m_B^2\lambda (1+3\widehat{r}_\rho \widehat{s})B_7t_1^{}\}`$
$`+`$ $`768m_B^3\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\lambda (A_1+C_1)C_T^{}t_1^{}`$
$``$ $`192m_B\widehat{m}_{\mathrm{}}\widehat{s}[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )](B_1+D_1)C_{TE}^{}t_1^{}`$
$`+`$ $`192m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1+3\widehat{r}_\rho \widehat{s})\lambda (B_2+D_2)C_{TE}^{}t_1^{}\},`$
where $`\widehat{s}=q^2/m_B^2`$, $`\widehat{r}_\rho =m_\rho ^2/m_B^2`$ and $`\lambda (a,b,c)=a^2+b^2+c^22ab2ac2bc`$, $`\widehat{m}_{\mathrm{}}=m_{\mathrm{}}/m_B`$, $`v=\sqrt{14\widehat{m}_{\mathrm{}}^2/\widehat{s}}`$ is the final lepton velocity.
Using the matrix element for the $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay, our next problem is to calculate the nine doubleโlepton polarization asymmetries. For this aim we introduce the spin projection operators
$`\mathrm{\Lambda }_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\gamma _5\overline{)}s_i^{}),`$
$`\mathrm{\Lambda }_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\gamma _5\overline{)}s_i^+)`$
for the lepton $`\mathrm{}^{}`$ and antiโlepton $`\mathrm{}^+`$, where $`i=L,N,T`$ correspond to the longitudinal, normal and transversal polarizations, respectively. Firstly we define the following orthogonal unit vectors $`s_i^{\pm \mu }`$ in the rest frame of $`\mathrm{}^\pm `$ (see also ),
$`s_L^\mu `$ $`=`$ $`(0,\stackrel{}{e}_L^{})=(0,{\displaystyle \frac{\stackrel{}{p}_{}}{\left|\stackrel{}{p}_{}\right|}}),`$
$`s_N^\mu `$ $`=`$ $`(0,\stackrel{}{e}_N^{})=(0,{\displaystyle \frac{\stackrel{}{p}_\rho \times \stackrel{}{p}_{}}{\left|\stackrel{}{p}_\rho \times \stackrel{}{p}_{}\right|}}),`$
$`s_T^\mu `$ $`=`$ $`(0,\stackrel{}{e}_T^{})=(0,\stackrel{}{e}_N^{}\times \stackrel{}{e}_L^{}),`$
$`s_L^{+\mu }`$ $`=`$ $`(0,\stackrel{}{e}_L^+)=(0,{\displaystyle \frac{\stackrel{}{p}_+}{\left|\stackrel{}{p}_+\right|}}),`$
$`s_N^{+\mu }`$ $`=`$ $`(0,\stackrel{}{e}_N^+)=(0,{\displaystyle \frac{\stackrel{}{p}_\rho \times \stackrel{}{p}_+}{\left|\stackrel{}{p}_\rho \times \stackrel{}{p}_+\right|}}),`$
$`s_T^{+\mu }`$ $`=`$ $`(0,\stackrel{}{e}_T^+)=(0,\stackrel{}{e}_N^+\times \stackrel{}{e}_L^+),`$ (16)
where $`\stackrel{}{p}_{}`$ and $`\stackrel{}{p}_\rho `$ are the threeโmomenta of the leptons $`\mathrm{}^{}`$ and $`\rho `$ meson in the center of mass frame (CM) of $`\mathrm{}^{}\mathrm{}^+`$ system, respectively. Transformation of unit vectors from the rest frame of the leptons to CM frame of leptons can be done by the Lorentz boost. Boosting of the longitudinal unit vectors $`s_L^{\pm \mu }`$ leads to
$`\left(s_L^\mu \right)_{CM}`$ $`=`$ $`({\displaystyle \frac{\left|\stackrel{}{p}_{}\right|}{m_{\mathrm{}}}},{\displaystyle \frac{E_{\mathrm{}}\stackrel{}{p}_{}}{m_{\mathrm{}}\left|\stackrel{}{p}_{}\right|}}),`$ (17)
where $`\stackrel{}{p}_+=\stackrel{}{p}_{}`$, $`E_{\mathrm{}}`$ and $`m_{\mathrm{}}`$ are the energy and mass of leptons in the CM frame, respectively. The remaining two unit vectors $`s_N^{\pm \mu }`$, $`s_T^{\pm \mu }`$ are unchanged under Lorentz boost.
We can now define the doubleโlepton polarization asymmetries as in :
$`P_{ij}(\widehat{s})={\displaystyle \frac{\left({\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+){\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+)\right)\left({\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+){\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+)\right)}{\left({\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+)+{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+)\right)+\left({\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+)+{\displaystyle \frac{d\mathrm{\Gamma }}{d\widehat{s}}}(\stackrel{}{s}_i^{},\stackrel{}{s}_j^+)\right)}},`$ (18)
where $`i,j=L,N,T`$, and the first subindex $`i`$ corresponds lepton while the second subindex $`j`$ corresponds to antilepton, respectively.
After lengthy calculations we get the following results for the doubleโpolarization asymmetries.
$`P_{LL}`$ $`=`$ $`{\displaystyle \frac{m_B^2}{3\widehat{r}_\rho \widehat{s}\mathrm{\Delta }}}\text{Re}\{12m_B\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_1D_1)(B_4^{}B_5^{})`$ (19)
$``$ $`24m_B^2\widehat{m}_{\mathrm{}}^2\widehat{s}\lambda \left[B_4B_5^{}+(B_1D_1)(B_3^{}D_3^{})\right]`$
$`+`$ $`12m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1\widehat{r}_\rho )\left[(B_2D_2)(B_4^{}B_5^{})+2m_B\widehat{m}_{\mathrm{}}(B_2D_2)(B_3^{}D_3^{})\right]`$
$``$ $`32m_B^5\widehat{m}_{\mathrm{}}\widehat{s}\lambda ^2(B_2+D_2)B_7^{}C_{TE}^{}`$
$`+`$ $`3m_B^2\widehat{s}^2\lambda (1+v^2)(\left|B_4\right|^2+\left|B_5\right|^2)`$
$``$ $`8m_B^4\widehat{r}_\rho \widehat{s}^2\lambda (1+3v^2)(\left|A_1\right|^2+\left|C_1\right|^2)`$
$`+`$ $`12m_B^3\widehat{m}_{\mathrm{}}\widehat{s}^2\lambda (B_3D_3)(B_4^{}B_5^{})`$
$`+`$ $`12m_B^4\widehat{m}_{\mathrm{}}^2\widehat{s}^2\lambda \left|B_3D_3\right|^2`$
$`+`$ $`32m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1\widehat{r}_\rho \widehat{s})\left[(B_1+D_1)B_7^{}C_{TE}^{}+2(B_2+D_2)B_6^{}C_{TE}^{}\right]`$
$`+`$ $`8m_B^2\widehat{m}_{\mathrm{}}^2\lambda (44\widehat{r}_\rho \widehat{s})(B_1D_2^{}+B_2D_1^{})`$
$``$ $`64m_B\widehat{m}_{\mathrm{}}\widehat{s}(\lambda +12\widehat{r}_\rho \widehat{s})(B_1+D_1)B_6^{}C_{TE}^{}`$
$``$ $`16m_B^2\widehat{s}\left[4\widehat{m}_{\mathrm{}}^2(\left|C_T\right|^2+8\left|C_{TE}\right|^2)\widehat{s}(\left|C_T\right|^2+4\left|C_{TE}\right|^2)\right]\left[m_B^4\lambda ^2\left|B_7\right|^2+4(\lambda +12\widehat{r}_\rho \widehat{s})\left|B_6\right|^2\right]`$
$``$ $`32\widehat{m}_{\mathrm{}}^2(\lambda +3\widehat{r}_\rho \widehat{s})B_1D_1^{}`$
$``$ $`8m_B^4\widehat{m}_{\mathrm{}}^2\lambda [\lambda +3(1\widehat{r}_\rho )^2]B_2D_2^{}`$
$`+`$ $`8m_B^2\lambda [\widehat{s}\widehat{s}(\widehat{r}_\rho +\widehat{s})3\widehat{m}_{\mathrm{}}^2(22\widehat{r}_\rho \widehat{s})](B_1B_2^{}+D_1D_2^{})`$
$``$ $`64m_B^4\widehat{s}^2\lambda (1\widehat{r}_\rho \widehat{s})\left[v^2\left|C_T\right|^24(12v^2)\left|C_{TE}\right|^2\right]B_6B_7^{}`$
$``$ $`m_B^4\widehat{s}\lambda [\lambda (1+3v^2)3(1\widehat{r}_\rho )^2(1v^2)](\left|B_2\right|^2+\left|D_2\right|^2)`$
$`+`$ $`4[6\widehat{m}_{\mathrm{}}^2(\lambda +6\widehat{r}_\rho \widehat{s})\widehat{s}(\lambda +12\widehat{r}_\rho \widehat{s})](\left|B_1\right|^2+\left|D_1\right|^2)`$
$``$ $`1024m_B^2\{12\widehat{r}_\rho \widehat{s}(1\widehat{r}_\rho )^2(12v^2)\lambda \widehat{s}[4\widehat{r}_\rho \widehat{s}(12v^2)]\}\left|t_1\right|^2\left|C_{TE}\right|^2`$
$`+`$ $`256m_B^2\widehat{s}\{\lambda (\widehat{s}v^28\widehat{r}_\rho )+12\widehat{r}_\rho v^2[\lambda +(1\widehat{r}_\rho )^2]\}\left|t_1\right|^2\left|C_T\right|^2`$
$``$ $`256m_B^2\widehat{s}^2[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )]\left[v^2\left|C_T\right|^24(12v^2)\left|C_{TE}\right|^2\right]B_6t_1^{}`$
$`+`$ $`128m_B^4\widehat{s}^2\lambda (1+3\widehat{r}_\rho \widehat{s})\left[v^2\left|C_T\right|^24(12v^2)\left|C_{TE}\right|^2\right]B_7t_1^{}`$
$`+`$ $`128m_B\widehat{m}_{\mathrm{}}\widehat{s}[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )](B_1+D_1)t_1^{}C_{TE}^{}`$
$``$ $`128m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1+3\widehat{r}_\rho \widehat{s})(B_2+D_2)t_1^{}C_{TE}^{}`$
$``$ $`512m_B^3\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\lambda (A_1+C_1)t_1^{}C_T^{}`$
$``$ $`64m_B^4\widehat{m}_{\mathrm{}}^2\widehat{r}_\rho \widehat{s}\lambda A_1C_1^{}\},`$
$`P_{LN}`$ $`=`$ $`{\displaystyle \frac{\pi m_B^2}{2\widehat{r}_\rho \mathrm{\Delta }}}\sqrt{{\displaystyle \frac{\lambda }{\widehat{s}}}}\text{Im}\{4m_B^2\widehat{m}_{\mathrm{}}^2\lambda [B_2B_4^{}+B_5D_2^{}+8(B_1D_1)B_7^{}C_{TE}^{}]`$ (20)
$``$ $`4m_B^4\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )\left[B_2D_2^{}+8m_B\widehat{m}_{\mathrm{}}(B_2D_2)B_7^{}C_{TE}^{}\right]`$
$`+`$ $`2m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda \left[B_2B_3^{}8(B_4B_5)B_7^{}C_{TE}^{}16m_B\widehat{m}_{\mathrm{}}(B_3D_3)B_7^{}C_{TE}^{}\right]`$
$``$ $`2m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda \left[B_3D_2^{}+(B_2+D_2)D_3^{}\right]`$
$``$ $`2m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1+3\widehat{r}_\rho \widehat{s})\left(B_1B_2^{}D_1D_2^{}32B_5C_{TE}^{}t_1^{}\right)`$
$`+`$ $`32m_B^3\widehat{r}_\rho \widehat{s}^2v^2\left[(A_1C_1)B_6^{}C_T^{}2(A_1+C_1)B_6^{}C_{TE}^{}\right]`$
$``$ $`m_B^3\widehat{s}\lambda (1+v^2)\left(B_2B_5^{}+B_4D_2^{}\right)`$
$``$ $`4\widehat{m}_{\mathrm{}}(1\widehat{r}_\rho \widehat{s})\left\{B_1D_1^{}+m_B\widehat{m}_{\mathrm{}}\left[B_1(B_4^{}+16B_6^{}C_{TE}^{})D_1(B_5^{}+16B_6^{}C_{TE}^{})\right]\right\}`$
$`+`$ $`64m_B^3\widehat{m}_{\mathrm{}}^2(1\widehat{r}_\rho )(1\widehat{r}_\rho \widehat{s})(B_2D_2)B_6^{}C_{TE}^{}`$
$``$ $`2m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1\widehat{r}_\rho \widehat{s})(B_1+D_1)(B_3^{}D_3^{})`$
$`+`$ $`32m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1\widehat{r}_\rho \widehat{s})\left[(B_4B_5)B_6^{}C_{TE}^{}+2m_B\widehat{m}_{\mathrm{}}(B_3D_3)B_6^{}C_{TE}^{}\right]`$
$`+`$ $`m_B\widehat{s}(1\widehat{r}_\rho \widehat{s})(1+v^2)\left(B_1B_5^{}+B_4D_1^{}\right)`$
$`+`$ $`2m_B^2\widehat{m}_{\mathrm{}}[\lambda +(1\widehat{r}_\rho )(1\widehat{r}_\rho \widehat{s})]\left(B_2D_1^{}+B_1D_2^{}\right)`$
$``$ $`128m_B^3\widehat{m}_{\mathrm{}}^2(1\widehat{r}_\rho )(1+3\widehat{r}_\rho \widehat{s})(B_2D_2)t_1^{}C_{TE}^{}`$
$``$ $`64m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1+3\widehat{r}_\rho \widehat{s})\left[B_4C_{TE}^{}t_1^{}+2m_B\widehat{m}_{\mathrm{}}(B_3D_3)C_{TE}^{}t_1^{}\right]`$
$``$ $`64m_B^3\widehat{r}_\rho \widehat{s}(1\widehat{r}_\rho )v^2\left[A_1(C_T^{}2C_{TE}^{})t_1^{}C_1(C_T^{}+2C_{TE}^{})t_1^{}\right]`$
$``$ $`32m_B\widehat{s}\left[(1+3\widehat{r}_\rho \widehat{s})D_1C_{TE}^{}t_1^{}+2\widehat{r}_\rho v^2D_1C_T^{}t_1^{}(1\widehat{r}_\rho \widehat{s})v^2D_1C_{TE}^{}t_1^{}\right]`$
$`+`$ $`32m_B\widehat{s}[(1+3\widehat{r}_\rho \widehat{s})B_1C_{TE}^{}t_1^{}2\widehat{r}_\rho v^2B_1C_T^{}t_1^{}(1\widehat{r}_\rho \widehat{s})v^2B_1C_{TE}^{}t_1^{}]\},`$
$`P_{NL}`$ $`=`$ $`{\displaystyle \frac{\pi m_B^2}{2\widehat{r}_\rho \mathrm{\Delta }}}\sqrt{{\displaystyle \frac{\lambda }{\widehat{s}}}}\text{Im}\{4m_B^3\widehat{m}_{\mathrm{}}^2\lambda [B_2B_5^{}+B_4D_2^{}8(B_1D_1)B_7^{}C_{TE}^{}]`$ (21)
$`+`$ $`4m_B^4\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )\left[B_2D_2^{}+8m_B\widehat{m}_{\mathrm{}}(B_2D_2)B_7^{}C_{TE}^{}\right]`$
$``$ $`2m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda \left[B_2B_3^{}8(B_4B_5)B_7^{}C_{TE}^{}16m_B\widehat{m}_{\mathrm{}}(B_3D_3)B_7^{}C_{TE}^{}\right]`$
$`+`$ $`2m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda \left[B_3D_2^{}+(B_2+D_2)D_3^{}\right]`$
$`+`$ $`2m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1+3\widehat{r}_\rho \widehat{s})\left(B_1B_2^{}D_1D_2^{}32B_5C_{TE}^{}t_1^{}\right)`$
$``$ $`32m_B^3\widehat{r}_\rho \widehat{s}^2v^2\left[(A_1C_1)B_6^{}C_T^{}+2(A_1+C_1)B_6^{}C_{TE}^{}\right]`$
$``$ $`m_B^3\widehat{s}\lambda (1+v^2)\left(B_2B_4^{}+B_5D_2^{}\right)`$
$`+`$ $`4\widehat{m}_{\mathrm{}}(1\widehat{r}_\rho \widehat{s})\left\{B_1D_1^{}m_B\widehat{m}_{\mathrm{}}\left[B_1(B_5^{}16B_6^{}C_{TE}^{})D_1(B_4^{}16B_6^{}C_{TE}^{})\right]\right\}`$
$`+`$ $`64m_B^3\widehat{m}_{\mathrm{}}^2(1\widehat{r}_\rho )(1\widehat{r}_\rho \widehat{s})(B_2D_2)B_6^{}C_{TE}^{}`$
$`+`$ $`2m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1\widehat{r}_\rho \widehat{s})(B_1+D_1)(B_3^{}D_3^{})`$
$``$ $`32m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1\widehat{r}_\rho \widehat{s})\left[(B_4B_5)B_6^{}C_{TE}^{}+2m_B\widehat{m}_{\mathrm{}}(B_3D_3)B_6^{}C_{TE}^{}\right]`$
$`+`$ $`m_B\widehat{s}(1\widehat{r}_\rho \widehat{s})(1+v^2)\left(B_1B_4^{}+B_5D_1^{}\right)`$
$``$ $`2m_B^2\widehat{m}_{\mathrm{}}[\lambda +(1\widehat{r}_\rho )(1\widehat{r}_\rho \widehat{s})]\left(B_2D_1^{}+B_1D_2^{}\right)`$
$`+`$ $`128m_B^3\widehat{m}_{\mathrm{}}^2(1\widehat{r}_\rho )(1+3\widehat{r}_\rho \widehat{s})(B_2D_2)t_1^{}C_{TE}^{}`$
$`+`$ $`64m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1+3\widehat{r}_\rho \widehat{s})\left[B_4C_{TE}^{}t_1^{}+2m_B\widehat{m}_{\mathrm{}}(B_3D_3)C_{TE}^{}t_1^{}\right]`$
$`+`$ $`64m_B^3\widehat{r}_\rho \widehat{s}(1\widehat{r}_\rho )v^2\left[A_1(C_T^{}+2C_{TE}^{})t_1^{}C_1(C_T^{}2C_{TE}^{})t_1^{}\right]`$
$``$ $`32m_B\widehat{s}\left[(1+3\widehat{r}_\rho \widehat{s})B_1C_{TE}^{}t_1^{}+2\widehat{r}_\rho v^2B_1C_T^{}t_1^{}(1\widehat{r}_\rho \widehat{s})v^2B_1C_{TE}^{}t_1^{}\right]`$
$`+`$ $`32m_B\widehat{s}[(1+3\widehat{r}_\rho \widehat{s})D_1C_{TE}^{}t_1^{}2\widehat{r}_\rho v^2D_1C_T^{}t_1^{}(1\widehat{r}_\rho \widehat{s})v^2D_1C_{TE}^{}t_1^{}]\},`$
$`P_{LT}`$ $`=`$ $`{\displaystyle \frac{\pi m_B^2v}{\widehat{r}_\rho \mathrm{\Delta }}}\sqrt{{\displaystyle \frac{\lambda }{\widehat{s}}}}\text{Re}\left\{m_B^4\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )\right|B_2D_2|^2`$ (22)
$``$ $`8m_B^2\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\left(A_1B_1^{}C_1D_1^{}\right)`$
$``$ $`m_B^3\widehat{s}\lambda \left(B_2B_5^{}+B_4D_2^{}m_B\widehat{m}_{\mathrm{}}B_2B_3^{}\right)`$
$``$ $`8m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_4+B_5)B_7^{}C_{TE}^{}`$
$``$ $`m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda \left(B_2D_3^{}+B_3D_2^{}D_2D_3^{}\right)`$
$`+`$ $`16m_B^3\widehat{r}_\rho \widehat{s}^2\left[A_1B_6^{}(C_T^{}2C_{TE}^{})+C_1B_6^{}(C_T^{}+2C_{TE}^{})\right]`$
$`+`$ $`\widehat{m}_{\mathrm{}}(1\widehat{r}_\rho \widehat{s})\left|B_1D_1\right|^2`$
$`+`$ $`m_B\widehat{s}(1\widehat{r}_\rho \widehat{s})[B_1B_5^{}+B_4D_1^{}+16m_B\widehat{m}_{\mathrm{}}(B_4+B_5)B_6^{}C_{TE}^{}`$
$``$ $`m_B\widehat{m}_{\mathrm{}}(B_1D_1)(B_3^{}D_3^{})]`$
$``$ $`m_B^2\widehat{m}_{\mathrm{}}[\lambda +(1\widehat{r}_\rho )(1\widehat{r}_\rho \widehat{s})](B_1D_1)(B_2^{}D_2^{})`$
$``$ $`1024m_B^2\widehat{m}_{\mathrm{}}\widehat{r}_\rho (1\widehat{r}_\rho )\left(\left|C_T\right|^2+4\left|C_{TE}\right|^2\right)\left|t_1\right|^2`$
$`+`$ $`512m_B^2\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\left(\left|C_T\right|^2+4\left|C_{TE}\right|^2\right)B_6t_1^{}`$
$``$ $`32m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1+3\widehat{r}_\rho \widehat{s})(B_4+B_5)C_{TE}^{}t_1^{}`$
$``$ $`32m_B^3\widehat{r}_\rho \widehat{s}(1\widehat{r}_\rho )\left[A_1(C_T^{}2C_{TE}^{})t_1^{}+C_1(C_T^{}+2C_{TE}^{})t_1^{}\right]`$
$``$ $`32m_B\widehat{r}_\rho \widehat{s}[B_1(C_T^{}2C_{TE}^{})t_1^{}D_1(C_T^{}+2C_{TE}^{})t_1^{}]\},`$
$`P_{TL}`$ $`=`$ $`{\displaystyle \frac{\pi m_B^2v}{\widehat{r}_\rho \mathrm{\Delta }}}\sqrt{{\displaystyle \frac{\lambda }{\widehat{s}}}}\text{Re}\left\{m_B^4\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )\right|B_2D_2|^2`$ (23)
$`+`$ $`8m_B^2\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\left(A_1B_1^{}C_1D_1^{}\right)`$
$`+`$ $`m_B^3\widehat{s}\lambda \left(B_2B_4^{}+B_5D_2^{}+m_B\widehat{m}_{\mathrm{}}B_2B_3^{}\right)`$
$`+`$ $`8m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_4+B_5)B_7^{}C_{TE}^{}`$
$``$ $`m_B^4\widehat{m}_{\mathrm{}}\widehat{s}\lambda \left(B_2D_3^{}+B_3D_2^{}D_2D_3^{}\right)`$
$`+`$ $`16m_B^3\widehat{r}_\rho \widehat{s}^2\left[A_1B_6^{}(C_T^{}+2C_{TE}^{})+C_1B_6^{}(C_T^{}2C_{TE}^{})\right]`$
$`+`$ $`\widehat{m}_{\mathrm{}}(1\widehat{r}_\rho \widehat{s})\left|B_1D_1\right|^2`$
$``$ $`m_B\widehat{s}(1\widehat{r}_\rho \widehat{s})[B_1B_4^{}+B_5D_1^{}+16m_B\widehat{m}_{\mathrm{}}(B_4+B_5)B_6^{}C_{TE}^{}`$
$`+`$ $`m_B\widehat{m}_{\mathrm{}}(B_1D_1)(B_3^{}D_3^{})]`$
$``$ $`m_B^2\widehat{m}_{\mathrm{}}[\lambda +(1\widehat{r}_\rho )(1\widehat{r}_\rho \widehat{s})](B_1D_1)(B_2^{}D_2^{})`$
$``$ $`1024m_B^2\widehat{m}_{\mathrm{}}\widehat{r}_\rho (1\widehat{r}_\rho )\left(\left|C_T\right|^2+4\left|C_{TE}\right|^2\right)\left|t_1\right|^2`$
$`+`$ $`512m_B^2\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\left(\left|C_T\right|^2+4\left|C_{TE}\right|^2\right)B_6t_1^{}`$
$`+`$ $`32m_B^2\widehat{m}_{\mathrm{}}\widehat{s}(1+3\widehat{r}_\rho \widehat{s})(B_4+B_5)C_{TE}^{}t_1^{}`$
$``$ $`32m_B^3\widehat{r}_\rho \widehat{s}(1\widehat{r}_\rho )\left[A_1(C_T^{}+2C_{TE}^{})t_1^{}+C_1(C_T^{}2C_{TE}^{})t_1^{}\right]`$
$`+`$ $`32m_B\widehat{r}_\rho \widehat{s}[B_1(C_T^{}+2C_{TE}^{})t_1^{}D_1(C_T^{}2C_{TE}^{})t_1^{}]\},`$
$`P_{NT}`$ $`=`$ $`{\displaystyle \frac{2m_B^2v}{3\widehat{r}_\rho \mathrm{\Delta }}}\text{Im}\{4\lambda \{B_1D_1^{}+m_B^4\lambda [B_2D_2^{}2m_B\widehat{m}_{\mathrm{}}B_2B_7^{}(C_T^{}4C_{TE}^{})`$ (24)
$``$ $`2m_B\widehat{m}_{\mathrm{}}D_2B_7^{}(C_T^{}+4C_{TE}^{})]\}`$
$``$ $`6m_B\widehat{m}_{\mathrm{}}\lambda (B_1D_1)(B_4^{}+B_5^{})`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )(B_2D_2)(B_4^{}+B_5^{})`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_3D_3)(B_4^{}+B_5^{})`$
$``$ $`4m_B^2\lambda (1\widehat{r}_\rho \widehat{s})\left[B_1D_2^{}+B_2D_1^{}+32m_B^2\widehat{s}\text{Re}[B_6B_7^{}]C_TC_{TE}^{}\right]`$
$`+`$ $`8m_B^3\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho \widehat{s})[(B_1B_7^{}+2B_2B_6^{})(C_T^{}4C_{TE}^{})`$
$`+`$ $`(B_7^{}D_1+2B_6^{}D_2)(C_T^{}+4C_{TE}^{})]`$
$`+`$ $`32m_B^2\widehat{s}\left[4(\lambda +12\widehat{r}_\rho \widehat{s})\left|B_6\right|^2+\lambda ^2m_B^4\left|B_7\right|^2\right]C_TC_{TE}^{}`$
$`+`$ $`2m_B^2\widehat{s}\lambda \left(3B_4B_5^{}8m_B^2\widehat{r}_\rho A_1C_1^{}\right)`$
$``$ $`16m_B\widehat{m}_{\mathrm{}}\left\{\lambda \left[B_1B_6^{}(C_T^{}4C_{TE}^{})+D_1B_6^{}(C_T^{}+4C_{TE}^{})\right]+12\widehat{r}_\rho \widehat{s}(B_1+D_1)B_6^{}C_T^{}\right\}`$
$`+`$ $`32m_B\widehat{m}_{\mathrm{}}\left\{12\widehat{r}_\rho (1\widehat{r}_\rho )(B_1+D_1)C_T^{}t_1^{}+\lambda \left[B_1(C_T^{}4C_{TE}^{})t_1^{}+D_1(C_T^{}+4C_{TE}^{})t_1^{}\right]\right\}`$
$``$ $`256m_B^3\widehat{m}_{\mathrm{}}\widehat{r}_\rho \lambda \left[A_1(C_T^{}+2C_{TE}^{})t_1^{}C_1(C_T^{}2C_{TE}^{})t_1^{}\right]`$
$``$ $`32m_B^3\widehat{m}_{\mathrm{}}\lambda (1+3\widehat{r}_\rho \widehat{s})\left[B_2(C_T^{}4C_{TE}^{})t_1^{}+D_2(C_T^{}+4C_{TE}^{})t_1^{}\right]`$
$`+`$ $`256m_B^2\widehat{s}\left\{2[\lambda +12\widehat{r}_\rho (2+2\widehat{r}_\rho \widehat{s})]\right|t_1|^22[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )]\text{Re}[B_6t_1^{}]`$
$`+`$ $`m_B^2\lambda (1+3\widehat{r}_\rho \widehat{s})\text{Re}[B_7t_1^{}]\}C_TC_{TE}^{}\},`$
$`P_{TN}`$ $`=`$ $`{\displaystyle \frac{2m_B^2v}{3\widehat{r}_\rho \mathrm{\Delta }}}\text{Im}\{4\lambda \{B_1D_1^{}+m_B^4\lambda [B_2D_2^{}+2m_B\widehat{m}_{\mathrm{}}B_2B_7^{}(C_T^{}+4C_{TE}^{})`$ (25)
$`+`$ $`2m_B\widehat{m}_{\mathrm{}}D_2B_7^{}(C_T^{}4C_{TE}^{})]\}`$
$``$ $`6m_B\widehat{m}_{\mathrm{}}\lambda (B_1D_1)(B_4^{}+B_5^{})`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )(B_2D_2)(B_4^{}+B_5^{})`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_3D_3)(B_4^{}+B_5^{})`$
$`+`$ $`4m_B^2\lambda (1\widehat{r}_\rho \widehat{s})\left[B_1D_2^{}+B_2D_1^{}32m_B^2\widehat{s}\text{Re}[B_6B_7^{}]C_TC_{TE}^{}\right]`$
$`+`$ $`8m_B^3\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho \widehat{s})[(B_1B_7^{}+2B_2B_6^{})(C_T^{}+4C_{TE}^{})`$
$`+`$ $`(B_7^{}D_1+2B_6^{}D_2)(C_T^{}4C_{TE}^{})]`$
$`+`$ $`32m_B^2\widehat{s}\left[4(\lambda +12\widehat{r}_\rho \widehat{s})\left|B_6\right|^2+\lambda ^2m_B^4\left|B_7\right|^2\right]C_TC_{TE}^{}`$
$`+`$ $`2m_B^2\widehat{s}\lambda \left(3B_4B_5^{}+8m_B^2\widehat{r}_\rho A_1C_1^{}\right)`$
$``$ $`16m_B\widehat{m}_{\mathrm{}}\left\{\lambda \left[B_1B_6^{}(C_T^{}+4C_{TE}^{})+D_1B_6^{}(C_T^{}4C_{TE}^{})\right]+12\widehat{r}_\rho \widehat{s}(B_1+D_1)B_6^{}C_T^{}\right\}`$
$`+`$ $`32m_B\widehat{m}_{\mathrm{}}\left\{12\widehat{r}_\rho (1\widehat{r}_\rho )(B_1+D_1)C_T^{}t_1^{}+\lambda \left[B_1(C_T^{}+4C_{TE}^{})t_1^{}+D_1(C_T^{}4C_{TE}^{})t_1^{}\right]\right\}`$
$`+`$ $`256m_B^3\widehat{m}_{\mathrm{}}\widehat{r}_\rho \lambda \left[A_1(C_T^{}2C_{TE}^{})t_1^{}C_1(C_T^{}+2C_{TE}^{})t_1^{}\right]`$
$``$ $`32m_B^3\widehat{m}_{\mathrm{}}\lambda (1+3\widehat{r}_\rho \widehat{s})\left[B_2(C_T^{}+4C_{TE}^{})t_1^{}+D_2(C_T^{}4C_{TE}^{})t_1^{}\right]`$
$`+`$ $`256m_B^2\widehat{s}\left\{2[\lambda +12\widehat{r}_\rho (2+2\widehat{r}_\rho \widehat{s})]\right|t_1|^22[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )]\text{Re}[B_6t_1^{}]`$
$`+`$ $`m_B^2\lambda (1+3\widehat{r}_\rho \widehat{s})\text{Re}[B_7t_1^{}]\}C_TC_{TE}^{}\},`$
$`P_{NN}`$ $`=`$ $`{\displaystyle \frac{2m_B^2}{3\widehat{r}_\rho \mathrm{\Delta }}}\text{Re}\{24\widehat{m}_{\mathrm{}}^2\widehat{r}_\rho (|B_1|^2+|D_1|^2)`$ (26)
$``$ $`6m_B\widehat{m}_{\mathrm{}}\lambda (B_1D_1)(B_4^{}B_5^{})`$
$``$ $`48m_B^5\widehat{m}_{\mathrm{}}\lambda ^2(B_2+D_2)B_7^{}C_{TE}^{}`$
$`+`$ $`6m_B^2\widehat{m}_{\mathrm{}}^2\lambda \left[\left|B_4\right|^2+\left|B_5\right|^22B_1(B_2^{}+B_3^{}D_3^{})+2D_1(B_3^{}D_2^{}D_3^{})\right]`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho )\left[(B_2D_2)(B_4^{}B_5^{})+2m_B\widehat{m}_{\mathrm{}}(B_2D_2)(B_3^{}D_3^{})\right]`$
$`+`$ $`m_B^2\widehat{s}\lambda \left[16m_B^2\widehat{r}_\rho v^2A_1C_1^{}3(1+v^2)B_4B_5^{}\right]`$
$`+`$ $`6m_B^4\widehat{m}_{\mathrm{}}^2\lambda (2+2\widehat{r}_\rho \widehat{s})(\left|B_2\right|^2+\left|D_2\right|^2)`$
$`+`$ $`6m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_3D_3)(B_4^{}B_5^{})`$
$`+`$ $`6m_B^4\widehat{m}_{\mathrm{}}^2\widehat{s}\lambda \left|B_3D_3\right|^2`$
$`+`$ $`48m_B^3\widehat{m}_{\mathrm{}}\lambda (1\widehat{r}_\rho \widehat{s})\left[(B_1+D_1)B_7^{}C_{TE}^{}+2(B_2+D_2)B_6^{}C_{TE}^{}\right]`$
$``$ $`96m_B\widehat{m}_{\mathrm{}}(1\widehat{r}_\rho \widehat{s})^2(B_1+D_1)B_6^{}C_{TE}^{}`$
$`+`$ $`8m_B^4\widehat{s}\lambda \left[\lambda m_B^2\left|B_7\right|^24(1\widehat{r}_\rho \widehat{s})B_6B_7^{}\right]\left[v^2\left|C_T\right|^24(32v^2)\left|C_{TE}\right|^2\right]`$
$`+`$ $`m_B^2\lambda [3(22\widehat{r}_\rho \widehat{s})v^2(22\widehat{r}_\rho +\widehat{s})](B_1D_2^{}+B_2D_1^{})`$
$``$ $`m_B^4\lambda \left[(3+v^2)\lambda +3(1v^2)(1\widehat{r}_\rho )^2\right]B_2D_2^{}`$
$``$ $`2[6\widehat{r}_\rho \widehat{s}(1v^2)+\lambda (3v^2)]B_1D_1^{}`$
$`+`$ $`32m_B^2\widehat{s}\left\{(\lambda +12\widehat{r}_\rho \widehat{s})v^2\left|C_T\right|^24[\lambda (32v^2)+12\widehat{r}_\rho \widehat{s}]\left|C_{TE}\right|^2\right\}\left|B_6\right|^2`$
$``$ $`192m_B^3\widehat{m}_{\mathrm{}}\lambda (1+3\widehat{r}_\rho \widehat{s})(B_2+D_2)C_{TE}^{}t_1^{}`$
$`+`$ $`192m_B\widehat{m}_{\mathrm{}}[\lambda +4\widehat{r}_\rho (1\widehat{r}_\rho )](B_1+D_1)C_{TE}^{}t_1^{}`$
$`+`$ $`128m_B^2\widehat{s}v^2[\lambda +12\widehat{r}_\rho (2+2\widehat{r}_\rho \widehat{s})]\left|C_T\right|^2\left|t_1\right|^2`$
$``$ $`512m_B^2\widehat{s}[\lambda (32v^2)+12\widehat{r}_\rho (2+2\widehat{r}_\rho \widehat{s})]\left|C_{TE}\right|^2\left|t_1\right|^2`$
$``$ $`128m_B^2\widehat{s}\left\{[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )]v^2\left|C_T\right|^24[\lambda (32v^2)+12\widehat{r}_\rho (1\widehat{r}_\rho )]\left|C_{TE}\right|^2\right\}B_6t_1^{}`$
$`+`$ $`64m_B^4(1+3\widehat{r}_\rho \widehat{s})\widehat{s}\left[v^2\right|C_T|^24(32v^2)\left|C_{TE}|^2\right]B_7t_1^{}\},`$
$`P_{TT}`$ $`=`$ $`{\displaystyle \frac{2m_B^2}{3\widehat{r}_\rho \widehat{s}\mathrm{\Delta }}}\text{Re}\{8m_B^4\widehat{r}_\rho \widehat{s}\lambda [4\widehat{m}_{\mathrm{}}^2(|A_1|^2+|C_1|^2)+2\widehat{s}A_1C_1^{}]`$ (27)
$`+`$ $`6m_B\widehat{m}_{\mathrm{}}\widehat{s}\lambda (B_1D_1)(B_4^{}B_5^{})`$
$``$ $`16m_B^5\widehat{m}_{\mathrm{}}\widehat{s}\lambda ^2(B_2+D_2)B_7^{}C_{TE}^{}`$
$``$ $`6m_B^2\widehat{m}_{\mathrm{}}^2\widehat{s}\lambda \left[\left|B_4\right|^2+\left|B_5\right|^22(B_1D_1)(B_3^{}D_3^{})\right]`$
$``$ $`6m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1\widehat{r}_\rho )\left[(B_2D_2)(B_4^{}B_5^{})+2m_B\widehat{m}_{\mathrm{}}(B_2D_2)(B_3^{}D_3^{})\right]`$
$``$ $`6m_B^3\widehat{m}_{\mathrm{}}\widehat{s}^2\lambda (B_3D_3)(B_4^{}B_5^{})`$
$``$ $`6m_B^4\widehat{m}_{\mathrm{}}^2\widehat{s}^2\lambda \left|B_3D_3\right|^2`$
$`+`$ $`16m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1\widehat{r}_\rho \widehat{s})\left[(B_1+D_1)B_7^{}C_{TE}^{}+2(B_2+D_2)B_6^{}C_{TE}^{}\right]`$
$`+`$ $`4m_B^2\widehat{m}_{\mathrm{}}^2\lambda (44\widehat{r}_\rho \widehat{s})(B_1B_2^{}+D_1D_2^{})`$
$`+`$ $`2\widehat{s}[6\widehat{r}_\rho \widehat{s}(1v^2)+\lambda (13v^2)]B_1D_1^{}`$
$``$ $`2m_B^4\widehat{m}_{\mathrm{}}^2\lambda [\lambda +3(1\widehat{r}_\rho )^2](\left|B_2\right|^2+\left|D_2\right|^2)`$
$``$ $`m_B^2\widehat{s}\lambda [22\widehat{r}_\rho +\widehat{s}3v^2(22\widehat{r}_\rho \widehat{s})](B_1D_2^{}+B_2D_1^{})`$
$``$ $`8\widehat{m}_{\mathrm{}}^2(\lambda 3\widehat{r}_\rho \widehat{s})(\left|B_1\right|^2+\left|D_1\right|^2)`$
$``$ $`32m_B\widehat{m}_{\mathrm{}}\widehat{s}(\lambda 12\widehat{r}_\rho \widehat{s})(B_1+D_1)B_6^{}C_{TE}^{}`$
$``$ $`8m_B^4\widehat{s}^2\lambda \left[\lambda m_B^2\left|B_7\right|^24(1\widehat{r}_\rho \widehat{s})B_6B_7^{}\right]\left[v^2\left|C_T\right|^2+4(12v^2)\left|C_{TE}\right|^2\right]`$
$`+`$ $`3m_B^2\widehat{s}^2\lambda (1+v^2)B_4B_5^{}`$
$``$ $`m_B^4\widehat{s}\lambda \left[(1+3v^2)\lambda 3(1v^2)(1\widehat{r}_\rho )^2\right]B_2D_2^{}`$
$``$ $`32m_B^2\widehat{s}^2\left\{(\lambda +12\widehat{r}_\rho \widehat{s})v^2\left|C_T\right|^2+4[\lambda (12v^2)12\widehat{r}_\rho \widehat{s}]\left|C_{TE}\right|^2\right\}\left|B_6\right|^2`$
$``$ $`128m_B^2\left\{4\lambda [\lambda (1\widehat{r}_\rho )^2]+8\widehat{s}(1\widehat{r}_\rho )[\lambda 6\widehat{r}_\rho (1\widehat{r}_\rho )]+8\lambda \widehat{s}v^2(8\widehat{r}_\rho \widehat{s})\right\}\left|C_{TE}\right|^2\left|t_1\right|^2`$
$`+`$ $`128m_B^2\left\{16\lambda \widehat{r}_\rho \widehat{s}\lambda [\lambda (1\widehat{r}_\rho )^2]v^22\widehat{s}v^2[\lambda (1+3\widehat{r}_\rho )+6\widehat{r}_\rho (1\widehat{r}_\rho )^2]\right\}\left|C_T\right|^2\left|t_1\right|^2`$
$`+`$ $`128m_B^2\widehat{s}^2\left\{[\lambda +12\widehat{r}_\rho (1\widehat{r}_\rho )]v^2\left|C_T\right|^2+4[\lambda (12v^2)12\widehat{r}_\rho (1\widehat{r}_\rho )]\right\}B_6t_1^{}`$
$``$ $`64m_B^4\widehat{s}^2\lambda (1+3\widehat{r}_\rho \widehat{s})\left[v^2\left|C_T\right|^2+4(12v^2)\left|C_{TE}\right|^2\right]B_7t_1^{}`$
$`+`$ $`64m_B\widehat{m}_{\mathrm{}}\widehat{s}[\lambda 12\widehat{r}_\rho (1\widehat{r}_\rho )](B_1+D_1)C_{TE}^{}t_1^{}`$
$`+`$ $`512m_B^3\widehat{m}_{\mathrm{}}\widehat{r}_\rho \widehat{s}\lambda (A_1+C_1)C_T^{}t_1^{}`$
$``$ $`64m_B^3\widehat{m}_{\mathrm{}}\widehat{s}\lambda (1+3\widehat{r}_\rho \widehat{s})(B_2+D_2)C_{TE}^{}t_1^{}\},`$
## 3 Numerical analysis
In this section we analyze the effects of the Wilson coefficients on the polarized $`FB`$ asymmetry. The input parameters we use in our numerical calculations are: $`m_\rho =0.77GeV`$, $`m_\tau =1.77GeV`$, $`m_\mu =0.106GeV`$, $`m_b=4.8GeV`$, $`m_B=5.26GeV`$ and $`\mathrm{\Gamma }_B=4.22\times 10^{13}GeV`$. For the values of the Wilson coefficients we use $`C_7^{SM}=0.313,C_9^{SM}=4.344`$ and $`C_{10}^{SM}=4.669`$. It should be noted that the aboveโpresented value for $`C_9^{SM}`$ corresponds only to short distance contributions. In addition to the short distance contributions, it receives long distance contributions which result from the conversion of $`\overline{u}u`$, $`\overline{d}d`$ and $`\overline{c}c`$ to the lepton pair. In order to minimize the hadronic uncertainties we will discard the regions around low lying resonances $`\rho `$, $`w`$, $`J/\psi `$, $`\psi ^{}`$, $`\psi ^{\prime \prime }`$, by dividing the $`q^2`$ region to low and high dilepton mass intervals:
$`\begin{array}{cc}\text{Region I:}\hfill & 1GeV^2q^28GeV^2,\hfill \\ \text{Region II:}\hfill & 14.5GeV^2q^2(m_Bm_\rho )^2GeV^2,\hfill \end{array}`$ (30)
where the contributions of the higher $`\psi `$ resonances do still exist in the second region. The form factors we have used in the present work are more refined ones predicted by the light cone QCD sum rules . The $`q^2`$ dependence of the form factors for the $`B\rho `$ transition can be represented in the following form:
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{r_1}{1q^2/m_{\mathrm{res}}^2}}+{\displaystyle \frac{r_2}{1q^2/m_{\mathrm{fit}}^2}},`$ (31)
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{r_2}{1q^2/m_{\mathrm{fit}}^2}},`$ (32)
$`F(q^2)`$ $`=`$ $`{\displaystyle \frac{r_1}{1q^2/m_{\mathrm{fit}}^2}}+{\displaystyle \frac{r_2}{(1q^2/m_{\mathrm{fit}}^2)^2}},`$ (33)
with the three independent parameters $`r_1`$, $`r_2`$ and $`m_{\mathrm{fit}}`$ being listed listed in Table 1. The dominant poles at $`q^2=m_{\mathrm{res}}^2`$ correspond to the resonances
$`J^P=\{\begin{array}{cc}1^{}\hfill & \text{for }V,\hfill \\ 0^{}\hfill & \text{for }A_0,\hfill \\ 1^+\hfill & \text{for }A_1,A_2,A_3\text{ and }T_2,T_3.\hfill \end{array}`$ (37)
The values of the parameters $`r_1`$, $`r_2`$ and $`m_{\mathrm{fit}}`$ for various form factors are presented in Table-1.
Note that $`T_3`$ entering into Eqs. (2) and (2) is related to $`\stackrel{~}{T}_3`$ as follows::
$`T_3={\displaystyle \frac{m_B^2m_\rho ^2}{q^2}}(\stackrel{~}{T}_3T_2).`$
In the numerical analysis the values of the new Wilson coefficients which describe the new physics beyond the SM are needed. In our calculations the new Wilson coefficients are varied in the range $`|C_{10}^{SM}||C_X||C_{10}^{SM}|`$. The experimental results on the branching ratio of the $`BK^{}(K)\mathrm{}^+\mathrm{}^{}`$ decay and the upper limit on the branching ratio of $`B\mu ^+\mu ^{}`$ suggests that that this is the right order of magnitude for the new Wilson coefficients.
It follows from the expressions of all nine doubleโlepton polarization asymmetries that depend both on $`q^2`$ and the new Wilson coefficients $`C_X`$. Therefore, it may experimentally be difficult to study these dependencies at the same time. For this reason, we eliminate $`q^2`$ dependence by performing integration over $`q^2`$ in the allowed region, i.e., we consider the averaged doubleโlepton polarization asymmetries. The averaging over $`q^2`$ is defined as
$`P_{ij}={\displaystyle \frac{{\displaystyle _{R_i}}P_{ij}{\displaystyle \frac{d}{d\widehat{s}}}๐\widehat{s}}{{\displaystyle _{R_i}}{\displaystyle \frac{d}{d\widehat{s}}}๐\widehat{s}}},`$
where $`R_i=`$ Regions I or II, over which the integrations are calculated. We present our analysis in a series of figures.
In Figs. (1) and (2) we present the dependence of $`P_{LL}`$ on $`C_X`$ for the $`B\rho \mu ^+\mu ^{}`$ decay in the regions I and II, respectively. The intersection of all curves corresponds to the SM case. From these figures we see that $`P_{LL}`$ exhibits strong dependence only on the tensor interactions $`C_T`$ and $`C_{TE}`$, and has practically symmetric behavior in regard to its dependence on $`C_T`$ and $`C_{TE}`$ with respect to zero position. Furthermore, $`P_{LL}`$ seems to be independent of all remaining new Wilson coefficients.
We depict from Figs. (3) and (4) the dependence of $`P_{LT}`$ on $`C_X`$ for the $`B\rho \mu ^+\mu ^{}`$ decay in the regions I and II, respectively.We observe from these figures that $`P_{LT}`$ is sensitive to to the existence of scalar $`C_{LRLR},C_{RLLR}`$ and tensor interactions $`C_T,C_{TE}`$ and it shows weak dependence on all remaining coefficients. A striking feature of its behavior is that $`P_{LT}`$ changes its sign in the aboveโmentioned region of the new Wilson coefficients, while in the SM case its sign never changes. For this reason study of the magnitude and sign of $`P_{LT}`$ can serve as a good test for looking new physics beyond the SM.
The dependence of $`P_{TL}`$ on $`C_X`$ for the $`B\rho \mu ^+\mu ^{}`$ decay is presented in Fig. (5) in Region I and Fig. (6) in Region II, respectively. In both regions $`P_{TL}`$ exhibits strong dependence on scalar $`C_{RLRL}`$ and $`C_{LRRL}`$ and tensor interaction coefficients. Moreover, when $`C_{RLRL}(C_{LRRL})`$ is negative (positive), $`P_{TL}`$ is positive (negative). When $`C_{TE}<0.8(>0)`$ and $`C_T<0(>2)`$, $`P_{TL}`$ is negative and positive otherwise. Hence determination of the magnitude and sign of $`P_{TL}`$ gives unambiguous confirmation of the existence of new physics due to scalar and tensor interactions.
In Figs. (7) and (8) we present the dependence of $`P_{TT}`$ on $`C_X`$ for the $`B\rho \mu ^+\mu ^{}`$ decay. In region I (see Fig. (7)) $`P_{TT}`$ is strongly dependent on vector type interactions $`C_{LR},C_{RR}`$ and for the negative values of $`C_{LL}`$ and $`C_T`$. On the other hand, in Region II, $`P_{TT}`$ is strongly dependent only on tensor interaction. In Region I $`P_{TT}`$ is positive (negative) for negative values of $`C_{LR}(C_{RR})`$ and it attains at negative (positive) values for positive values of $`C_{LR}(C_{RR})`$. In the second region the sign of $`P_{TT}`$ changes only for the vector interaction $`C_{RR}`$.
Depicted in Figs. (9) and (10) are the dependence of $`P_{NN}`$ on the new Wilson coefficients. The situation is quite similar to the previous case for the $`P_{TT}`$. The only difference being, $`P_{NN}`$ in Region I depends strongly on $`C_{TE}`$ rather than $`C_T`$, for their negative values, compared to that for the $`P_{TT}`$ case.
All remaining doubleโlepton polarization asymmetries for the $`B\rho \mu ^+\mu ^{}`$ decay are very small numerically and therefore we do not present them.
Through Figs. (11)โ(14) we study the dependence of $`P_{ij}`$ on the new Wilson coefficients for the $`B\rho \tau ^+\tau ^{}`$ decay, which provides richer information about the new physics effects.
In Fig. (11) the dependence of $`P_{LL}`$ on $`C_X`$ is given. We observe from this figure that $`P_{LL}`$ is very sensitive to all new Wilson coefficients except $`C_{RL}`$. It changes its sign only for the variations in $`C_T`$ and for all rest of the new Wilson coefficients $`P_{LL}`$ does not seem to change its sign. Therefore investigation of the sign of $`P_{LL}`$ can give important clue about the existence of the tensor interaction.
In Fig. (12) we present the dependence of of $`P_{LT}`$ on the new Wilson coefficients. Noting that $`P_{TL}`$ exhibits similar behavior, except several scalar coefficients, $`P_{LT}`$ is sensitive to all remaining Wilson coefficients. Similar to the $`P_{LL}`$ case, $`P_{LT}`$ changes its sign in the presence of the tensor interaction and therefore this circumstance can be quite useful in looking for new physics beyond the SM.
The dependence of $`P_{LN}P_{NL}`$ on $`C_X`$ is presented in Fig. (13). We see from this figure that $`P_{LN}`$ is very sensitive to all new Wilson coefficients, especially to the vector interaction coefficients $`C_{LL}`$ and $`C_{LR}`$.
In Fig. (14) we present the dependence of $`P_{NN}P_{TT}`$ on the new Wilson coefficients. We observe from this figure that when $`C_X`$ is negative $`P_{NN}>P_{NN}^{SM}`$ for the coefficients $`C_{LR}`$, $`C_{LL}`$, $`C_{LRRL}`$ and $`C_T`$, and $`P_{NN}>P_{NN}^{SM}`$ for the coefficients $`C_{RL}`$, $`C_{RR}`$, $`C_{RLLR}`$ and $`C_{TE}`$. On the other hand, when $`C_X`$ is positive the situation changes to the contrary, except for the tensor interaction (neglecting the narrow region for the coefficient $`C_{TE}`$). The numerical analysis for the rest of the remaining doubleโlepton polarization asymmetries for the $`B\rho \tau ^+\tau ^{}`$ decay are not presented in this work due to their negligible smallness.
It follows from the present analysis that few of the doubleโlepton polarization asymmetries show considerable departure from the SM predictions and these ones are strongly dependent on different types of interactions. Hence, the study of these quantities can play crucial role in establishing new physics beyond the SM.
At the end of this section, we would like to discuss the following problem. Could there be a case in which the branching ratio coincides with that of the SM result, while doubleโlepton polarization asymmetry does not? In order to answer this question we study the correlation between the $`P_{ij}`$ and the branching ratio $``$. We can briefly summarize the results of our numerical analysis as follows: For the $`B\rho \mu ^+\mu ^{}`$ decay, except for a very narrow region of $`C_{RR}`$, such a region is absent for all new Wilson coefficients for all of the asymmetries $`P_{ij}`$.
The $`B\rho \tau ^+\tau ^{}`$ decay is more informative for this aim, which are measurable in the experiments. In Figs. (15) and (16) we present the dependence of $`P_{LL}`$ and $`P_{LT}`$ on the branching ratio. It follows from these figures that, there indeed exists certain regions of $`C_X`$ for which the doubleโlepton polarization asymmetry differs from the SM prediction, while the branching ratio coincides with that of the SM result. We also note that, such a region exists for the remaining doubleโlepton polarization asymmetries for the tensor interaction as well.
In conclusion, in the present work we investigate the doubleโlepton polarization asymmetries when both leptons are polarized, using a general, model independent form of the effective Hamiltonian. We obtain that various doubleโlepton polarization asymmetries can serve as a good test in looking for new physics beyond the SM. We also study the correlation between $`P_{ij}`$ and the branching ratio for the $`B\rho \tau ^+\tau ^{}`$ decay and find out that there exist regions of the new Wilson coefficients for which the doubleโlepton polarization asymmetry differs considerably from the SM prediction, while the branching ratio coincides with the SM prediction. Therefore in these regions the new physics effects can be established just by measuring the doubleโlepton polarizations.
## Figure captions
Fig. (1) The dependence of the averaged doubleโlepton polarization asymmetry $`P_{LL}`$ on the new Wilson coefficients $`C_X`$, for the $`B\rho \mu ^+\mu ^{}`$ decay, in Region I.
Fig. (2) The same as in Fig. (1), but in Region II.
Fig. (3) The same as in Fig. (1), but for the averaged doubleโlepton polarization asymmetry $`P_{LT}`$.
Fig. (4) The same as in Fig. (3), but in Region II.
Fig. (5) The same as in Fig. (1), but for the averaged doubleโlepton polarization asymmetry $`P_{TL}`$.
Fig. (6) The same as in Fig. (5), but in Region II.
Fig. (7) The same as in Fig. (1), but for the averaged doubleโlepton polarization asymmetry $`P_{TT}`$.
Fig. (8) The same as in Fig. (7), but in Region II.
Fig. (9) The same as in Fig. (1), but for the averaged doubleโlepton polarization asymmetry $`P_{NN}`$.
Fig. (10) The same as in Fig. (9), but in Region II.
Fig. (11) The dependence of the averaged doubleโlepton polarization asymmetry $`P_{LL}`$ on the new Wilson coefficients $`C_X`$, for the $`B\rho \tau ^+\tau ^{}`$ decay, in Region II.
Fig. (12) The same as in Fig. (11), but for the $`P_{LT}`$.
Fig. (13) The same as in Fig. (11), but for the $`P_{LN}`$.
Fig. (14) The same as in Fig. (11), but for the $`P_{NN}`$.
Fig. (15) Parametric plot of the correlation between the averaged doubleโlepton polarization asymmetry $`P_{LL}`$ and the branching ratio for the $`B\rho \tau ^+\tau ^{}`$ decay, in Region II.
Fig. (16) Parametric plot of the correlation between the averaged doubleโlepton polarization asymmetry $`P_{LT}`$ and the branching ratio for the $`B\rho \tau ^+\tau ^{}`$ decay, in Region II. |
warning/0506/cond-mat0506673.html | ar5iv | text | # Radiation-induced magnetotransport in high-mobility two-dimensional systems: Role of electron heating
## I Introduction
The discovery of microwave induced magnetorersistance oscillations (MIMOs) and zero-resistance states (ZRS) in high-mobility two-dimensional (2D) electron gas (EG)Zud01 ; Ye ; Mani ; Zud03 has stimulated tremendous experimentalYang ; Dor ; Mani04 ; Will ; Zud04 ; Stud ; Kovalev ; Mani-apl ; Du04 ; Dor04 and theoreticalRyz ; Ryz86 ; Anderson ; Koul ; Andreev ; Durst ; Xie ; Dmitriev ; Lei03 ; Lei04 ; Ryz03 ; Vav ; Mikh ; Dietel ; Torres ; Dmitriev04 ; Inar ; Ryz0411 ; Joas interest in radiation related magneto-transport in 2D electron systems. Since theoretically it has been shown that the ZRS can be the result of the instability induced by absolute negative resistivity,Andreev the majority of microscopic models focus mainly on MIMOs in spatially uniform cases and identify the region where an negative dissipative magnetoresistance develops as that of measured zero resistance. Most of previous investigations concentrated on the range of low magnetic fields $`\omega _c/\omega 1`$ ($`\omega _c`$ stands for the cyclotron frequency) subject to a radiation of frequency $`\omega /2\pi 100`$ GHz, where MIMOs show up strongly and Shubnikov-de Haas oscillations (SdHOs) are rarely appreciable. In spite of the fact that both MIMOs and SdHOs are magnetoresistance related phenomena appearing in overlapping field regimes, little attention was paid to the influence of a microwave radiation on SdHO until a recent experimental finding at higher frequency.Kovalev Further observations clearly show that the amplitudes of SdHOs are strongly affected by microwave radiations of different frequency in both low ($`\omega _c/\omega <1`$) and high ($`\omega _c/\omega >1`$) magnetic field ranges.Mani-apl ; Du04 ; Dor04 Kovalev et al.Kovalev observed a suppression of the SdHOs around cyclotron resonance $`\omega _c\omega `$ induced by a radiation of 285 GHz. Du et al.Du04 found strong modulations of SdHO in an ultra-clean 2D sample subjected to microwaves of 146 GHz, clearly showing, in addition to the first node at $`\omega /\omega _c=1`$, higher order nodes around $`\omega /\omega _c=2`$ and 3. ManiMani-apl reported strong modulation in the amplitude of SdHOs accompanying MIMOs and zero-resistance states excited by a 163.5-GHz radiation and large dropoff of the dissipative resistivity below its dark value at high ($`\omega _c/\omega >1`$) field side when subjected to low-frequency radiation. Very recently, Dorozhkin et al.Dor04 reported both the strong suppression of the magnetoresistance caused by radiation below 30 GHz and an interesting modulation of SdHOs in the range $`\omega _c/\omega >1`$. They found that SdHOs are generally strongly damped by the radiation but there is a narrow magnetic field range in between allowed ranges of inter- and intra-Landau level transitions, where the amplitude of SdHO is insensitive to the microwave irradiation. These observations provide a more complicated and appealing picture of the microwave-related transport phenomena, which must be accounted for in any theoretical model for MIMOs.
We propose that these SdHO modulations come from the electron heating induced by the microwave radiation. Under the illumination of microwave the electron system, which continuously absorbs energy from the radiation field, would certainly be heated. Unfortunately, the electron heating has so far been ignored in most of the theoretical treatments. The electron-acoustic phonon interaction was previously considered to contribute to Landau-level broadeningLei04 or to act as a dampingInar for the orbit movement, providing a mechanism for the suppression of MIMOs when the lattice temperature increases. Besides the inelastic electron-phonon scattering also plays another important role to dissipate energy from the electron system to the lattice. The energy absorption rate is indeed small in high-mobility electron systems at low temperature as in the experiments. This, however, does not imply a negligible electron heating, since the electron energy-dissipation rate is also small because of weak electron-phonon scattering at temperature $`T1`$ K. To deal with SdHO, which is very sensitive to the smearing of the electron distribution, one has to carefully calculate the electron heating due to microwave irradiation in a uniform model.
On the other hand, microwave irradiation heats the electrons and thus greatly strengthens the thermalizing trend of the system by enhancing the electron-electron scattering rate at this low temperature regime. This enables us to describe these high-mobility 2D electron systems with a quasi-equilibrium distribution in a moving reference frame.
In this paper we pursue a theoretical investigation on MIMOs and SdHOs taking account of the electron heating under microwave irradiation. We generalize the balance equation approach to radiation-induced magnetotransport in high mobility two-dimensional electron systems. By carefully calculating the electron heating based on the balance of the energy absorption from the radiation field and the energy dissipation to the lattice through electron-phonon interactions in a typical GaAs-based heterosystem and taking into account the electrodynamic effect, we are able not only to reproduce the interesting phenomena of MIMOs in quantitative agreement with experiments in amplitudes, phases and radiation dependence of the oscillation, but also to obtain SdHO modulations observed in the experiments.
## II Formulation
### II.1 Force- and energy-balance equations
This paper is concerned with the magnetotransport in a microscopically homogeneous 2D system, and refers the measured zero resistance to the macroscopic consequence of the instability due to the occurrence of negative dissipative resistivity.
We consider $`N_\mathrm{e}`$ electrons in a unit area of an infinite quasi-2D system in the $`x`$-$`y`$ plane with a confining potential $`V(z)`$ in the $`z`$ direction. These electrons, in addition to interacting with each other, are scattered by random impurities and/or disorders and by phonons in the lattice. Within the magnetic field range relevant to MIMO phenomenon, the experiments exclude the onset of the quantum Hall effect, thus allowing us to assume that the 2D electrons are in extended states.
To include possible elliptically polarized microwave illumination we assume that a dc electric field $`๐_0`$ and a high-frequency (HF) ac field of angular frequency $`\omega `$,
$$๐(t)๐_s\mathrm{sin}(\omega t)+๐_c\mathrm{cos}(\omega t),$$
(1)
are applied inside the 2D system in the $`x`$-$`y`$ plane, together with a magnetic field $`๐=(0,0,B)`$ along the $`z`$ direction. The spatial homogeneity of the fields and the parabolic band structure allows to describe the transport of this system in terms of its center-of-mass (c.m.) motion and the relative motion, i.e. the motion of electrons in the reference frame moving with the c.m.Ting ; Lei85 ; Lei851 The center-of-mass momentum and coordinate of the 2D electron system are defined as $`๐_j๐ฉ_j`$ and $`๐N_\mathrm{e}^1_j๐ซ_j`$ with $`๐ฉ_j(p_{jx},p_{jy})`$ and $`๐ซ_j(x_j,y_j)`$ being the momentum and coordinate of the $`j`$th electron in the 2D plane, respectively, and the relative electron momentum and coordinate are defined as $`๐ฉ_j^{}๐ฉ_j๐/N_\mathrm{e}`$ and $`๐ซ_j^{}๐ซ_j๐`$, respectively. In terms of these variables, the Hamiltonian of the system, $`H`$, can be written as the sum of a center-of-mass part $`H_{\mathrm{cm}}`$ and a relative electron part $`H_{\mathrm{er}}`$ ($`๐(๐ซ)`$ is the vector potential of the $`๐`$ field),
$`H_{\mathrm{cm}}`$ $`=`$ $`{\displaystyle \frac{1}{2N_\mathrm{e}m}}(๐N_\mathrm{e}e๐(๐))^2N_\mathrm{e}e(๐_0+๐(t))๐,`$ (2)
$`H_{\mathrm{er}}`$ $`=`$ $`{\displaystyle \underset{j}{}}\left[{\displaystyle \frac{1}{2m}}\left(๐ฉ_j^{}e๐(๐ซ_j^{})\right)^2+{\displaystyle \frac{p_{jz}^2}{2m_z}}+V(z_j)\right]`$ (3)
$`+{\displaystyle \underset{i<j}{}}V_c(๐ซ_i^{}๐ซ_j^{},z_i,z_j),`$
together with electron-impurity and electron-phonon interactions
$`H_{\mathrm{ei}}`$ $`=`$ $`{\displaystyle \underset{j,a,๐ช_{}}{}}u(๐ช_{},z_a)e^{i๐ช_{}(๐+๐ซ_j^{}๐ซ_a)},`$ (4)
$`H_{\mathrm{ep}}`$ $`=`$ $`{\displaystyle \underset{j,๐ช_{}}{}}M(๐ช_{},q_z)(b_๐ช+b_๐ช^{})e^{i๐ช_{}(๐+๐ซ_j^{})}.`$ (5)
Here $`m`$ and $`m_z`$ are, respectively, the electron effective mass parallel and perpendicular to the 2D plane, and $`V_c`$ stands for the electron-electron Coulomb interaction; $`u(๐ช_{},z_a)`$ is the potential of the $`a`$th impurity locating at $`(๐ซ_a,z_a)`$; $`b_๐ช^{}(b_๐ช)`$ are the creation (annihilation) operators of the bulk phonon with wavevector $`๐ช=(๐ช_{},q_z)`$ and $`M(๐ช_{},q_z)`$ is the matrix element of the electron-phonon interaction in the 3D plane-wave representation. Note that the uniform electric field (dc and ac) appears only in $`H_{\mathrm{cm}}`$, and that $`H_{\mathrm{er}}`$ is just the Hamiltonian of a quasi-2D system subjected to a magnetic field without an electric field. The coupling between the center-of-mass and the relative electrons appears only in the exponential factor $`\mathrm{exp}(i๐ช_{}๐)`$ inside the 2D momemtum $`๐ช_{}`$ summation in $`H_{\mathrm{ei}}`$ and $`H_{\mathrm{ep}}`$.Lei851 The balance equation treatment starts with the Heisenberg operator equation for the rate of change of the center-of-mass velocity $`\dot{๐}=i[๐,H]+๐/t`$ with $`๐=i[๐,H]`$, and that for the rate of change of the relative electron energy $`\dot{H}_{\mathrm{er}}=\mathrm{i}[H_{\mathrm{er}},H]`$. Then we proceed with the determination of their statistical averages.
As proposed in Ref. Lei85, , the c.m. coordinate operator $`๐`$ and velocity operator $`๐`$ can be treated classically, i.e. as the time-dependent expectation values of c.m. coordinate and velocity, $`๐(t)`$ and $`๐(t)`$, such that $`๐(t)๐(t^{})=_t^{}^t๐(s)๐s`$. We are concerned with the steady transport state under an irradiation of single frequency and focus on the photon-induced dc resistivity and the energy absorption of the HF field. These quantities are directly related to the time-averaged and/or base-frequency oscillating components of the c.m. velocity. Although higher harmonics of the current may affect the dc and lower harmonic terms of the drift velocity through entering the damping force and energy exchange rates in the resulting equations, in an ordinary semiconductor the power of even the third harmonic current is rather weak as compared to the fundamental current. For the HF field intensity in the MIMO experiments, the effect of higher harmonic current is safely negligible. Hence, it suffices to assume that the c.m. velocity, i.e. the electron drift velocity, consists of a dc part $`๐ฏ_0`$ and a stationary time-dependent part $`๐ฏ(t)`$ of the form
$$๐(t)=๐ฏ_0๐ฏ_1\mathrm{cos}(\omega t)๐ฏ_2\mathrm{sin}(\omega t).$$
(6)
With this, the exponential factor in the operator equations can be expanded in terms of Bessel functions $`J_n(x)`$,
$`\mathrm{e}^{\mathrm{i}๐ช_{}_t^{}^t๐(s)๐s}={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}J_n^2(\xi )\mathrm{e}^{\mathrm{i}(๐ช_{}๐ฏ_0n\omega )(tt^{})}+`$
$`{\displaystyle \underset{m0}{}}\mathrm{e}^{\mathrm{i}m(\omega t\phi )}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}J_n(\xi )J_{nm}(\xi )\mathrm{e}^{\mathrm{i}(๐ช_{}๐ฏ_0n\omega )(tt^{})}.`$
Here the argument in the Bessel functions
$$\xi \frac{1}{\omega }\left[(๐ช_{}๐ฏ_1)^2+(๐ช_{}๐ฏ_2)^2\right]^{\frac{1}{2}},$$
(7)
and $`\mathrm{tan}\phi =(๐ช_{}๐ฏ_2)/(๐ช_{}๐ฏ_1)`$.
Under the influence of a modest-strength HF electric field the electron system is far from equilibrium. However, the distribution function of relative electrons, which experience no electric field directly, may be close to an quasi-equilibrium type distribution function. For the experimental GaAs-based ultra-clean 2D electron systems having carrier mobility of the order of $`2000`$ m<sup>2</sup>/Vs, the elastic momentum scattering rate is around $`\tau _m^110`$ mK. In these systems, the thermalization time $`\tau _{\mathrm{th}}`$ (i.e. the time for system to return to its internal equilibrating state when it is deviated from), estimated conservatively using electron-electron (e-e) interaction related inelastic scattering time $`\tau _{ee}`$ calculated with an equilibrium distribution function at temperature $`T=1`$ K, is also around $`\tau _{\mathrm{th}}^1\tau _{ee}^110`$ mK. The illumination of microwave certainly heats the electrons. Even an electron heating comparable to a couple of degrees temperature rise would greatly enhance $`\tau _{ee}^1`$, such that the thermalization time $`\tau _{\mathrm{th}}`$ would become much shorter than the momentum relaxation time $`\tau _m`$ under microwave irradiation.note-life The relative electron systems subject to a modest radiation would rapidly thermalize and can thus be described reasonably by a Fermi-type distribution function at an average electron temperature $`T_\mathrm{e}`$ in the reference frame moving with the center-of-mass. This allows us to carry out the statistical average of the operator equations for the rates of changes of the c.m. velocity $`๐`$ and relative electron energy $`H_{\mathrm{er}}`$ to the leading order in $`H_{\mathrm{ei}}`$ and $`H_{\mathrm{ep}}`$ with succinct forms.
For the determination of unknown parameter $`๐ฏ_0`$, $`๐ฏ_1`$, $`๐ฏ_2`$, and $`T_\mathrm{e}`$, it suffices to know the damping force up to the base frequency oscillating term $`๐
(t)=๐
_0+๐
_s\mathrm{sin}(\omega t)+๐
_c\mathrm{cos}(\omega t)`$, and the energy-related quantities up to the time-average terms. We finally obtain the force and energy balance equations,
$`N_\mathrm{e}e๐_0+N_\mathrm{e}e(๐ฏ_0\times ๐)+๐
_0=0,`$ (9)
$`๐ฏ_1={\displaystyle \frac{e๐_s}{m\omega }}+{\displaystyle \frac{๐
_s}{N_\mathrm{e}m\omega }}{\displaystyle \frac{e}{m\omega }}(๐ฏ_2\times ๐),`$
$``$ $`๐ฏ_2={\displaystyle \frac{e๐_c}{m\omega }}+{\displaystyle \frac{๐
_c}{N_\mathrm{e}m\omega }}{\displaystyle \frac{e}{m\omega }}(๐ฏ_1\times ๐),`$ (11)
$`N_\mathrm{e}e๐_0๐ฏ_0+S_\mathrm{p}W=0.`$
Here
$`๐
_0={\displaystyle \underset{๐ช_{}}{}}\left|U(๐ช_{})\right|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}๐ช_{}J_n^2(\xi )\mathrm{\Pi }_2(๐ช_{},\omega _0n\omega )`$
$`+{\displaystyle \underset{๐ช}{}}\left|M(๐ช)\right|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}๐ช_{}J_n^2(\xi )\mathrm{\Lambda }_2(๐ช,\omega _0+\mathrm{\Omega }_๐ชn\omega )`$ (12)
is the time-averaged damping force,
$`S_\mathrm{p}={\displaystyle \underset{๐ช_{}}{}}\left|U(๐ช_{})\right|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n\omega J_n^2(\xi )\mathrm{\Pi }_2(๐ช_{},\omega _0n\omega )`$
$`+{\displaystyle \underset{๐ช}{}}\left|M(๐ช)\right|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}n\omega J_n^2(\xi )\mathrm{\Lambda }_2(๐ช,\omega _0+\mathrm{\Omega }_๐ชn\omega )`$ (13)
is the time-averaged rate of the electron energy absorption from the HF field, and
$$W=\underset{๐ช}{}\left|M(๐ช)\right|^2\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{\Omega }_๐ชJ_n^2(\xi )\mathrm{\Lambda }_2(๐ช,\omega _0+\mathrm{\Omega }_๐ชn\omega )$$
(14)
is the time-averaged rate of the electron energy dissipation to the lattice due to electron-phonon scatterings. The oscillating frictional force amplitudes $`๐
_s๐
_{22}๐
_{11}`$ and $`๐
_c๐
_{21}+๐
_{12}`$ are given by ($`\mu =1,2`$)
$`๐
_{1\mu }={\displaystyle \underset{๐ช_{}}{}}๐ช_{}\eta _\mu |U(๐ช_{})|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[J_n^2(\xi )\right]^{}\mathrm{\Pi }_1(๐ช_{},\omega _0n\omega )`$
$`{\displaystyle \underset{๐ช}{}}๐ช_{}\eta _\mu |M(๐ช)|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[J_n^2(\xi )\right]^{}\mathrm{\Lambda }_1(๐ช,\omega _0+\mathrm{\Omega }_๐ชn\omega ),`$ (15)
$`๐
_{2\mu }={\displaystyle \underset{๐ช_{}}{}}๐ช_{}{\displaystyle \frac{\eta _\mu }{\xi }}|U(๐ช_{})|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}2nJ_n^2(\xi )\mathrm{\Pi }_2(๐ช_{},\omega _0n\omega )`$
$`+{\displaystyle \underset{๐ช}{}}๐ช_{}{\displaystyle \frac{\eta _\mu }{\xi }}|M(๐ช)|^2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}2nJ_n^2(\xi )\mathrm{\Lambda }_2(๐ช,\omega _0+\mathrm{\Omega }_๐ชn\omega ).`$ (16)
In these expressions, $`\eta _\mu ๐ช_{}๐ฏ_\mu /\omega \xi `$; $`\omega _0๐ช_{}๐ฏ_0`$; $`U(๐ช_{})`$ and $`M(๐ช)`$ are effective impurity and phonon scattering potentials (including effects of the spatial distribution of impurities and the form factor of quasi-2D electrons).Lei851 $`\mathrm{\Pi }_2(๐ช_{},\mathrm{\Omega })`$ and $`\mathrm{\Lambda }_2(๐ช,\mathrm{\Omega })=2\mathrm{\Pi }_2(๐ช_{},\mathrm{\Omega })[n(\mathrm{\Omega }_๐ช/T)n(\mathrm{\Omega }/T_\mathrm{e})]`$ (with $`n(x)1/(\mathrm{e}^x1)`$) are the imaginary parts of the electron density correlation function and electron-phonon correlation function in the presence of the magnetic field. $`\mathrm{\Pi }_1(๐ช_{},\mathrm{\Omega })`$ and $`\mathrm{\Lambda }_1(๐ช,\mathrm{\Omega })`$ are the real parts of these two correlation functions.
Effects of a microwave radiation on electron transport first come from the HF field induced c.m. motion (electron drift motion) and the related change of the electron distribution. In addition to this, the HF field also enters via the argument $`\xi `$ of the Bessel functions in $`๐
_0`$, $`๐
_{\mu \nu }`$, $`W`$ and $`S_\mathrm{p}`$. Compared with that without a HF field, we see that in an electron gas having impurity and/or phonon scatterings (otherwise homogeneous), a HF field of frequency $`\omega `$ opens additional channels for electron transition: an electron in a state can absorb or emit one or several photons of frequency $`\omega `$ and scattered to a different state with the help of impurities and/or phonons. The sum over $`|n|1`$ represents contributions of real single and multiple photon participating processes. The role of these processes is two folds. On the one hand, they contribute additional damping force to the moving electrons, giving rise directly to photoresistance, and at the same time, transfer energy from the HF field to the electron system, resulting in electron heating, i.e. another change (smearing) in the electron distribution.Lei98 Furthermore, the radiation field, showing up in the term with $`J_0(\xi )`$ in $`๐
_0`$, $`๐
_{\mu \nu }`$ and $`W`$, gives rise to another effective change of damping forces and energy-loss rate, without emission or absorption of real photons. This virtual photon process also contributes to photoresistance.Lei-apl All these effects are carried by parameters $`๐ฏ_0`$, $`๐ฏ_1`$, $`๐ฏ_2`$ and $`T_\mathrm{e}`$. Eqs. (9)-(11) form a closed set of equations for the determination of these parameters when $`๐_0`$, $`๐_c`$ and $`๐_s`$ are given in a 2D system subjected to a magnetic field $`B`$ at temperature $`T`$.
### II.2 Longitudinal and transverse resistivities
The nonlinear resistivity in the presence of a high-frequency field is easily obtained from Eq. (9). Taking $`๐ฏ_0`$ to be in the $`x`$ direction, $`๐ฏ_0=(v_{0x},0,0)`$, we immediately get the transverse and longitudinal resistivities,
$`R_{xy}`$ $``$ $`{\displaystyle \frac{E_{0y}}{N_\mathrm{e}ev_{0x}}}={\displaystyle \frac{B}{N_\mathrm{e}e}},`$ (17)
$`R_{xx}`$ $``$ $`{\displaystyle \frac{E_{0x}}{N_\mathrm{e}ev_{0x}}}={\displaystyle \frac{F_0}{N_\mathrm{e}^2e^2v_{0x}}}.`$ (18)
The linear magnetoresistivity is the weak dc current limit ($`v_{0x}0`$):
$`R_{xx}`$ $`=`$ $`{\displaystyle \underset{๐ช_{}}{}}q_x^2{\displaystyle \frac{|U(๐ช_{})|^2}{N_\mathrm{e}^2e^2}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}J_n^2(\xi ){\displaystyle \frac{\mathrm{\Pi }_2}{\mathrm{\Omega }}}|_{\mathrm{\Omega }=n\omega }`$ (19)
$`{\displaystyle \underset{๐ช}{}}q_x^2{\displaystyle \frac{\left|M(๐ช)\right|^2}{N_\mathrm{e}^2e^2}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}J_n^2(\xi ){\displaystyle \frac{\mathrm{\Lambda }_2}{\mathrm{\Omega }}}|_{\mathrm{\Omega }=\mathrm{\Omega }_๐ช+n\omega }.`$
Note that although according to Eqs. (12), (18) and (19), the longitudinal magnetoresistivity $`R_{xx}`$ can be formally written as the sum of contributions from various individual scattering mechanisms, all the scattering mechanisms have to be taken into account simultaneously in solving the momentum- and energy-balance equations (9), (11) and (11) for $`๐ฏ_1`$, $`๐ฏ_2`$ and $`T_\mathrm{e}`$, which enter the Bessel functions and other parts in the expression of $`R_{xx}`$.
### II.3 Landau-level broadening
In the present model the effects of interparticle Coulomb screening are included in the electron complex density correlation function $`\mathrm{\Pi }(๐ช_{},\mathrm{\Omega })=\mathrm{\Pi }_1(๐ช_{},\mathrm{\Omega })+i\mathrm{\Pi }_2(๐ช_{},\mathrm{\Omega })`$, which, in the random phase approximation, can be expressed as
$$\mathrm{\Pi }(๐ช_{},\mathrm{\Omega })=\frac{\mathrm{\Pi }_0(๐ช_{},\mathrm{\Omega })}{ฯต(๐ช_{},\mathrm{\Omega })},$$
(20)
where
$$ฯต(๐ช_{},\mathrm{\Omega })1V(q_{})\mathrm{\Pi }_0(๐ช_{},\mathrm{\Omega })$$
(21)
is the complex dynamical dielectric function,
$$V(q_{})=\frac{e^2}{2ฯต_0\kappa q_{}}H(q_{})$$
(22)
is the effective Coulomb potential with $`\kappa `$ the dielectric constant of the material and $`H(q_{})`$ is a 2D wavefunction-related overlapping integration,Lei851 $`\mathrm{\Pi }_0(๐ช_{},\mathrm{\Omega })=\mathrm{\Pi }_{01}(๐ช_{},\mathrm{\Omega })+i\mathrm{\Pi }_{02}(๐ช_{},\mathrm{\Omega })`$ is the complex density correlation function of the independent electron system in the presence of the magnetic field. With this dynamically screened density correlation function the collective plasma modes of the 2DES are incorporated. Disregard these collective modes one can just use a static screening $`ฯต(๐ช_{},0)`$ instead.
The $`\mathrm{\Pi }_{02}(๐ช_{},\mathrm{\Omega })`$ function of a 2D system in a magnetic field can be written in terms of Landau representation:Ting
$`\mathrm{\Pi }_{02}(๐ช_{},\mathrm{\Omega })={\displaystyle \frac{1}{2\pi l_\mathrm{B}^2}}{\displaystyle \underset{n,n^{}}{}}C_{n,n^{}}(l_\mathrm{B}^2q_{}^2/2)\mathrm{\Pi }_2(n,n^{},\mathrm{\Omega }),`$ (23)
$`\mathrm{\Pi }_2(n,n^{},\mathrm{\Omega })={\displaystyle \frac{2}{\pi }}{\displaystyle ๐\epsilon \left[f(\epsilon )f(\epsilon +\mathrm{\Omega })\right]}`$
$`\times \mathrm{Im}G_n(\epsilon +\mathrm{\Omega })\mathrm{Im}G_n^{}(\epsilon ),`$ (24)
where $`l_\mathrm{B}=\sqrt{1/|eB|}`$ is the magnetic length,
$$C_{n,n+l}(Y)n![(n+l)!]^1Y^le^Y[L_n^l(Y)]^2$$
(25)
with $`L_n^l(Y)`$ the associate Laguerre polynomial, $`f(\epsilon )=\{\mathrm{exp}[(\epsilon \mu )/T_\mathrm{e}]+1\}^1`$ the Fermi distribution function, and $`\mathrm{Im}G_n(\epsilon )`$ is the imaginary part of the electron Greenโs function, or the density of states (DOS), of the Landau level $`n`$. The real part function $`\mathrm{\Pi }_{01}(๐ช_{},\mathrm{\Omega })`$ and corresponding $`\mathrm{\Lambda }_{01}(๐ช_{},\mathrm{\Omega })`$ function can be derived from their imaginary parts via the Kramers-Kronig relation.
In principle, to obtain the Greenโs function $`\mathrm{Im}G_n(\epsilon )`$, a self-consistent calculation has to be carried out from the Dyson equation for the self-energy with all the impurity, phonon and e-e scatterings included. The resultant $`G_n`$ is generally a complicated function of the magnetic field, temperature, and Landau-level index $`n`$, also dependent on the different kinds of scatterings. Such a calculation is beyond the scope of the present study. In this paper we model the DOS function with a Gaussian-type form ($`\epsilon _n`$ is the energy of the $`n`$-th Landau level):Ando82 ; Raikh93
$$\mathrm{Im}G_n(\epsilon )=\frac{\sqrt{2\pi }}{\mathrm{\Gamma }}\mathrm{exp}\left[\frac{2(\epsilon \epsilon _n)^2}{\mathrm{\Gamma }^2}\right]$$
(26)
with a broadening width given by
$$\mathrm{\Gamma }=\left(\frac{8e\omega _c\alpha }{\pi m\mu _0}\right)^{1/2},$$
(27)
where $`\mu _0`$ is the linear mobility in the absence of the magnetic field and $`\alpha `$ is a semiempirical parameter to take into account the difference of the transport scattering time $`\tau _m`$ determining the mobility $`\mu _0`$, from the single particle lifetime $`\tau _s`$ related to Landau level broadening. The latter depends on elastic scatterings of different types and their relative strengths, as well as contributions of electron-phonon and electron-electron scatterings. $`\alpha `$ will be served as the only adjustable parameter in the present investigation. Unlike the semielliptic function, which can model only separated Landau-level case, a Gaussian-type broadening function can reasonably cover both the separated-level and overlapping-level regimes.
### II.4 Effect of radiative decay
The HF electric field $`๐(t)`$ appearing in Eqs. (8) and (9) is the total (external and induced) field really acting on the 2D electrons. Experiments are always performed under the condition of giving external radiation. In this paper we assume that the electromagnetic wave is incident perpendicularly (along $`z`$-axis) upon 2DEG from the vacuum with the incident electric field of
$$๐_\mathrm{i}(t)=๐_{\mathrm{i}s}\mathrm{sin}(\omega t)+๐_{\mathrm{i}c}\mathrm{cos}(\omega t)$$
(28)
at plane $`z=0`$. The relation between $`๐(t)`$ and $`๐_\mathrm{i}(t)`$ is easily obtained by solving the Maxwell equations connecting both sides of the 2DEG which is carrying a sheet current density $`N_\mathrm{e}e๐ฏ(t)`$. If the 2DEG locates under the surface plane at $`z=0`$ of a thick (treated as semi-infinite) semiconductor substrate having a refraction index $`n_s`$, we haveChiu ; Liu
$$๐(t)=\frac{N_\mathrm{e}e๐ฏ(t)}{(1+n_s)ฯต_0c}+\frac{2}{1+n_s}๐_\mathrm{i}(t).$$
(29)
If the 2DEG is contained in a thin sample suspended in vacuum at the plane $`z=0`$, then
$$๐(t)=\frac{N_\mathrm{e}e๐ฏ(t)}{2ฯต_0c}+๐_\mathrm{i}(t).$$
(30)
In the numerical calculation of this paper we consider the latter case and use Eq. (30) for the total selfconsistent field $`๐(t)`$ in Eqs. (9) and (11). This electrodynamic effect,Chiu ; Liu recently refered as radiative decay,Mikh gives rise to an additional damping in the 2DEG response to a given incident HF field. The induced damping turns out to be much stronger than the intrinsic damping due to scattering-related forces $`๐
_s`$ and $`๐
_c`$ for the experimental high-mobility systems at low temperatures. For almost all the cases pertinent to MIMO experiments we can neglect the forces $`๐
_s`$ and $`๐
_c`$ completely in solving $`๐ฏ_1(v_{1x},v_{1y})`$ and $`๐ฏ_2(v_{2x},v_{2y})`$ from Eqs.(9) and (11) for given incident fields $`๐_{\mathrm{i}s}`$ and $`๐_{\mathrm{i}c}`$, and obtain explicitly
$$\begin{array}{ccc}v_{1x}& =& (a\chi _{sx}+b\chi _{sy})/\mathrm{\Delta }\\ v_{1y}& =& (a\chi _{sy}b\chi _{sx})/\mathrm{\Delta }\\ v_{2x}& =& (a\chi _{cx}b\chi _{cy})/\mathrm{\Delta }\\ v_{2y}& =& (a\chi _{cy}+b\chi _{cx})/\mathrm{\Delta }\end{array}$$
(31)
with $`\mathrm{\Delta }=(1\delta _\omega ^2+\gamma _\omega ^2)^2+(2\gamma _\omega \delta _\omega )^2`$, and
$`\begin{array}{ccc}\chi _{sx}& =& \nu _{sx}\delta _\omega \nu _{cy}+\gamma _\omega \nu _{cx}\\ \chi _{sy}& =& \nu _{sy}+\delta _\omega \nu _{cx}+\gamma _\omega \nu _{cy}\\ \chi _{cx}& =& \nu _{cx}+\delta _\omega \nu _{sy}\gamma _\omega \nu _{sx}\\ \chi _{cy}& =& \nu _{cy}\delta _\omega \nu _{sx}+\gamma _\omega \nu _{sy}\end{array}`$ (36)
Here
$$\nu _\eta \frac{eE_{\mathrm{i}\eta }}{m\omega }(\eta =sx,sy,cx,cy),$$
(37)
$`\delta _\omega \omega _c/\omega `$ and $`\gamma _\omega \gamma /\omega `$ with
$$\gamma =\frac{N_\mathrm{e}e^2}{2mฯต_0c}.$$
(38)
With these $`๐ฏ_1`$ and $`๐ฏ_2`$, the argument $`\xi `$ entering the Bessel functions is obtained. All the transport quantities, such as $`S_\mathrm{p}`$, $`W`$ and $`R_{xx}`$, can be calculated directly with the electron temperature $`T_\mathrm{e}`$ determined from the energy balance equation (11).
## III Numerical results for GaAs-based systems
As in the experiments, we focus our attention on high mobility 2DEGs formed by GaAs/AlGaAs heterojunctions. For these systems at temperature $`T1`$ K, the dominant contributions to the energy absorption $`S_\mathrm{p}`$ and photoresistivity $`R_{xx}R_{xx}(0)`$ come from the impurity-assisted photon-absorption and emission process. At different magnetic field strength, this process is associated with electron transitions between either inter-Landau level states or intra-Landau-level states. According to (26), the width of each Landau level is about $`2\mathrm{\Gamma }`$. The condition for inter-Landau level transition with impurity-assisted single-photon processnote1 is $`\omega >\omega _c2\mathrm{\Gamma }`$, or $`\omega _c/\omega <a_{\mathrm{inter}}=(\beta +\sqrt{\beta ^2+4})^2/4`$; and that for impurity-assisted intra-Landau level transition is $`\omega <2\mathrm{\Gamma }`$, or $`\omega _c/\omega >a_{\mathrm{intra}}=\beta ^2`$, here $`\beta =(32e\alpha /\pi m\mu _0\omega )^{\frac{1}{2}}`$. However, since the DOS of each Landau level is assumed to be Gaussian rather than a clear cutoff function and the multi-photon processes also play roles, the transition boundaries between different regimes may be somewhat smeared.
As indicated by experiments,Umansky although long range scattering due to remote donors always exists in the 2D heterostructures, in ultra-clean GaAs-based 2D samples having mobility of order of $`10^3`$ m<sup>2</sup>/Vs, the remote donor scattering is responsible for merely $`10\%`$ or less of the total momentum scattering rate. The dominant contribution to the momentum scattering rate comes from short-range scatterers such as residual impurities or defects in the background. Furthermore, even with the same momentum scattering rate the remote impurity scattering is much less efficient in contributing to microwave-induced magnetoresistance oscillations than short-ranged background impurities or defects.Lei0409219 Therefore, in the numerical calculations in this paper we assume that the elastic scatterings are due to short-range impurities randomly distributed throughout the GaAs region. The impurity densities are determined by the requirement that electron total linear mobility at zero magnetic field equals the giving value at lattice temperature $`T`$. Possibly, long-range remote donnor scattering may give rise to important contribution to the Landau-level broadening. This effect, together with the role of electron-phonon and electron-electron scatterings, is included in the semiempirical parameter $`\alpha `$ in the expression (27).
In order to obtain the energy dissipation rate from the electron system to the lattice, $`W`$, we take into account scatterings from bulk longitudinal acoustic (LA) and transverse acoustic (TA) phonons (via the deformation potential and piezoelectric couplings), as well as from longitudinal optical (LO) phonons (via the Frรถhlich coupling) in the GaAs-based system. The relevant matrix elements are well known.Lei851 The material and coupling parameters for the system are taken to be widely accepted values in bulk GaAs: electron effective mass $`m=0.068m_\mathrm{e}`$ ($`m_\mathrm{e}`$ is the free electron mass), transverse sound speed $`v_{\mathrm{st}}=2.48\times 10^3`$ m/s, longitudinal sound speed $`v_{\mathrm{sl}}=5.29\times 10^3`$ m/s, acoustic deformation potential $`\mathrm{\Xi }=8.5`$ eV, piezoelectric constant $`e_{14}=1.41\times 10^9`$ V/m, dielectric constant $`\kappa =12.9`$, material mass density $`d=5.31`$ g/cm<sup>3</sup>.
### III.1 100 GHz
Figure 1 shows the calculated energy absorption rate $`S_\mathrm{p}`$, the electron temperature $`T_\mathrm{e}`$ and the longitudinal magnetoresistivity $`R_{xx}`$ as functions of $`\omega _c/\omega `$ for a 2D system having an electron density of $`N_\mathrm{e}=3.0\times 10^{15}`$ m<sup>-2</sup>, a linear mobility of $`\mu _0=2000`$ m<sup>2</sup>/Vs and a broadening parameter of $`\alpha =10`$, subject to linearly $`x`$-direction polarized incident microwave radiations of frequency $`\omega /2\pi =100`$ GHz having four different amplitudes $`E_{\mathrm{i}s}=2.2,3,4`$ and 5 V/cm at a lattice temperature of $`T=1`$ K. The energy absorption rate $`S_\mathrm{p}`$ exhibits a broad main peak at cyclotron resonance $`\omega _c/\omega =1`$ and secondary peaks at harmonics $`\omega _c/\omega =1/2,1/3,1/4`$. The electron heating has similar feature: $`T_\mathrm{e}`$ exhibits peaks around $`\omega _c/\omega =1,1/2,1/3,1/4`$. For this GaAs system $`\beta =0.65`$, $`a_{\mathrm{inter}}=1.6`$ and $`a_{\mathrm{intra}}=4.7`$. We can see that, at lower magnetic fields, especially $`\omega _c/\omega <1.4`$, the system absorbs enough energy from the radiation field via inter-Landau level transitions and $`T_\mathrm{e}`$ is significantly higher than $`T`$, with the maximum as high as 21 K around $`\omega _c/\omega =1`$. With increasing strength of the magnetic field the inter-Landau level transition weakens (impurity-assisted single-photon process is mainly allowed when $`\omega _c/\omega <a_{\mathrm{inter}}=1.6`$) and the absorbed energy decreases rapidly. Within the range $`2<\omega _c/\omega <4`$ before intra-Landau level transitions can take place, $`S_\mathrm{p}`$ is two orders of magnitude smaller than that in the low magnetic field range. Correspondingly the electron temperature $`T_\mathrm{e}`$ is only slightly higher than the lattice temperature $`T`$. The magnetoresistivity $`R_{xx}`$ showing in the upper part of Fig. 1, exhibits interesting features. MIMOs (with fixed points rather than extrema at $`\omega _c/\omega =1,1/2,1/3,1/4`$) clearly appear at lower magnetic fields, which are insensitive to the electron heating even at $`T_\mathrm{e}`$ of order of 20 K. SdHOs appearing in the higher magnetic field side, however, are damped due to the rise of the electron temperature $`T_\mathrm{e}>1`$ K as compared with that without radiation. With an increase in the microwave amplitude from $`E_{\mathrm{i}s}=2.2`$ V/cm to $`5`$ V/cm, MIMOs become much stronger and SdHOs are further damped. But the radiation-induced SdHO damping is always relatively smaller within $`2.4<\omega _c/\omega <4`$ between allowed ranges of inter- and intra-Landau level transitions.
It is worth noting that the predicted MIMOs here exhibit much improved agreement with experiments over previous theoretical models. The maxima of $`R_{xx}`$ oscillation locate at $`\omega /\omega _c=j\delta _{}`$ and minima at $`\omega /\omega _c=j+\delta _+`$, with $`\delta _\pm 0.230.25`$ for $`j=2,3,4\mathrm{}`$ and $`\delta _\pm 0.160.18`$ for $`j=1`$ (see Fig. 2). These phase details, as well as the absolute (rather than reduced) magnitudes of the oscillation amplitudes and the required incident microwave strengths to induce oscillations are in good quantitative agreement with experiments.Mani ; Zud03 ; Yang ; Dor ; Mani04
The MIMOs depend on the polarization of the incident microwave field in respect to the dc field $`๐_0`$. Physically this is clear in the present model since it is through the c.m. motion that a HF field affects the photoresistivity of the 2D electron system. Under the influence of a magnetic field perpendicular to the plane, the c.m. performs a cyclic motion of frequency $`\omega _c`$ in the 2D plane. A perpendicularly incident circularly-polarized microwave would accelerate or decelerate this cyclic motion depending on the HF electric field circling with or against it. Thus, at fix incident power, a left-polarized microwave would yield much stronger effect on the $`R_{xx}`$ oscillation than a right-polarized one and this effect is apparently strongest in the vicinity of cylcotron resonance $`\omega _c/\omega =1`$. The difference between the $`x`$-direction linearly polarized wave and the $`y`$-direction linearly polarized wave, however, comes mainly from the the different angle of radiation-induced c.m. motion with respect to the dc current, and thus not so sensitive to that of the $`\omega _c/\omega `$ range. In Fig. 2 we plot the calculated $`R_{xx}`$ versus $`\omega /\omega _c`$ for the same system as described in Fig. 1, subject to a 100 GHz microwave radiation having a fixed incident power of $`P_\mathrm{i}=210`$ W/m<sup>2</sup> (equivalent to an incident amplitude $`E_{\mathrm{i}s}=4`$ V/cm of linear polarization) but four different polarizations: linear $`x`$-polarizaton, linear $`y`$-polarization, left circular polarization and right circular polarization. Their difference is clearly seen.
### III.2 50 GHz and lower frequency
Figure 3 shows the energy absorption rate $`S_\mathrm{p}`$, the electron temperature $`T_\mathrm{e}`$ and the longitudinal magnetoresistivity $`R_{xx}`$ as functions of $`\omega _c/\omega `$ for a 2D system having an electron density of $`N_\mathrm{e}=3.0\times 10^{15}`$ m<sup>-2</sup>, a linear mobility of $`\mu _0=2500`$ m<sup>2</sup>/Vs and a broadening parameter of $`\alpha =12.5`$, subject to linearly $`x`$-direction polarized incident microwave radiations of frequency $`\omega /2\pi =50`$ GHz having four different amplitudes $`E_{\mathrm{i}s}=0.8,1.2,2.0`$ and 3.5 V/cm at a lattice temperature of $`T=1`$ K. For this GaAs system at 50 GHz $`a_{\mathrm{inter}}=1.9`$ and $`a_{\mathrm{intra}}=2.4`$. The intra-Landau level single-photon transitions are allowed when $`\omega _c/\omega >2.4`$, yielding, at the high $`\omega _c/\omega `$ side, an absorption rate $`S_\mathrm{p}`$ somewhat larger, an electron temperature $`T_\mathrm{e}`$ somewhat higher, and a SdHO damping stronger than those in the 100-GHz case (Fig. 1). On the other hand, at equivalent HF field strength the multiphoton processes are more important at lower frequency. This helps to enhance the absorption $`S_\mathrm{p}`$ in the range $`1.9<\omega _c/\omega <2.4`$, where the single-photon process is forbidden and to increase the two-photon resonance in $`S_\mathrm{p}`$ and $`T_\mathrm{e}`$ around $`\omega /\omega _c=1.5,2.5`$ and 3.5 (see $`S_\mathrm{p}`$ and $`T_\mathrm{e}`$ curves corresponding to $`E_{\mathrm{i}s}=3.5`$ V/cm in Fig. 3). The effect of the two-photon process can also be seen clearly in the $`R_{xx}`$-vs-$`\omega /\omega _c`$ curves as shown in Fig. 4, where the $`R_{xx}`$ curve of $`E_{\mathrm{i}s}=3.5`$ V/cm exhibits obvious shoulders around $`\omega /\omega _c=1.5,2.5`$ and 3.5, and the descends down around $`\omega /\omega _c=0.6`$. This kind of two-photon process was clearly seen in the experiments.Zud03 ; Zud04
At even lower frequency, such as 30 GHz and 20 GHz, the ranges for intra-Landau level and inter-Landau level single-photon transitions overlap. The enhanced effect of the virtual photon process, together with enhanced multiphoton-assisted electron transition, pushes the resistivity $`R_{xx}`$ remarkably down below the average of its oscillatory curve without radiation, resulting in a strong suppression of dissipative magnetoresistance across a wide magnetic field range as shown in Fig. 5, in agreement with experimental observations.Mani-apl ; Dor04
### III.3 150 and 280 GHz
The radiation-induced SdHO modulation can be seen clearly in the low magnetic field region $`\omega /\omega _c>1`$ with higher radiation frequency. Figure 6 shows the calculated electron temperature $`T_\mathrm{e}`$ and magnetoresistivity $`R_{xx}`$ as functions of $`\omega /\omega _c`$ for a 2D system of electron density $`N_\mathrm{e}=3.0\times 10^{15}`$ m<sup>-2</sup>, linear mobility $`\mu _0=2000`$ m<sup>2</sup>/Vs and $`\alpha =3`$, subject to a 150-GHz microwave radiation of three different amplitudes $`E_{\mathrm{i}s}=0.1,0.6`$ and 2 V/cm at a lattice temperature of $`T=0.5`$ K. Low-power microwave illumination ($`E_{\mathrm{i}s}=0.1`$ V/cm) already yields sufficient $`T_\mathrm{e}`$ oscillation with maxima at $`\omega /\omega _c=1,2,3,4`$, giving rise to clear SdHO modulations having nodes at $`T_\mathrm{e}`$ maxima. At higher microwave power ($`E_{\mathrm{i}s}=0.6`$ V/cm) when the MIMO shows up, the $`T_\mathrm{e}`$ maxima gets higher, suppressing the SdHO in the vicinities of $`\omega /\omega _c=1,2,3,4`$, but a strong amplitude modulation of SdHOs is still seen. In the case of $`E_{\mathrm{i}s}=2`$ V/cm, $`R_{xx}`$ shows strong MIMO and the electron temperature further grows so that most of SdHOs almost disappear in the range of $`\omega /\omega _c>2`$. Note that the small $`T_\mathrm{e}`$ peaks at $`\omega /\omega _c=1.5`$ and 2.5 are due to the absorption rate $`S_\mathrm{p}`$ maxima induced by two-photon processes, which gives rise to additional nodes in the SdHOs.
Another example of the SdHO modulation appearing simultaneously with MIMO is plotted in Fig. 7, where the energy absorption rate $`S_\mathrm{p}`$, the electron temperature $`T_\mathrm{e}`$, and the magnetoresistivity $`R_{xx}`$ are shown as functions of $`\omega /\omega _c`$ for a 2D system having an electron density of $`N_\mathrm{e}=3.0\times 10^{15}`$ m<sup>-2</sup>, a linear mobility of $`\mu _0=1000`$ m<sup>2</sup>/Vs, and a broadening parameter of $`\alpha =2`$, subject to linearly $`x`$-direction polarized incident microwave radiations of frequency $`\omega /2\pi =280`$ GHz and amplitude $`E_{\mathrm{i}s}=3.5`$ V/cm. The energy absorption rate $`S_\mathrm{p}`$ has broad large peaks at $`\omega /\omega _c=1,2,3,4,5`$ (due to single-photon resonant process) and small peaks at $`\omega /\omega _c=1.5,2.5`$ (due to two-photon resonant process), giving rise to the oscillation of the electron temperature $`T_\mathrm{e}`$. One can clearly see the peaks of the electron temperature $`T_\mathrm{e}`$ and the nodes of SdHO modulation at $`\omega /\omega _c=1,2,3,4`$ and 5, together with MIMOs. These are in agreement with the experimental observation reported in Ref.Kovalev, .
### III.4 Discussion
Note that in GaAs-based systems at a temperature around $`T1`$ K, LA phonons generally give larger contribution to the electron energy dissipation $`W`$ than that from TA phonons and LO phonons are usually frozen. However, in the case of high radiation power or in the vicinity of $`\omega \omega _c`$, where the resonantly absorbed energy can be relatively large and the electron temperature can rise up above 20 K, a weak emission of LO phonons takes place. Though at this temperature the number of excited LO phonons is still very small and their contribution to momentum relaxation (resitivity) is negligible in comparison with acoustic phonons, they can already provide an efficient energy dissipation because each excited LO phonon contributes a huge energy transfer of $`\mathrm{\Omega }_{\mathrm{LO}}400`$ K. With a continuing rise of electron temperature the LO-phonon contribution increases rapidly. This effectively prevents the electron temperature from going much higher than 20 K, such that the $`T_\mathrm{e}`$-vs-$`\omega _c/\omega `$ curve of large incident microwave power in Fig.1 exhibits a flat top around $`\omega _c/\omega =1`$.
In this paper, we did not consider the role of surface or interface phonons in the GaAs heterostructure. Depending on sample geometry, the surface phonons may be important in dissipating electron energy thus decreasing the electron temperature.
Acknowledgements
We thank Dr. V.I. Ryzhii, Dr. R.G. Mani and Dr. R.R. Du for helpful discussions. This work was supported by Projects of the National Science Foundation of China, the Special Funds for Major State Basic Research Project, and the Shanghai Municipal Commission of Science and Technology. |
warning/0506/q-bio0506017.html | ar5iv | text | # Built to evolve
## 1 Introduction
We consider the probabilities of obtaining genomes by random mutations and natural selection, and obtain the conditions that are required for successfull evolution. Examples are presented for viruses, bacteria, eukaryote cells, and multi-cellular organisms. Our conclusions are collected in Section 10.
## 2 The odds
The human genome has about 30 thousand useful genes, each with an average coding region of 20 thousand base pairs coding about 6666 aminoacids. What is the probability of writing a specific sequence of $`30000\times 6666`$ words (aminoacids) chosen at random from a list of 20? The answer is $`20^{30000\times 6666}10^{260000000}`$, i.e. zero for all practical purposes. For that matter, what is the probability of getting the correct genome of any one of the 30 million species, with one try every second since the Big Bang, by every one of the $`10^{31}`$ (or so) bacteria on Earth? The answer is $`3\times 10^7\times 410^{17}\times 10^{31}\times 20^{30000\times 6666}`$ which is still the same result: $`10^{260000000}`$ (only the last two digits in the exponent are changed), i.e. zero for all practical purposes.
Note that the โsimplestโ living organism, the prokaryote bacteria, is almost as complex as a human being: it has of order 1000 active genes with an average of $`1500`$ base pairs! There are many missing links between organic molecules (sugars, lipids, bases, aminoacids, etc) and the simplest forms of life.
Perhaps the trick is to write little pieces of genetic code at a time. So, let us turn the question around: What is the largest gene that can be produced at random with a finite probability (say, $`10^4`$) with one try every second during 100 million years by each of the $`10^{31}`$ bacteria on Earth? The answer is one gene encoding, at most, $`39`$ aminoacids, or $`117`$ base pairs. This is roughly the limit for undirected random evolution.
## 3 Evolution
Let us play a game called โEvolutionโ. We sit in front of a key board and hit keys at random. The probability of obtaining โRomeo and Julietโ is zero for all practical purposes. In fact, the probability of obtaining any meaningfull novel in any known language is zero for all practical purposes. Now introduce โmutationsโ, i.e. replace random letters by random hits of the keyboard. Still no meaningfull novel will ever be obtained for all practical purposes.
Now suppose we are allowed to select which letter (or small set of letters) to mutate at a rate much higher than the background mutation rate of the other letters, and we are allowed to stop the hypermutations when a particular outcome is obtained. For example, choose to hypermutate the 10th letter until โeโ is obtained, or choose to hypermutate the 10th, 11th and 12th letters until โdogโ is obtained. Of course, now we can write any novel at all: the game has become trivial. If the choice of which letter to hypermutate is perfectly specific, and the choice of which outcome to select is perfectly specific, the game becomes trivial, i.e. the outcome becomes certain, even tho the keys are hit at random. If the environment were perfectly specific, it would have perfect control over evolution (even if mutations occur at random).
Evolution lies somewhere in between. The choice of which bases on which genes to hypermutate is not perfectly specific and is incomplete, and the outcome that stops the hypermutations is also not perfectly specific. However, with enough specificity it is possible to write little pieces of survivable genome. Then the pieces can be combined. In our example we could copy โdogโ, reverse โdogโ, concatenate โdogโ and โcatโ, interchange โdogโ and โcatโ (as in horizontal gene transfer, or in meiosis), introduce words and phrases from other books, and so on. The interchange of genetic material within, and between, bacteria, viruses and eukaryote cells plays a major role in evolutionary change.
Survivable pieces of the genome will often involve negative (stabilizing) feedback loops. A hierarchy of negative feedback loops, within negative feedback loops, within โฆ, is self organizing.
The steps in evolution, from simple to complex, might have been as shown in Figure 1. At all levels of complexity the environment must have directed evolution with sufficient specificity for the steps to have a non-zero probability.
## 4 Adaptive mutations
Consider the (simplified) metabolic pathway of a cell shown in Figure 2. Precursor A is converted to end-product B by enzyme C. Enzyme C is encoded by gene 1. The transcription of gene 1 (and synthesis of enzyme C) is repressed (directly or indirectly) by the end-product B. This negative feedback loop regulates the concentration of B, and is efficient because enzyme C is produced only when needed. This feedback loop therefore has a value for survival, and would have been selected by nature. Now starve the cell of precursor A in the presence of a similar precursor A. The result is a reduction of the concentration of the end-product B, and a derepression of gene 1. The rate of transcription of gene 1 increases (by a factor that can exceed 1000). During transcription, mRNA copies one strand of DNA exposing the other strand. Single strand DNA is prone to mutations due to the lack of hydrogen bond stabilization between complementary bases, the formation of loops and other secondary structures, <sup>1</sup><sup>1</sup>1Segments of the single strand DNA may stick to other segments with mostly complementary bases, resulting in unpaired or mispaired bases. Unpaired bases are prone to deamination, deletion or replacement. Cytosine deaminates to uracil at a rate 100 times larger in single strand DNA than in double strand DNA. Mutations also occur in the end-loops where bases have no complement, and in the stem. These errors are immortalized during DNA duplication or repair. and supercoiling. As a result, gene 1 acquires a high rate of mutations and begins synthesizing enzymes similar to C, until one of them is able to convert precursor A into end-product B. Gene 1 is then repressed by end-product B and hypermutation stops. The resulting negative feedback loop resumes control of the concentration of the end-product B to the same original concentration. The net result of these processes is that a change of the environment (the starvation for precursor A) triggers mutations of a specific gene of the cell, until the cell is capable of substituting precursor A by precursor A. So the environment can direct evolution in very specific ways.
Let us briefly describe examples.
## 5 B-lymphocytes
Let us consider B-lymphocytes of the immune system (see Figure 3), which have been studied in considerable detail. These B-lymphocytes have antibody proteins (called immunoglobulins) attached to their membrane. An invading bacteria has antigen proteins attached to its membrane. The antibody can bind to a very specific set of antigens. This binding triggers a series of complex steps (including helper T-lymphocytes) that activate mitosis of the B-cell, expresses several genes that code immunoglobulins, and differentiates the B-lymphocytes into antibody secreting cells and memory cells. Proliferanting B-cells in germinal centers show high rates of point mutations in genes coding immunoglobulins ($`10^3`$ to $`10^6`$ times higher than the spontaneous rate of other genes). The result is the synthesis of immunoglobulins with small differences. The binding of these antibodies to the antigens (with the intervention of follicular dendritic cells) produce signals that rescue the B-lymphocytes from programmed cell death. At the latter stages of the infection the concentration of invading bacteria becomes low, so only B-cells producing very high affinity antibodies can bind to the antigens and survive. This phenomenon is called โaffinity maturationโ.
The net result of these complex processes is that a change of the environment (the invading bacteria) triggers hypermutations of specific sections (those that code for the binding site of the immunoglobulins) of specific genes of specific B-cells, and selects mutations producing immunoglobulins with the highest affinity to the antigen. Note that each one of these steps is very specific.
Since B-lymphocytes are somatic, the selected mutations are not passed on to the next generation, so this is an example of evolution of B-cells in one individual, not evolution of the species.
The hypermutation associated to affinity maturation of B-cells is caused by the induction of mutagenic genes such as cytidine deaminase (which causes C to U transitions) and error prone DNA polymerases. The presence of the same genes in other eukaryots may indicate that similar hypermutagenic processes may occur in eukarya, and, as in prokaryots or B-cells, the induction of these genes may be triggered by environmental stress such as starvation.
Let us mention that B-lymphocytes with hypermutation and recombination of V, D and J cassettes of genes coding immunoglobulins, produce B-cells that synthesize of order $`10^{11}`$ different antibodies capable of binding to as many different antigens. So, with just a handfull of genes in the genome, B-cells are able to synthesize a much larger number of different proteins!
## 6 Enterobacter arogenes
Let us briefly describe the metabolic pathway of Enterobacter arogenes shown in Figure 4. In the wild strain 5P14, ribitol induces gene 1 to synthesize ribitol dehydrogenase. This enzyme metabolizes ribitol, or, with low specific activity, xylitol. The wild bacteria can therefore metabolize xylitol only if ribitol is present. If the bacteria is starved for ribitol in the presence of xylitol, a mutation occurs in a gene 2 that causes the expression of gene 1 even in the absence of ribitol. This strain, called X1, appears in 4.1 hours. Mutations in gene 1 produce strains X2 in 1.7 hours, and then X3 in 0.9 hours. These strains synthesize modified enzymes with increasing specific activity on xylitol.
So a change of the environment (the removal of ribitol in the presence of xylitol) results in specific mutations of two specific genes. The specificity is so great that the experiment is repeatable!
## 7 Sex
Mate the largest dogs of different litters for several generations, and you end up with huge Great Danes. This is artificial selection. A Great Dane has many โbig-dog-genesโ. These big-dog-genes were already in the genetic pool of the population of dogs. The largest dog of a litter probably has more big-dog-genes than each parent due to the mixing of genes during sexual reproduction (meiosis). The smallest dog of the litter probably has less big-dog-genes than each of the parents.
Natural selection can work in a similar way. An ecological niche attracts individuals specially adapted to that niche. <sup>2</sup><sup>2</sup>2Mimetism is an example: a green insect that chooses a green environment to hide in has a better chance of survival. These individuals have an enhanced probability of mating with each other. If the population has a gene pool with several genes that favor the niche, then, after a few generations, these genes can come together and we obtain individuals specially adapted to the niche. If these individuals no longer mate and reproduce with the general population, then a sub-species has formed.
Even cultural preferences can bias mating, resulting in genomes specially adapted to these cultural preferences. This is known as the Baldwin effect, after the description of this phenomenon by James Mark Baldwin in 1896.
## 8 Viruses
A virus is composed of genetic material packaged (mostly) in proteins. The genetic material may be linear or circular, single or double stranded, haploid or diploid, monopartite or multipartite, DNA and/or RNA. The RNA can be โpositiveโ and serve as a messenger RNA to directly synthesize proteins (using the tRNA and ribosomes of the host cell), or it can be โnegativeโ and require a transcription to +RNA before protein synthesis.
What proteins are coded by the DNA or RNA? In order to reproduce, the virus must code the proteins that form part of the virus itself: structural proteins and enzymes needed prior to protein synthesis (such as enzymes used by retroviruses to synthesize DNA from RNA, enzymes used by negative strand RNA viruses to transcribe -RNA to +RNA, enzymes used by double stranded RNA viruses to make single strands, etc). Depending on the type of virus, other proteins may be coded as well. Examples are enzymes to transcribe +RNA to -RNA, DNA to RNA at various starting sites, proteins that block defenses of the host cell, proteins that cleave other proteins at special sites (so one mRNA of the virus can code many proteins linked together at cleaving sites).
In addition, and of particular interest for evolution, the virus may encode proteins that can turn on or turn off hypermutations, and enzymes used to recombine RNA within the virus, among different viruses (even of different species), and between the virus and the host cell. Let us quote from : โThe two major forces acting upon viral genomes to generate diversity that can be tested for environmental survival and replicative fitness are mutations and recombination. Some viruses have a good deal of control over their own rates of mutation and even the frequency of recombination. They exert control by encoding viral enzymesโ for replicative and recombinational functions.
## 9 Other examples
A strain (known as FC40) of Escherichia coli can not digest lactose due to a frameshift mutation in gene lacZ that does not allow the synthetis of $`\beta `$-galactosidase in sufficient quantity. It has been observed that this frameshift mutation undergoes reversion when lactose becomes the sole source of energy. It is important to note that most reversions occur after exposure to lactose.
Starvation of Escherichia coli induces the production of alternative polymerase enzymes (DinB and UmuD$`{}_{2}{}^{}{}_{}{}^{}`$C) which are capable of replicating badly damaged sequences of DNA, and, in the process, produce high rates of mutations. Homologs of these alternative enzymes have been found in Saccharomyces cerevisiae, mice and humans.
Escherichia coli is able to mutate even when not dividing or replicating its DNA, and these mutations may be its main source of genetic variation.
Some plants switch from asexual proliferation (rhizomes) to sexual reproduction in conditions of stress. In doing so they speed up evolution by trying new combinations of genetic material in the process of meiosis.
Snails switch from hermaphrodite reproduction to bi-sexual reproduction in conditions of stress due to parasites. This strategy speeds up evolution when needed.
## 10 Conclusions
Evolution appears to be hopelessly improbable unless random mutations are limited to no more than about 100 bases of specific genes, and the selection of the outcomes are sufficiently specific. We therefore propose that all living organisms, in addition to being self-organizing and reproducing (autopoyetic), are built to evolve in selective ways. It appears that viruses, prokaryots and eukaryots have considerable control over the rates and spectrum of mutations and recombinations. There are built-in mechanisms to control hypermutations of selective regions of selective genes, and sufficiently selective mechanisms to choose the outcome of these mutations. The high selectivity of these and other mechanisms are required for evolution to be successfull. We have given examples of viruses, prokaryots, eukaryots, and multicellular organisms, where this is indeed the case. |
warning/0506/hep-th0506012.html | ar5iv | text | # Supertube Dynamics in Diverse Backgrounds
## 1 Introduction
The formulation of string theory in various time-dependent and cosmological backgrounds remains as one of the most important and exciting open problems for theoretical physicists. In doing so, one expects a better understanding of some of the outstanding problems and confusions of quantum gravity and cosmology. In an attempt along this direction, Sen proposed a rolling tachyon solution โthe decay of the unstable D-brane (brane-anti brane pair) when the tachyon rolls in the valley into the bottom of its potential (vacuum without D-branes). It has also been argued in that the final state of such a decay process leads to a classical matter state, with the equation of motion of a non-interacting and pressureless dust known as โtachyon materโ with no obvious open string excitations. These solutions, in general, are constructed by perturbing the boundary conformal field theory that describes the D-brane by an exact marginal deformation. The real time tachyon dynamics shows that the effective field theory for Dirac-Born-Infeld type captures surprisingly well many aspects of rolling tachyon solutions of full string theory. (See and references therein for a detailed review of the tachyon dynamics)
Recently, Kutasov gave a โgeometric realizationโ of the open string tachyon, in the form of the rolling D-braneโthe time dependent dynamics of the D-brane in the throat of NS5-branes. It has been shown that the decay process resembles astonishingly, that of the decay of unstable D-branes, when we restricts ourselves to the case when the distance between the D-brane and NS5-brane is of the order of the string length $`(l_s)`$. One wonders whether the above identification is an artifact of low energy effective field theory or it can be extended to the full string theory. This was analyzed in in the form of the boundary states of the in falling D-branes into the throat region of a stack of NS5-branes (the CHS geometry). The electric deformation of the process has also been realized in terms of the decay of the electrified D-brane into the NS5 throat. This has further been extended in to the case of D$`p`$-brane background. Later on in , this ideas were extended to the case of the Non-BPS probe branes and a careful analysis revealed the existence of a symmetry (possibly broken explicitly at the Lagrangian level ). Hence a natural guess will be that the dynamics of the D-branes in the curved backgrounds has a much more richer structure in the geometry and in the dynamics, and hence likely to add some input in our understanding of string theory in time dependent curved backgrounds. See ( \- ) for related studies of the D-brane dynamics and cosmological applications in various backgrounds.
In another context, the supertubeโthe bunch of straight strings with D0-brane blown up into a supersymmetric tubular D2-brane with electric and magnetic fieldโ has been instrumental in understanding black hole physics. For example, the quantum states of a supertube counted from the direct quantization of the Born-Infeld action near the geometrically allowed microstates with some fixed charges are shown to be in one to one correspondence with some black holes. The stability of this bound state is achieved by the non-zero angular momentum generated by the Born-Infeld electric and magnetic charges. Recently the stability of the supertubes in other curved backgrounds has also been studied in . Hence in the recent surge of studying the D-brane dynamics in various backgrounds, in connection with finding out the time dependent solutions of string theory and tachyon dynamics, it seems very interesting to study the dynamics of the supertube in diverse backgrounds. But looking at the rather complicated analysis of the problem in the usual effective theory approach, we would like to take advantage of the Hamiltonian formulation. This approach has been very instructive in investigating the D-branes in the strong coupling limit, by formally taking the zero tension limit, for example. By doing this we achieve a formal Hamiltonian for the D-brane motion in curved background, and the beauty of this formalism makes us comfortable to use the constraints of equations of motion in appropriate ways. We try to be as general as possible in the beginning, but while studying the motion of the tube in the background of various macroscopic objects, we use some properties of the supergravity backgrounds relevant for studying string theory, to make the analysis simpler.
The layout of the paper is as follows. In section-2, we start by deriving the Dirac-Born-Infeld action using the Hamiltonian formulation. The rather complicated action can be made simpler by making use of an appropriate gauge fixing. After deriving an action for the dynamics of the tube in rather general curved backgrounds, in section-3, we focus our attention to the case of Dp-branes, NS5-branes and the fundamental string backgrounds. We study in detail the effective potential and the motion of the tube in the vicinity of these backgrounds generated by the macroscopic objects. We present our conclusions in section-4.
## 2 Hamiltonian dynamics for Dp-brane in general background and its gauge fixed version
This section is devoted to the formulation of the Hamiltonian formalism for Dp-brane in general curved background. We will mainly follow the very nice analysis presented in . The Dirac-Born-Infeld action for a Dp-brane in a general bosonic background is given by the following usual form<sup>1</sup><sup>1</sup>1We do not consider the Wess-Zummino term since in the examples studied below the coupling to Ramond-Ramond fields is not important.
$$S_p=\tau _pd^{p+1}\xi e^\mathrm{\Phi }\sqrt{det๐},$$
(1)
where
$$๐_{\mu \nu }=\gamma _{\mu \nu }+F_{\mu \nu },\mu ,\nu =0,\mathrm{},p,$$
(2)
and where $`\tau _p`$ is Dp-brane tension. The induced metric $`\gamma _{\mu \nu }`$ and the induced field strengths on the worldvolume $`F_{\mu \nu }`$ are given by the following expressions
$`\gamma _{\mu \nu }`$ $`=`$ $`_\mu X^M_\nu X^Ng_{MN},`$ (3)
$`F_{\mu \nu }`$ $`=`$ $`_\mu X^M_\nu X^Nb_{MN}+_\mu A_\nu _\nu A_\mu .`$ (5)
Where $`g_{MN}(X(\xi )),b_{MN}(X(\xi )),M,N=0,\mathrm{},9`$ are metric and NS-NS two form field that are in general functions of the embedding coordinates $`X^M(\xi )`$. Let us now calculate the conjugate momenta from (1)
$`P_M(\xi )`$ $`=`$ $`{\displaystyle \frac{\delta S}{\delta _0X^M(\xi )}}={\displaystyle \frac{e^\mathrm{\Phi }}{2}}\sqrt{det๐}_\nu X^N\left(g_{MN}(๐^1)^{(\nu 0)}+b_{MN}(๐^1)^{[\nu 0]}\right),`$ (6)
$`\pi ^a(\xi )`$ $`=`$ $`{\displaystyle \frac{\delta S}{\delta _0A_a(\xi )}}={\displaystyle \frac{e^\mathrm{\Phi }}{2}}\sqrt{det๐}(๐^1)^{[a0]},a,b=1,\mathrm{},p,`$ (8)
$`\pi ^0(\xi )`$ $`=`$ $`{\displaystyle \frac{\delta S}{\delta _0A_0(\xi )}}={\displaystyle \frac{e^\mathrm{\Phi }}{2}}\left((๐^1)^{00}(๐^1)^{00}\right)\sqrt{det๐}=0,`$
where we have introduced the symmetric and antisymmetric form of any (p+1)$`\times `$(p+1) matrix $`\mathrm{A}`$ as
$$\mathrm{A}^{(\mu \nu )}=\mathrm{A}^{\mu \nu }+\mathrm{A}^{\nu \mu },\mathrm{A}^{[\mu \nu ]}=\mathrm{A}^{\mu \nu }\mathrm{A}^{\nu \mu }.$$
(11)
For our purpose it is useful to define the conjugate momenta $`\mathrm{\Pi }_M`$ as
$$\mathrm{\Pi }_M=P_M+P^aB_{aM}=\frac{e^\mathrm{\Phi }}{2}V\sqrt{det๐}\gamma _{M\nu }(๐^1)^{(\nu 0)},$$
(12)
where
$$\gamma _{M\nu }=g_{MN}X_\nu ^N,B_{aM}=_aX^NB_{NM}.$$
(13)
Using these expressions it is easy to see that following constraints are obeyed
$`\mathrm{\Pi }_M_aX^M+\pi ^aF_{ab}=0,`$
$`\mathrm{\Pi }_Mg^{MN}\mathrm{\Pi }_N+\pi ^a\gamma _{ab}\pi ^b+e^{2\mathrm{\Phi }}\tau _p^2det๐_{ab}=0,`$
$`\pi ^0=0.`$
The corresponding Hamiltonian for Dp-brane in curved background takes the following form
$`H={\displaystyle d^p\xi (\xi )},`$ (15)
where the Hamiltonian density $`()`$ is given by
$``$ $`=`$ $`\pi ^i_iA_0+\sigma \pi ^0+\rho ^a(\mathrm{\Pi }_M_aX^M+F_{ab}\pi ^b)`$ (16)
$`+`$ $`\lambda (\mathrm{\Pi }_Mg^{MN}\mathrm{\Pi }_N+\pi ^a\gamma _{ab}\pi ^b+e^{2\mathrm{\Phi }}\tau _p^2det๐_{ab}),`$ (18)
where $`\sigma ,\rho ^a`$ and $`\lambda `$ are Lagrange multipliers for the constraints. More precisely, the final Hamiltonian is just the sum of constraints, in agreement with the diffeomorphism invariance of the original Lagrangian.
### 2.1 Gauge fixing problem
Since the original DBI action is diffeomorphism invariant, it is convenient to use this symmetry to reduce the number of independent equations of motions. This procedure is commonly known as gauge fixing.
To proceed we can without loss of generality presume that the metric takes diagonal form. Then we consider the equation of motion for $`X^M`$
$`_0X^M={\displaystyle \frac{\delta H}{\delta P_M(\xi )}}=\rho ^a_aX^M+2\lambda g^{MN}\mathrm{\Pi }_N.`$
Usually the static gauge is imposed at the level of the action and Lagrangian. In our case, however, we would like to perform the similar procedure at the level of canonical equations of motion. Since in the Lagrangian formalism the static gauge is imposed by the constraints $`X^\mu =\xi ^\mu ,\mu =0,\mathrm{},p`$, it is natural to solve the canonical equations of motion for $`X^\mu `$ by the ansatz
$$_0X^0=1,_aX^b=\delta _a^b,a=1,\mathrm{},p.$$
(20)
Then the equations of motion for $`X^\mu `$ simplify as
$`\mathrm{\Pi }_0={\displaystyle \frac{g_{00}}{2\lambda }},\mathrm{\Pi }_a={\displaystyle \frac{\rho ^ag_{aa}}{2\lambda }}(\mathrm{No}\mathrm{summation}\mathrm{over}\mathrm{a}).`$
Using these results we get
$$\mathrm{\Pi }_M_aX^M=\frac{\rho ^bg_{ba}}{2\lambda }+\mathrm{\Pi }_s_aX^s$$
(22)
where $`s,r,t\mathrm{}`$ label the directions transverse to Dp-brane. If we ignore the terms in the Hamiltonian that enforce Gauss law constraints and the vanishing of $`\pi ^0`$ we obtain <sup>2</sup><sup>2</sup>2To see this, note that the equation of motion for $`A_0`$ leads to the โGauss lawโ constraints
$$\frac{\delta H}{\delta A_0(๐ฑ)}=0_a\pi ^a=0$$
(23) and variation of the Hamiltonian with respect to $`\sigma `$ implies $`\pi ^0=0`$.
$`={\displaystyle \frac{g_{00}}{4\lambda }}{\displaystyle \frac{\rho ^ag_{ab}\rho ^b}{4\lambda }}+\rho ^ab_a+\lambda D,`$
where
$`b_a=\mathrm{\Pi }_s_aX^s+F_{ab}\pi ^b,`$
$`D=\mathrm{\Pi }_rg^{rs}\mathrm{\Pi }_s+\pi ^a\gamma _{ab}\pi ^b+e^{2\mathrm{\Phi }}\tau _p^2det๐_{ab}.`$
In order to enforce the variation with the Lagrange multiplies we should consider the Dp-brane action
$`S_p={\displaystyle d^{p+1}\xi }`$ (26)
where we express $``$ as a Lagrange transformation of the hamiltonian density (2.1) so that we obtain the action in the form
$$S_p=d^{p+1}\xi \left(P_M_0X^M+\pi _a_0A^a\right)$$
(27)
Now we replace $`P_\mu `$ with the help of $`\mathrm{\Pi }_\mu `$ given above. Using the fact that $`_0X^\mu =\delta _0^\mu `$ we obtain an action for $`P_s,X^s`$ that contains the Lagrange multipliers $`\lambda ,\rho ^a`$
$`S_p={\displaystyle d^{p+1}\xi \left(P_s_0X^s+\pi ^a_0A_a\pi ^aB_{a0}+\frac{g_{00}}{4\lambda }+\frac{\rho ^ag_{ab}\rho ^b}{4\lambda }\lambda D\rho ^ab_a\right)}.`$
Now the variation with respect to $`\rho ^a`$ implies ( remember that $`g_{ab}`$ is diagonal)
$$\rho ^a=2\lambda g^{ab}b_b.$$
(29)
If we then insert (29) back into the action we get
$$S=d^{p+1}\xi \left(P_s_0X^s+\pi ^a_0A_a\pi ^aB_{a0}+\frac{g_{00}}{4\lambda }\lambda (D+b_ag^{ab}b_b)\right).$$
(30)
Now the variation with respect to $`\lambda `$ implies
$$\lambda =\frac{1}{2}\sqrt{\frac{g_{00}}{b_ag^{ab}b_b+D}}.$$
(31)
Inserting back to the action we get finally
$$S=d^{p+1}\xi \left(P_s_0X^s+\pi ^a_0A_a\sqrt{g_{00}}\sqrt{b_ag^{ab}b_b+D}\pi ^aB_{a0}\right)$$
(32)
from which we obtain the Hamiltonian density for transverse variables $`Y^s,P_s`$ and gauge field $`p^a,A_a`$ in the form
$``$ $`=`$ $`\sqrt{g_{00}}\sqrt{๐ฆ}+\mathrm{\Pi }^aB_{a0},`$ (33)
$`๐ฆ`$ $`=`$ $`\mathrm{\Pi }_rg^{rs}\mathrm{\Pi }_s+\pi ^a\gamma _{ab}\pi ^b+b_ag^{ab}b_b+e^{2\mathrm{\Phi }}\tau _p^2det๐_{ab},`$ (35)
$`\mathrm{\Pi }_s`$ $`=`$ $`P_s+\pi ^aB_{as},b_a=\mathrm{\Pi }_s_aX^s+F_{ab}\pi ^b.`$ (37)
Now we are ready, using the Hamiltonian density (33) to study the dynamics of the supertube in the various supergravity background.
## 3 Supertube in Diverse Backgrounds
By standard treatment, the supertube in flat spacetime is the D2-brane that is stretched in one particular direction, say $`z`$ direction and that have arbitrary shape in the transverse $`R^8`$ space. In the case when we study the supertube dynamics in the background of some macroscopic objects (Dp-branes, NS5-branes or fundamental strings) the situation is slightly more complicated. On the other hand since these backgrounds have generally manifest rotational invariance of the transverse $`R^{9k}`$ space, where $`k`$ is the spatial dimension of given objects, it is natural to simplify the analysis by presuming that the supertube has its base in the $`(x^8,x^9)`$-plane <sup>3</sup><sup>3</sup>3 We restrict ourselves to the case of background Dp-branes with $`p<5`$ so that the coupling of the D2-brane to the background Ramond-Ramond field vanishes.. In this plane we introduce the polar coordinates
$$x^8=R\mathrm{cos}\varphi ,x^9=R\mathrm{sin}\varphi .$$
(38)
Then the embedding coordinates are $`R(\sigma ,t),X^m,m=p+1,\mathrm{},7`$, that parameterize the position of D2-brane in the directions transverse to given macroscopic objects of spatial dimension $`p`$ which are however transverse to the $`(x^8,x^9)`$ plane. We also introduce the coordinates $`Y^u,(u,v=2,\mathrm{},p)`$, where $`Y^u`$ parameterize the position of supertube in the space parallel with Dp-brane. Since the $`B`$ field is also zero, the Hamiltonian for such a Dp-brane takes the form
$$H=๐t๐z๐\sigma =๐t๐z๐\sigma \left[\sqrt{g_{_{00}}}\sqrt{๐ฆ}\right],$$
(39)
with
$`๐ฆ`$ $`=`$ $`\pi _ag^{ab}\pi _b+p__Rg^{RR}p__R+p_mg^{mn}p_n+p_ug^{uv}p_v+b_ag^{ab}b_b`$ (40)
$`+`$ $`(\pi ^a_aR)g_{_{RR}}(\pi ^b_bR)+(\pi ^a_aY^u)g_{_{uv}}(\pi ^b_bY^v)`$ (42)
$`+`$ $`(\pi ^a_aX^m)g_{mn}(\pi ^b_bX^n)+e^{2\mathrm{\Phi }}\tau _2^2det๐_{ab},`$ (44)
and
$$๐_{ab}=g_{_{ab}}+g_{{}_{R}{}^{}_R}_aR_bR+g_{mn}_aX^m_bX^n+g_{_{uv}}_aY^u_bY^v+F_{_{ab}},a,b=z,\varphi $$
(45)
with
$$b_a=F_{_{ab}}\pi ^b+_aRp__R+_aX^mp_m+_aY^up__u.$$
(46)
Using the Hamiltonian (39) the equation of motions for $`A_a`$ and $`\pi ^a`$ take the form
$$\dot{A}_a(๐ฑ)=\frac{\delta H}{\delta \pi ^a(๐ฑ)}=\sqrt{g_{00}}\frac{g_{ab}\pi ^b+_aX^mg_{mn}(\pi ^b_bX^n)+_aY^ug_{uv}(\pi ^b_bY^v)+F_{ab}g^{bc}b_c}{\sqrt{๐ฆ}},$$
(47)
$`\dot{\pi }^a(๐ฑ)`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta A_a(๐ฑ)}}=_c\left[{\displaystyle \frac{\sqrt{g_{00}}\pi ^ag^{cb}b_b}{\sqrt{๐ฆ}}}\right]_c\left[{\displaystyle \frac{\sqrt{g_{00}}\pi ^cg^{ab}b_b}{\sqrt{๐ฆ}}}\right]`$
$`+`$ $`{\displaystyle \frac{1}{2}}_c\left[{\displaystyle \frac{e^{2\mathrm{\Phi }}\tau _2^2\sqrt{g_{00}}}{\sqrt{๐ฆ}}}(๐^1)^{ac}det๐_{ab}\right]{\displaystyle \frac{1}{2}}_c\left[{\displaystyle \frac{e^{2\mathrm{\Phi }}\tau _2^2\sqrt{g_{00}}}{\sqrt{๐ฆ}}}(๐^1)^{ca}det๐_{ab}\right].`$
Further, the equations of motion for the embedding coordinates are
$`\dot{R}(๐ฑ)`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta p_R(๐ฑ)}}={\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}\left(g^{RR}p_R+_iRg^{ab}b_b\right),`$ (51)
$`\dot{X}^m(๐ฑ)`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta p_m(๐ฑ)}}={\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}\left(g^{mn}p_n+_aX^mg^{ab}b_b\right),`$ (53)
$`\dot{Y}^u(๐ฑ)`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta p_u(๐ฑ)}}={\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}\left(g^{uv}p_v+_aX^ug^{ab}b_b\right),`$ (55)
while the equation of motion for $`p_m,p_u`$ and $`p_R`$ are
$`\dot{p}_m\left(๐ฑ\right)={\displaystyle \frac{\delta H}{\delta X^m\left(๐ฑ\right)}}=_a\left[{\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}\pi ^ag_{mn}\left(\pi ^b_bX^n\right)\right]`$ (56)
$`+`$ $`_a\left[{\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}e^{2\mathrm{\Phi }}\tau _2^2g_{mn}_bX^n\left(๐^1\right)^{ba}det๐_{ab}\right]+{\displaystyle \frac{\delta g_{_{00}}}{\delta X^m}}{\displaystyle \frac{\sqrt{๐ฆ}}{2\sqrt{g_{_{00}}}}}`$ (58)
$``$ $`{\displaystyle \frac{\sqrt{g_{_{00}}}}{2\sqrt{๐ฆ}}}(\pi ^a{\displaystyle \frac{\delta g_{ab}}{\delta X^m}}\pi ^b+p_R{\displaystyle \frac{\delta g^{RR}}{\delta X^m}}p_R+p_n{\displaystyle \frac{\delta g^{no}}{\delta X^m}}p_o+p_u{\displaystyle \frac{\delta g^{uv}}{\delta X^m}}p_v+{\displaystyle \frac{\delta \left[e^{2\mathrm{\Phi }}\right]}{\delta R}}\tau _2^2det๐_{ab}`$
$`+`$ $`e^{2\mathrm{\Phi }}\tau _2^2[{\displaystyle \frac{\delta g_{_{RR}}}{\delta X^m}}\left(_aR_bR\right)+{\displaystyle \frac{\delta g_{no}}{\delta X^m}}\left(_aX^n_bX^o\right)+{\displaystyle \frac{\delta g_{uv}}{\delta X^m}}\left(_aY^u_bY^v\right)]\left(๐^1\right)^{ba}det๐_{ab}),`$
$`\dot{p}_u(๐ฑ)={\displaystyle \frac{\delta H}{\delta Y^u(๐ฑ)}}`$ $`=`$ $`_a\left[{\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}\pi ^ag_{uv}(\pi ^b_bY^v)\right]`$
$`+`$ $`_a\left[{\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}e^{2\mathrm{\Phi }}\tau _2^2g_{uv}_bY^v(๐^1)^{ba}det๐_{ab}\right]`$
and
$`\dot{p}__R\left(๐ฑ\right)={\displaystyle \frac{\delta H}{\delta R\left(๐ฑ\right)}}=_a\left[{\displaystyle \frac{\sqrt{g_{_{00}}}}{\sqrt{๐ฆ}}}\pi ^ag_{_{RR}}\left(\pi ^b_bR\right)\right]`$ (66)
$`+`$ $`_a\left[{\displaystyle \frac{\sqrt{g_{_{00}}}}{\sqrt{๐ฆ}}}e^{2\mathrm{\Phi }}\tau _2^2g_{_{RR}}_bR\left(๐^1\right)^{ba}det๐_{ab}\right]+{\displaystyle \frac{\delta g_{_{00}}}{\delta R}}{\displaystyle \frac{\sqrt{๐ฆ}}{2\sqrt{g_{_{00}}}}}`$ (68)
$``$ $`{\displaystyle \frac{\sqrt{g_{_{00}}}}{2\sqrt{๐ฆ}}}(\pi ^a{\displaystyle \frac{\delta g_{ab}}{\delta R}}\pi ^b+p_R{\displaystyle \frac{\delta g^{RR}}{\delta R}}p_R+p_m{\displaystyle \frac{\delta g^{mn}}{\delta R}}p_n+p_u{\displaystyle \frac{\delta g^{uv}}{\delta R}}p_v+{\displaystyle \frac{\delta \left[e^{2\mathrm{\Phi }}\right]}{\delta R}}\tau _2^2det๐_{ab}`$
$`+`$ $`e^{2\mathrm{\Phi }}\tau _2^2[{\displaystyle \frac{\delta g_{_{RR}}}{\delta R}}\left(_aR_bR\right)+{\displaystyle \frac{\delta g_{mn}}{\delta R}}\left(_aX^m_bX^n\right)+{\displaystyle \frac{\delta g_{uv}}{\delta R}}\left(_aY^u_bY^v\right)]\left(๐^1\right)^{ba}det๐_{ab})`$
Note that generally the metric is function of the expression $`R^2+X^mX_m`$. On the other hand thanks to the symmetry of the problem with respect to $`z`$ direction it is natural to consider the modes that are $`z`$ independent. We also take the ansatz for the gauge field in the form
$$F=_0A_z(\varphi )dtdz+B(\varphi )dzd\varphi ,A_\varphi =Bz.$$
(73)
Then the matrix $`๐_{ab}`$ takes the form
$$๐_{ab}=\left(\begin{array}{cccc}g_{_{zz}}& & & B\\ B& & & g_{_{\varphi \varphi }}+g_{_{RR}}(_\varphi R)^2+g_{mn}_\varphi X^m_\varphi X^n+g_{_{uv}}_\varphi Y^u_\varphi Y^v\end{array}\right)$$
(74)
so that
$$det๐_{ab}=g_{_{zz}}\left(g_{_{\varphi \varphi }}+g_{_{RR}}(_\varphi R)^2+g_{mn}_\varphi X^m_\varphi X^n+g_{_{uv}}_\varphi Y^u_\varphi Y^v\right)+B^2$$
(75)
and
$$(๐^1)^{ab}=\frac{1}{det๐}\left(\begin{array}{cccc}g_{_{\varphi \varphi }}+g_{_{RR}}(_\varphi R)^2+g_{mn}_\varphi X^m_\varphi X^n+g_{_{uv}}_\varphi Y^u_\varphi Y^v& & & B\\ B& & & g_{_{zz}}\end{array}\right).$$
(76)
With this notation the equation of motion for $`\pi ^z`$ takes the form
$$\dot{\pi }^z=_\varphi \left[\frac{\sqrt{g_{_{00}}}}{\sqrt{๐ฆ}}g^{\varphi \varphi }(B\pi ^z+_\varphi X^mp__m+_\varphi Y^up_u)\right]+__\varphi \left[e^{2\mathrm{\Phi }}\tau _2^2\frac{\sqrt{g_{_{00}}}B}{\sqrt{๐ฆ}}\right]$$
(77)
As it is clear from the equation above $`\pi ^z`$ is generally time dependent. On the other hand we can show that the quantity $`n=๐\varphi \pi ^z`$ is conserved since
$$\dot{n}=๐\varphi \dot{\pi }^z=๐\varphi _\varphi [\mathrm{}]=0.$$
(78)
The similar case occurs for $`p_u`$ where $`\dot{p}_u(๐ฑ)0`$
$$\dot{p}_u=_\varphi \left[\frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}e^{2\mathrm{\Phi }}\tau _2^2g_{uv}(_\varphi Y^v)g_{zz}\right]$$
(79)
while the the total momentum $`P_u`$ is conserved
$$P_u๐\varphi p_u\dot{P}_u=๐\varphi \dot{p}_u=๐\varphi _\varphi [(\mathrm{})]=0$$
(80)
To simplify further the analysis of the time dependent evolution of supertube we will restrict ourselves to the case of homogenous fields on the worldvolume of D2-brane:
$$_\varphi R=_\varphi X^m=_\varphi Y^u=_\varphi \pi =0$$
(81)
As it is now clear from (79) and (77) $`\pi ^z\mathrm{\Pi }`$ and $`p_u`$ are conserved. On the other hand the equation of motion for $`A_z,A_\varphi `$ are equal to
$`\dot{A}_z=E={\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}(g_{zz}+g^{\varphi \varphi }B^2)\mathrm{\Pi }^2,`$
$`\dot{A}_\varphi =0,`$
where the second equation above implies that $`B=_zA_\varphi `$ is conserved as well. Then we get simple form of the Hamiltonian density $``$
$$=\sqrt{g_{00}}\sqrt{p_Rg^{RR}p_R+p_mg^{mn}p_n+p_ug^{uv}p_v+\mathrm{\Pi }g_{zz}\mathrm{\Pi }+b_\varphi g^{\varphi \varphi }b_\varphi +e^{2\mathrm{\Phi }}\tau _2^2\left(g_{zz}g_{\varphi \varphi }+B^2\right)},$$
(83)
where
$$b_\varphi =F_{\varphi z}\pi ^z=B\mathrm{\Pi }$$
(84)
For supertube in D$`p`$-brane background we have
$`g_{00}`$ $`=`$ $`H_p^{1/2},g_{zz}=H_p^{1/2},g_{\varphi \varphi }=H_p^{1/2}R^2,`$ (85)
$`g_{uv}`$ $`=`$ $`H_p^{1/2}\delta _{uv},g_{mn}=H_p^{1/2}\delta _{mn},g_{RR}=H_p^{1/2},`$ (87)
$`e^{2\mathrm{\Phi }}`$ $`=`$ $`H_p^{\frac{p3}{2}}`$ (89)
where the harmonic function for $`N`$ D$`p`$-branes is given by <sup>4</sup><sup>4</sup>4We work in units $`l_s=1`$.
$$H_p=1+\frac{N}{(R^2+X^mX_m)^{(7p)/2}}$$
(90)
and hence the Hamiltonian density takes the form
$$=\sqrt{\frac{\mathrm{\Pi }^2}{H_p}\left(1+\frac{B^2}{R^2}\right)+\frac{p_R^2}{H_p}+\frac{p_mp^m}{H_p}+p_up^u+\tau _2^2H_p^{\frac{p4}{2}}\left[R^2+B^2\right]}$$
(91)
that in the static case ($`p_u=p_R=p_m=0`$) reduces to the Hamiltonian studied in . For letter purposes it is also useful to define conserved energy $``$ through the relation
$$==2\pi =\frac{}{2\pi }.$$
(92)
As it is clear from the form of the harmonic function $`H_p`$ the background has $`SO(7p)`$ symmetry in the transverse subspace labelled with coordinates $`X^m`$. Then it is natural to simplify the analysis by restricting the dynamics of the supertube to the $`(x^6,x^7)`$ plane and introduce second polar coordinates as
$$x^6=\rho \mathrm{cos}\psi ,x^7=\rho \mathrm{sin}\psi $$
(93)
with corresponding metric components
$$g_{\rho \rho }=H_p^{1/2},g_{\psi \psi }=H_p^{1/2}\rho ^2.$$
(94)
Then the Hamiltonian density $``$ can be written as
$$=\sqrt{\frac{\mathrm{\Pi }^2}{H_p}\left(1+\frac{B^2}{R^2}\right)+\frac{p_R^2}{H_p}+\frac{p_\rho ^2}{H_p}+\frac{p_\psi ^2}{H_p\rho ^2}+p_up^u+\tau _2^2H_p^{\frac{p4}{2}}\left[R^2+B^2\right]}.$$
(95)
Since the Hamiltonian does not explicitly depend on $`\psi `$ we get immediately that $`p_\psi `$ is conserved:
$$\dot{p}_\psi =\frac{\delta H}{\delta \psi }=0.$$
(96)
Then the equations of motion for $`\rho ,R`$ take the form
$`\dot{R}`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta p_R}}={\displaystyle \frac{2\pi p_R}{H_p}},`$ (97)
$`\dot{\rho }`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta p_\rho }}={\displaystyle \frac{2\pi p_\rho }{H_p}}`$
On the other hand the equation of motion for $`p__R,p__\rho `$ are much more complicated thanks to the nontrivial dependence of $`H`$ on $`\rho ,R`$. In fact it is very complicated to solve these equations in the full generality. Let us rather consider the special case when we will study the dynamics of $`R`$ only. In order to do this we should find the stable values for $`\rho `$. Let us then consider the equation of motion for $`p_\rho `$
$`\dot{p}_\rho ={\displaystyle \frac{1}{2\sqrt{๐ฆ}}}{\displaystyle \frac{\delta ๐ฆ}{\delta H_p}}{\displaystyle \frac{\delta H_p}{\delta \rho }}+{\displaystyle \frac{p_\psi ^2}{\rho ^3H_p\sqrt{๐ฆ}}}={\displaystyle \frac{1}{\sqrt{๐ฆ}}}\left[{\displaystyle \frac{N(p7)\rho }{(R^2+\rho ^2)^{(8p)/2}}}{\displaystyle \frac{\delta ๐ฆ}{\delta H_p}}{\displaystyle \frac{p_\psi ^2}{\rho ^3H_p}}\right]`$
We see that the momentum $`p_\rho `$ is equal to zero for $`\rho =\psi =0`$. The question is whether there exist closed orbits with $`p_\psi 0`$ for which $`p_\rho =0`$. Looking at the equation above it is clear that we find $`\rho `$ as function of $`N`$, conserved momenta $`p_u,p_\psi `$ and, most importanly, as a function of $`R`$. Then it follows that the resulting Hamiltonian is complicated function of $`R`$. Even if it would be certainly interesting to study the properties of supertube with nonzero $`p_\psi `$ we will restrict ourselved in this paper to the case of $`p_\psi =\rho =0`$.
Analogously, we can consider the situation when $`p__R=0`$ that corresponds some particular value of $`R_{stat}=R(N,p_u,p_\psi ,\rho )`$ and study the time dependence of $`\rho `$. As in the previous case we leave the study of this problem for future.
Now let us consider the case when $`p_\rho =p_\psi =\rho =0`$. In this case the study of the supertube dynamics reduces to the study of the time dependence of $`R`$. Then it is natural to write the Hamiltonian density as as $`=\sqrt{\frac{p_R^2}{H_p}+V}`$ where $`\sqrt{V}=(p__R=0)`$. Then we can write the equation of motion for $`R`$ in the following form
$$\dot{R}=\frac{1}{\sqrt{H_p}}\sqrt{^2(2\pi )^2V},\frac{2}{2}\dot{R}^2+V_{\mathrm{eff}}=0,$$
(101)
where
$$V_{\mathrm{eff}}=\frac{1}{H_p}\left(\frac{(2\pi )^2V}{^2}1\right).$$
(102)
Then in order to obtain qualitative character of the dynamics of he supertube we use the observation that the equation given above corresponds to the conservation of energy for massive particle with mass ($`m=2`$) with the effective potential $`V_{\mathrm{eff}}`$ with total zero energy. We will be interested in two cases corresponding to $`p=2,4`$.
### 3.1 D2-brane background
In this case the effective potential takes the form
$$V=(R^2+B^2)\left(\frac{R^3}{R^5+N}\right)(\mathrm{\Pi }^2+\tau _2^2R^2)+p_up^u$$
(103)
and hence $`V_{\mathrm{eff}}`$ is equal to
$$V_{\mathrm{eff}}=\frac{R^5}{R^5+N}\left(\frac{4\pi ^2}{^2}(R^2+B^2)\frac{R^3}{R^5+N}(\mathrm{\Pi }^2+\tau _2^2R^2)+\frac{4\pi ^2p_up^u}{^2}1\right).$$
(104)
It is clear that the particle with zero energy can move in the interval between the points where $`V_{\mathrm{eff}}=0`$. First of all, the asymptotic behavior of $`V_{\mathrm{eff}}`$ is as follows
$`V_{\mathrm{eff}}{\displaystyle \frac{R^5}{N}}({\displaystyle \frac{4\pi ^2B^2\mathrm{\Pi }^2\tau _2R^3}{^2}}+{\displaystyle \frac{4\pi ^2p_up^u}{^2}}1),R0,`$
$`V_{\mathrm{eff}}{\displaystyle \frac{4\pi ^2}{^2}}R^2,R\mathrm{}.`$
Few comments are in order. From the above limits we can see that the potential approaches $`0`$ for $`R0`$ from below. Then since the potential blows up for $`R\mathrm{}`$ there should exists the point $`R_T`$ where the potential vanishes: $`V_{\mathrm{eff}}(R_T)=0`$. In order to study the dynamics around this point we introduce the variable $`r`$ through the substitution $`R=r+R_T`$ and insert it to the expression for conservation of energy. Using the fact that near the turning point we have
$`V_{\mathrm{eff}}(R)=V_{\mathrm{eff}}(R_T)+V_{\mathrm{eff}}^{}(R_T)r=V_{\mathrm{eff}}^{}(R_T)r`$ where we have used $`V_{\mathrm{eff}}(R_T)=0`$. Then we get the equation of motion for $`r`$
$$\dot{r}^2=V_{\mathrm{eff}}^{}(R_T)r\dot{r}=\pm \sqrt{V_{\mathrm{eff}}^{}(R_T)r}.$$
(106)
Since $`V_{\mathrm{eff}}^{}(R_T)>0`$ it follows from the equation above that $`r`$ should be negative. Integrating the above equation we get
$$r=\frac{1}{4}\left(r_0\sqrt{V_{\mathrm{eff}}^{}(R_T)}t\right)^2.$$
(107)
if we demand that for $`t=0`$ the particle reaches its turning point we have $`r_0=0`$ and hence we obtain
$$r=\frac{V_{\mathrm{eff}}}{4}t^2$$
(108)
so that for small negative $`t`$ (in order to trust $`r1`$) the particle approaches its turning point and then it moves back.
As follows from the properties of the effective potential $`V_{\mathrm{eff}}`$ <sup>5</sup><sup>5</sup>5Namely, since the potential reaches zero for $`R0`$ from below and since the potential blows up for $`R\mathrm{}`$ there should certainly exists the point where $`V_{\mathrm{eff}}^{}(R_m)=0`$. there should also exist the point where the potential reaches its local minimum corresponding to $`V_{eff}^{}(R_m)=0`$. Again, if we introduce $`r`$ as $`R=r+R_m`$ we get
$$V_{\mathrm{eff}}(R)=V_{\mathrm{eff}}(R_m)+\frac{1}{2}V_{\mathrm{eff}}^{\prime \prime }(R_m)r^2$$
(109)
where $`V_{\mathrm{eff}}(R_m)A<0`$ and $`\frac{1}{2}V_{\mathrm{eff}}^{\prime \prime }(R_m)B>0`$ and hence the conservation of energy implies the following differential equation
$`{\displaystyle \frac{dr}{\sqrt{1+\frac{B}{A}r^2}}}=\pm \sqrt{A}.`$ (110)
The above equation has a solution for $`r`$ that is
$`r=\sqrt{{\displaystyle \frac{|A|}{B}}}\mathrm{sin}\sqrt{B}t=\sqrt{{\displaystyle \frac{2|V_{\mathrm{eff}}(R_m)|}{V_{\mathrm{eff}}^{\prime \prime }(R_m)}}}\mathrm{sin}\sqrt{{\displaystyle \frac{V_{\mathrm{eff}}^{\prime \prime }(R_m)}{2}}}t.`$ (111)
In other words if we have a supertube inserted close to its local minimum position we observe that the supertube will fluctuate around this point with harmonic oscillations.
Finally, we will consider the case $`R0`$. Since now the effective potential is given by
$`V_{\mathrm{eff}}={\displaystyle \frac{R^5}{N}}\left({\displaystyle \frac{4\pi ^2p_up^u}{^2}}1\right),`$ (112)
the time evolution equation for $`R`$ is
$$\dot{R}^2=\frac{R^5}{N}\left(\frac{4\pi ^2p_up^u}{^2}1\right)\frac{1}{R^3}=\left(C\frac{3t}{2}\sqrt{\frac{1}{N}\left(1\frac{4\pi ^2p_up^u}{^2}\right)}\right)^2,C^2=\frac{1}{R_0^3},$$
(113)
whrere $`R_0`$ is the position of supertube in time $`t=0`$. However since we demand that $`R0`$ it is clear that the above solution is valid in case of large positive or negative $`t`$ and we obtain following assymptotic behaviour of $`R`$
$$R\frac{1}{t^{2/3}}\frac{1}{\left[\frac{9}{4N}\left(1\frac{4\pi ^2p_up^u}{^2}\right)\right]^{1/3}}.$$
(114)
This result shows that supertube approaches the worldvolume of $`N`$ D2-branes for $`t\mathrm{}`$.
### 3.2 D4-brane background
The situation of supertube in D4-brane background is similar to the case of D2-brane background. Namely, the potential $`V`$ takes the form
$$V=(R^2+B^2)\left[\frac{\mathrm{\Pi }^2R}{R^3+N}+\tau _2^2\right]+p_up^u$$
(115)
and hence the effective potential takes the form
$$V_{\mathrm{eff}}=\frac{R^3}{(N+R^3)}\left(\frac{4\pi ^2}{^2}(R^2+B^2)\left[\frac{\mathrm{\Pi }^2R}{R^3+N}+\tau _2^2\right]+\frac{4\pi ^2p_up^u}{^2}1\right).$$
(116)
from which we again obtain the asymptotic behavior
$`V_{\mathrm{eff}}{\displaystyle \frac{R^3}{N}}\left({\displaystyle \frac{4\pi ^2(B^2\tau _2^2+p_u^2)}{^2}}1\right),R0`$
$`V_{\mathrm{eff}}{\displaystyle \frac{4\pi ^2\tau _2^2}{^2}}R,R\mathrm{},`$
Since $`\frac{^2}{4\pi ^2}>(B^2\tau _2^2+p_u^2)`$, we see that the effective potential approaches zero in the limit $`R0`$ from below. In the same way we also see that the potential blows up for $`R\mathrm{}`$. It then follows that the behavior of the supertube in this background is the same as in the case of D2-brane background studied in the previous section. Therefore we skip the details of the discussions in this case.
### 3.3 NS5-brane background
As the next example we will consider supertube in the background of $`N`$ NS5-branes. The metric, the dilaton and NS-NS field are
$`ds^2`$ $`=`$ $`dx_\mu dx^\mu +H_{NS}(R)dx^mdx^m,`$ (118)
$`e^{2(\mathrm{\Phi }\mathrm{\Phi }_0)}`$ $`=`$ $`H_{NS}(R),`$ (120)
$`H_{mnp}`$ $`=`$ $`ฯต_{mnp}^q_q\mathrm{\Phi },`$
where $`H_{NS}=1+\frac{N}{R^2}`$ is the harmonic function in the transverse directions of $`N`$ NS5-branes. If we again restrict to the homogenous modes the Hamiltonian density takes the form
$`=\sqrt{p_R{\displaystyle \frac{1}{H_{NS}}}p_R+p_ug^{uv}p_v+\mathrm{\Pi }^2+b_\varphi {\displaystyle \frac{1}{R^2H_{NS}}}b_\varphi +{\displaystyle \frac{\tau _2^2}{H_{NS}}}\left(R^2H_{NS}+B^2\right)}`$
$`={\displaystyle \frac{1}{\sqrt{N+R^2}}}\sqrt{p_R^2R^2+p_up^u+(N+R^2+B^2)(\mathrm{\Pi }^2+\tau _2^2R^2)},`$
where we have used
$$b_\varphi =F_{\varphi z}\pi ^z=B\mathrm{\Pi }$$
(124)
and the fact that $`\pi ^z\mathrm{\Pi }`$ and $`p_u`$ are conserved. We must also stress that we consider the situation when modes $`X^m`$ that were defined in previous section, are not excited and equal to zero.
Now the equation of motions take the form
$`\dot{A}_z`$ $`=`$ $`E={\displaystyle \frac{\sqrt{g_{00}}}{\sqrt{๐ฆ}}}(g_{zz}+g^{\varphi \varphi }B^2)\mathrm{\Pi }^2,`$ (125)
$`\dot{A}_\varphi `$ $`=`$ $`0,`$
where the second equation above implies that $`B=_zA_\varphi `$ is conserved as well.
In order to study the dynamics of the radial mode we use as in the previous sections the fact that the energy is conserved. Firstly, if we write the Hamiltonian density as $`=\sqrt{p_Rg^{RR}p_R+V}`$ where
$$V=\frac{1}{N+R^2}\left(p_up^u+(N+R^2+B^2)(\mathrm{\Pi }^2+\tau _2^2R^2)\right)$$
(128)
Now we can express $`p_R`$ from $`=2\pi `$ as
$$p_R=\frac{\sqrt{g_{RR}}}{2\pi }\sqrt{^24\pi ^2V}$$
(129)
and hence the equation of motion for $`\dot{R}`$ is
$`\dot{R}={\displaystyle \frac{p_Rg^{RR}}{}}={\displaystyle \frac{1}{\sqrt{g_{RR}}}}\sqrt{^24\pi ^2V}`$ (130)
(131)
$`\dot{R}^2+V_{\mathrm{eff}}=0,`$ (132)
where
$`V_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{1}{g_{RR}}}\left({\displaystyle \frac{4\pi ^2V}{^2}}1\right)`$
$`=`$ $`{\displaystyle \frac{R^2}{N+R^2}}\left({\displaystyle \frac{4\pi ^2}{(N+R^2)^2}}\left(p_up^u+(N+R^2+B^2)(\mathrm{\Pi }^2+\tau _2^2R^2)\right)1\right).`$
For small and for large $`R`$ this potential reduces to
$`V_{\mathrm{eff}}{\displaystyle \frac{R^2}{N}}\left({\displaystyle \frac{4\pi ^2}{^2N}}(p_up^u+(N+B^2)\mathrm{\Pi }^2)1\right),R0,`$
$`V_{\mathrm{eff}}{\displaystyle \frac{4\pi ^2R^2}{^2}},R\mathrm{}.`$
Again we see that the potential $`V_{\mathrm{eff}}`$ approaches zero for $`R0`$ from below and blows up for $`R\mathrm{}`$. Consequently we obtain the same picture of the supertube dynamics as in the previous sections.
### 3.4 Macroscopic fundamental string background
In this section, we study the time dependent dynamics of the tube in the Macroscopic fundamental string background, where the metric, dilaton and the NS-NS charge of the F-background is given by
$`ds^2={\displaystyle \frac{1}{H_f(r)}}\left(dt^2+dz^2\right)+\delta _{mn}dx^mdx^n,B_{01}=H_f^11,e^{2\mathrm{\Phi }}=H_f^1,`$ (137)
where $`H_f=1+\frac{N}{r^6}`$ is the harmonic function in the transverse eight-space of the $`N`$ F-strings. $`r`$ denotes the spatial coordinate transverse to the macroscopic string.
Let us again consider the Hamiltonian density for D2-brane
$``$ $`=`$ $`\sqrt{g_{00}}\sqrt{๐ฆ}+\mathrm{\Pi }^aB_{a0},`$ (138)
$`๐ฆ`$ $`=`$ $`(p_r+\pi ^aB_{ar})g^{rs}(p_s+\pi ^aB_{as})+\pi ^a\gamma _{ab}\pi ^b+b_ag^{ab}b_b+e^{2\mathrm{\Phi }}\tau _2^2det๐_{ab},`$ (140)
$`b_a`$ $`=`$ $`\mathrm{\Pi }_s_aX^s+F_{ab}\pi ^b.`$
The Hamiltonian density (138) is the starting point for our calculation where we consider D2-brane supertube stretched in $`z`$ direction and that wind $`\varphi `$ direction in the plane $`(x^8,x^9)`$ in the space transverse to the fundamental string, where the modes that parameterize the embedding of the supertube, are time dependent only:
$$R=R(t),X^m=X^m(t)m=3,\mathrm{},7.$$
(143)
Now thanks to the manifest rotation symmetry $`SO(6)`$ of the subscpace $`(x^3,\mathrm{},x^7)`$ we can restrict ourselves to the motion of supertube in $`(x^3,x^4)`$ plane where we introduce polar coordinates as
$$X^3=\rho \mathrm{cos}\psi ,X^4=\rho \mathrm{sin}\psi .$$
(144)
Note also that the fact that we consider homogenous modes implies that $`A_z`$ and $`\mathrm{\Pi }`$ are time independent.
Now the spatial matrix $`๐_{ab}`$ takes the form
$$๐=\left(\begin{array}{cc}g_{zz}& B_{z\varphi }\\ B_{\varphi z}& g_{\varphi \varphi }\end{array}\right)=\left(\begin{array}{cc}H_f^1& B\\ B& R^2\end{array}\right)$$
(145)
hence its determinant is equal to
$$det๐_{ab}=H_f^1R^2+B^2,H_f=1+\frac{N}{(R^2+\rho ^2)^3}.$$
(146)
We also consider the electric flux in the $`z`$ direction only $`\pi ^z\mathrm{\Pi }`$ so that the term $`\pi ^aB_{aM}=\pi ^zB_{zM}=0`$ using the fact that all fields do not depend on $`z`$. In the same way $`b_\varphi =F_{\varphi z}\pi ^z=B\mathrm{\Pi }`$. Now the term $`๐ฆ`$ takes the form
$`๐ฆ=p_R^2+p_\rho ^2+{\displaystyle \frac{p_\psi ^2}{\rho ^2}}+{\displaystyle \frac{\mathrm{\Pi }^2}{H_f}}+{\displaystyle \frac{B^2\mathrm{\Pi }^2}{R^2}}+\tau _2^2(R^2+H_fB^2).`$
It is now easy to see that we can restrict ourselves to the motion when $`p_\rho =p_\psi =\rho =0`$. Then we get
$`๐ฆ={\displaystyle \frac{(R^8+B^2N+B^2R^6)}{R^6}}\left({\displaystyle \frac{R^6p_R^2}{R^8+B^2N+B^2R^6}}+{\displaystyle \frac{\mathrm{\Pi }^2R^4+B+R^6}{N+R^6}}\right).`$
Then finally the Hamiltonian density takes the form
$``$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{R^6+N}}}\sqrt{(R^8+B^2N+B^2R^6)\left({\displaystyle \frac{R^6p_R^2}{R^8+B^2N+B^2R^6}}+{\displaystyle \frac{\mathrm{\Pi }^2R^4+B+R^6}{N+R^6}}\right)}`$
$`+{\displaystyle \frac{\mathrm{\Pi }N}{N+R^6}}.`$
In order to simplify notation we will write the Hamiltonian density for supertube in Fstring background in the form
$$=\sqrt{F(R)p_R^2+G(R)}+\frac{\mathrm{\Pi }N}{R^6+N},$$
(152)
where
$$F=\frac{R^6}{R^6+N},G(R)=\frac{(R^8+R^6B^2+B^2N)(R^6+\mathrm{\Pi }^2R^4+N)}{(R^6+N)^2}.$$
(153)
Then the differential equation for $`R`$ is
$$\dot{R}=\frac{Fp_R}{\sqrt{(\mathrm{})}}=\frac{Fp_R}{\left(\frac{}{2\pi }\frac{\mathrm{\Pi }N}{R^6+N}\right)}$$
(154)
where we have expressed the square root using the conserved energy $``$. If we also express $`p_R`$ using $``$ as
$$p_R^2=\frac{1}{F}\left[\left(\frac{}{2\pi }\frac{\mathrm{\Pi }NF}{R^6}\right)^2G\right]$$
(155)
we get
$`\dot{R}^2`$ $`=`$ $`{\displaystyle \frac{F^2p_R^2}{\left(\frac{}{2\pi }\frac{\mathrm{\Pi }NF}{R^6}\right)^2}}={\displaystyle \frac{F}{\left(\frac{}{2\pi }\frac{\mathrm{\Pi }NF}{R^6}\right)^2}}\left[\left({\displaystyle \frac{}{2\pi }}{\displaystyle \frac{\mathrm{\Pi }NF}{R^6}}\right)^2G\right]`$ (156)
$`V_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{F}{\left(\frac{}{2\pi }\frac{\mathrm{\Pi }NF}{R^6}\right)^2}}\left[G\left({\displaystyle \frac{}{2\pi }}{\displaystyle \frac{\mathrm{\Pi }NF}{R^6}}\right)^2\right].`$
Now we can perform the analysis exactly in the same way as in the previous section. Namely, the asymptotic behavior of $`V_{\mathrm{eff}}`$ is
$`V_{\mathrm{eff}}`$ $``$ $`{\displaystyle \frac{R^6}{N}}{\displaystyle \frac{1}{(\frac{}{2\pi }\mathrm{\Pi })^2}}\left[B^2\left({\displaystyle \frac{}{2\pi }}\mathrm{\Pi }\right)^2\right],R0,`$ (159)
$`V_{\mathrm{eff}}`$ $``$ $`{\displaystyle \frac{R^24\pi ^2}{^2}},R\mathrm{}`$
The potential approaches zero from below for $`R0`$ on condition
$$\frac{}{2\pi }>B+\mathrm{\Pi }$$
(162)
that is again obeyed. Consequently the analysis is the same as in the examples studied in the previous sections.
## 4 Summary and Conclusion
In this paper, we have studied the time dependent dynamics of the supertube in various backgrounds, by using the Dirac-Born-Infeld effective field theory description in Hamiltonian formalism. This provide yet another example of time dependent solutions of string theory in terms of some exotic bound of various D-branes with Born-Infeld electric and magnetic fields. The stability of these bound states in various curved backgrounds that are discussed in the literature recently, has been very suggestive in studying such dynamics and to view the trajectories of the supertube. We adopt the Hamiltonian formalism for studying such dynamics of D-brane in general curved backgrounds in the presence of worldvolume gauge field. By using the crucial gauge fixing, we have studied the dynamics in Dp-brane, NS5-brane and fundamental string backgrounds and have analyzed the effective potential and the trajectory. It would be very interesting to study some other exotic bound states of branes (e.g. it has been argued in that a string network can be blown up into a D3-brane, the so called supertube) in general curved background. There the analysis in terms on Hamiltonian formalism seems much more complicated, but nevertheless a doable problem. It would also be very interesting to study the dynamics in Anti-de Sitter (AdS) spaces, and to study the effect in the dual conformal field theory by using the gauge-gravity duality. In the case of D-branes in AdS backgrounds, the tachyon like radion would be spatial dependent and hence studying the dynamics ought to shed light in studying the various properties of D-brane in curved backgrounds. We hope to return one or more of these issues in near future.
Acknowledgements: We would like to thank M. Bianchi, U. Lindstrom, Rikard von Unge and A. Sagnotti for various useful discussions. J.K. would like to thank the Dipartimento di Fisica, โTor Vergataโ for the kind hospitality where a part of the work was done. The work of J.K. was supported by the Czech Ministry of Education under Contract No. MSM 0021622409. The work of K.L.P. was supported in part by INFN, by the MIUR-COFIN contract 2003-023852, by the EU contracts MRTN-CT-2004-503369 and MRTN-CT-2004-512194, by the INTAS contract 03-51-6346, and by the NATO grant PST.CLG.978785. |
warning/0506/cs0506011.html | ar5iv | text | # On the dimensions of certain LDPC codes based on ๐-regular bipartite graphs
## I Introduction
Let $`V`$ be a $`4`$-dimensional vector space over the field $`๐
_q`$ of $`q`$ elements. We assume that $`V`$ has a nonsingular alternating bilinear form $`(v,v^{})`$ and denote by $`\mathrm{Sp}(V)`$ the group of linear automorphisms of $`V`$ which preserve this form. We choose a symplectic basis $`e_0`$, $`e_1`$, $`e_2`$, $`e_3`$ of $`V`$, with $`(e_i,e_{3i})=1`$, for $`i=0`$, $`1`$.
Let $`P=๐(V)`$ be the set of points of the projective space of $`V`$. A subspace of $`V`$ is said to be *totally isotropic* if $`(v,v^{})=0`$ whenever $`v`$ and $`v^{}`$ are both in the subspace. Let $`L`$ denote the set of totally isotropic $`2`$-dimensional subspaces of $`V`$, considered as lines in $`P`$. The pair $`(P,L)`$, together with the natural relation of incidence between points and lines, is called the symplectic generalized quadrangle. Except for in the appendix, the term โlineโ will always mean an element of $`L`$. It is easy to verify that $`(P,L)`$ satisfies the following quadrangle property: Given any line and any point not on the line, there is a unique line which passes through the given point and meets the given line.
Now fix a point $`p_0P`$ and a line $`\mathrm{}_0L`$ through $`p_0`$. We can assume that we chose our basis so that $`p_0=e_0`$ and $`\mathrm{}_0=e_0,e_1`$. For $`pP`$, denote by $`p^{}`$ the set of points on lines through $`p`$; $`p^{}p^{}`$ if and only if the subspace of $`V`$ spanned by $`p`$ and $`p^{}`$ is isotropic. Consider the set $`P_1=Pp_0^{}`$ of points not collinear with $`p_0`$, and the set $`L_1`$ of lines which do not meet $`\mathrm{}_0`$. Then we can also consider the incidence systems $`(P_1,L_1)`$, $`(P,L_1)`$ and $`(P_1,L)`$. Let $`M(P,L)`$ and $`M(P_1,L_1)`$ be the binary incidence matrices of the respective incidence systems, with rows indexed by points and columns by lines. The rows and columns of $`M(P,L)`$ have weight $`q+1`$ and, as a consequence of the quadrangle property, those of $`M(P_1,L_1)`$ have weight $`q`$.
If $`q`$ is odd we know by Theorem 9.4 of that the $`2`$-rank of $`M(P,L)`$ is $`(q^3+2q^2+q+2)/2`$. Here we prove the following theorem.
###### Theorem I.1
Assume $`q`$ is a power of an odd prime. The $`2`$-rank of $`M(P_1,L_1)`$ equals $`(q^3+2q^23q+2)/2`$.
In , a family of codes designated $`\mathrm{LU}(3,q)`$ was defined in the following way. Let $`P^{}`$ and $`L^{}`$ be sets in bijection with $`๐
_{q}^{}{}_{}{}^{3}`$, where $`q`$ is any prime power. An element $`(a,b,c)P^{}`$ is incident with an element $`[x,y,z]L^{}`$ if and only if
$$y=ax+b\text{and}z=ay+c.$$
(1)
The binary incidence matrix with rows indexed by $`L^{}`$ and columns indexed by $`P^{}`$ is denoted by $`H(3,q)`$ and the two binary codes having $`H(3,q)`$ and its transpose as parity check matrices are called $`\mathrm{LU}(3,q)`$ codes. The name comes from , where the bipartite graph with parts $`P^{}`$ and $`L^{}`$ and adjacency defined by the equations (1) had been studied previously.
It is not difficult to show that the incidence systems $`(P_1,L_1)`$ and $`(P^{},L^{})`$ are equivalent. A detailed proof is given in the appendix. Thus, $`M(P_1,L_1)`$ is a parity check matrix of the $`\mathrm{LU}(3,q)`$ code given by the transpose of $`H(3,q)`$ and Theorem I.1 has the following immediate corollary.
###### Corollary I.2
If $`q`$ is a power of an odd prime, the dimension of $`\mathrm{LU}(3,q)`$ is $`(q^32q^2+3q2)/2`$.
The corollary was conjectured in . There it was established that this number is a lower bound when $`q`$ is an odd prime.
## II Relative dimensions and a lower bound for $`\mathrm{LU}(3,q)`$
In this section $`q`$ is an arbitrary prime power.
Let $`๐
_2[P]`$ be the vector space of all $`๐
_2`$-valued functions on $`P`$. We can think of such a function as a vector in which the positions are indexed by the points of $`P`$, and the entries are the values of the function at the points. For $`pP`$, the characteristic function $`\chi _p`$ is the vector with $`1`$ in the position with index $`p`$ and zero in the other positions. The set of all characteristic functions of points forms a basis of $`๐
_2[P]`$. Let $`\mathrm{}L`$. Its characteristic function $`\chi _{\mathrm{}}๐
_2[P]`$ is the function which takes the value $`1`$ at the $`q+1`$ points of $`\mathrm{}`$ and zero at all other points. The subspace of $`๐
_2[P]`$ spanned by all the $`\chi _{\mathrm{}}`$ is the $`๐
_2`$-code of $`(P,L)`$, denoted by $`C(P,L)`$. One can think of $`C(P,L)`$ as the column space of $`M(P,L)`$. For brevity, we will sometimes blur the distinction between lines and their characteristic functions and speak, for instance, of the subspace of $`๐
_2[P]`$ spanned by a set of lines. Let $`C(P,L_1)`$ be the subspace of $`๐
_2[P]`$ spanned by lines in $`L_1`$. Let $`C(P_1,L_1)`$ denote the code of $`(P_1,L_1)`$, viewed as a subspace of $`๐
_2[P_1]`$, and let $`C(P_1,L)`$ be the larger subspace of $`๐
_2[P_1]`$ spanned by the restrictions to $`P_1`$ of the characteristic functions of all lines of $`L`$.
Consider the natural projection map
$$\pi _{P_1}:๐
_2[P]๐
_2[P_1]$$
(2)
given by restriction of functions. Its kernel will be denoted by $`\mathrm{ker}\pi _{P_1}`$.
Let $`ZC(P,L_1)`$ be a set of characteristic functions of lines in $`L_1`$ which maps bijectively under $`\pi _{P_1}`$ to a basis of $`C(P_1,L_1)`$. Let $`X`$ be the set of characteristic functions of the $`q+1`$ lines of $`L`$ through $`p_0`$ and let $`X_0=X\{\chi _\mathrm{}_0\}`$. Finally, choose any $`q`$ lines of $`L`$ which meet $`\mathrm{}_0`$ in the $`q`$ distinct points other than $`p_0`$ and let $`Y`$ be the set of their characteristic functions. It is clear that the sets $`X`$, $`Y`$ and $`Z`$ are disjoint and that $`X`$ is contained in $`\mathrm{ker}\pi _{P_1}`$.
###### Lemma II.1
$`ZX_0Y`$ is linearly independent over $`๐
_2`$.
###### Proof:
Each element of $`Y`$ contains in its support a point of $`\mathrm{}_0`$ which is not in the support of any other element of $`ZX_0Y`$. So it is enough to show that $`X_0Z`$ is linearly independent. This is true because $`X_0`$ is a linearly independent subset of $`\mathrm{ker}\pi _{P_1}`$ and $`Z`$ maps bijectively under $`\pi _{P_1}`$ to a linearly independent set. โ
We note that $`|Z|=dim_{๐
_2}C(P_1,L_1)`$ and $`|X_0Y|=2q`$.
###### Corollary II.2
Let $`q`$ be an arbitrary prime power. Then
$$dim_{๐
_2}\mathrm{LU}(3,q)q^3dim_{๐
_2}C(P,L)+2q.$$
(3)
###### Proof:
From the definition of $`\mathrm{LU}(3,q)`$ and the equivalence of $`(P^{},L^{})`$ with $`(P_1,L_1)`$, we have
$$dim_{๐
_2}\mathrm{LU}(3,q)=q^3dim_{๐
_2}C(P_1,L_1).$$
(4)
The corollary now follows from Lemma II.1. โ
## III Proof of Theorem I.1
In this section we assume that $`q`$ is odd. In view of Corollary II.2 and the known $`2`$-rank of $`M(P,L)`$ the proof of Theorem I.1 will be completed if we can show that $`ZX_0Y`$ spans $`C(P,L)`$ as a vector space over $`๐
_2`$.
###### Lemma III.1
Let $`\mathrm{}L`$. Then the sum of the characteristic functions of all lines which meet $`\mathrm{}`$ (excluding $`\mathrm{}`$ itself) is the constant function 1.
###### Proof:
The function given by the sum takes the value $`q1`$ $`(\mathrm{mod}2)`$ at any point of $`\mathrm{}`$ and value $`1`$ at any point off $`\mathrm{}`$, by the quadrangle property. โ
###### Lemma III.2
Let $`\mathrm{}L`$ be a line, other than $`\mathrm{}_0`$, which meets $`\mathrm{}_0`$ at a point $`p`$. Let $`\mathrm{\Phi }_{\mathrm{}}`$ be the sum of all the characteristic functions of lines in $`L_1`$ which meet $`\mathrm{}`$. Then
$$\mathrm{\Phi }_{\mathrm{}}(p^{})=\{\begin{array}{cc}0,\text{if }p^{}=p;\hfill & \\ q,\text{if }p^{}\mathrm{}\{p\};\hfill & \\ 0,\text{if }p^{}p^{}\mathrm{};\hfill & \\ 1,\text{if }p^{}Pp^{}.\hfill & \end{array}$$
(5)
###### Proof:
This is an immediate consequence of the quadrangle property. โ
###### Corollary III.3
Let $`p\mathrm{}_0`$ and let $`\mathrm{}`$, $`\mathrm{}^{}`$ be two lines through $`p`$, neither equal to $`\mathrm{}_0`$. Then $`\chi _{\mathrm{}}\chi _{\mathrm{}^{}}C(P,L_1)`$.
###### Proof:
Since $`q=1`$ in $`๐
_2`$, one easily check using Lemma III.2 that
$$\chi _{\mathrm{}}\chi _{\mathrm{}^{}}=\mathrm{\Phi }_{\mathrm{}}\mathrm{\Phi }_{\mathrm{}^{}}C(P,L_1).$$
(6)
We now come to our main technical lemma.
###### Lemma III.4
$`\mathrm{ker}\pi _{P_1}C(P,L)`$ has dimension $`q+1`$, with basis the set $`X`$ of characteristic functions of the $`q+1`$ lines through $`p_0`$.
###### Proof:
Let $`G_{p_0}`$ be the stabilizer in $`\mathrm{Sp}(V)`$ of $`p_0`$.
From the definition,
$$\mathrm{ker}\pi _{P_1}=๐
_2[p_0^{}]=๐
_2[\{p_0\}]๐
_2[p_{0}^{}{}_{}{}^{}\{p_0\}]$$
(7)
as an $`๐
_2G_{p_0}`$-module. Clearly $`๐
_2[\{p_0\}]`$ is a one-dimensional trivial $`๐
_2G_{p_0}`$-module. To find the structure of $`๐
_2[p_{0}^{}{}_{}{}^{}\{p_0\}]`$, we consider the following subgroups of $`G_{p_0}`$, which we will describe as matrix groups with respect to our chosen basis.
Let
$$Q=\{\left(\begin{array}{cccc}1& a& b& c\\ 0& 1& 0& b\\ 0& 0& 1& a\\ 0& 0& 0& 1\end{array}\right)a,b,c๐
_q\}$$
(8)
and
$$C=\{\left(\begin{array}{cccc}1& 0& 0& c\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)c๐
_q\}.$$
(9)
The group $`Q`$ is a normal subgroup of $`G_{p_0}`$ and $`C`$ is the center of $`Q`$, with $`Q/C`$ elementary abelian of order $`q^2`$. It is easy to see by matrix computations that $`C`$ acts trivially on $`p_0^{}`$ and that $`Q`$ stabilizes each line $`\mathrm{}`$ through $`p_0`$, acting transitively on the $`q`$ points of $`\mathrm{}\{p_0\}`$. These $`q`$ points have homogeneous coordinates of the form $`[d:x:y:0]`$, where $`[x:y]`$ are homogeneous coordinates of a fixed point on a projective line, and $`d`$ varies over $`๐
_q`$. It is clear that the subgroup $`Q[x:y]`$ of index $`q`$ in $`Q`$ consisting of matrices (8) in which $`ax+by=0`$ is the kernel of the action on $`\mathrm{}\{p_0\}`$ and so $`๐
_2[\mathrm{}\{p_0\}]`$ affords the regular representation of $`Q/Q[x:y]`$.
As $`[x:y]`$ varies over the projective line, we deduce that, $`๐
_2[p_0^{}\{p_0\}]`$ contains the trivial module of $`Q/C`$ with multiplicity $`q+1`$. Thus since $`Q`$ has odd order, we have a $`๐
_2G_{p_0}`$-module decomposition
$$๐
_2[p_0^{}\{p_0\}]=TW,$$
(10)
where $`T`$ is the $`(q+1)`$-dimensional space of $`Q`$-fixed points and $`W`$ has dimension $`q^21`$ and no $`Q`$-fixed points. Let $`E`$ be a splitting field for $`Q`$ over $`๐
_2`$, and consider the action of $`G_{p_0}`$ on the characters of $`Q/C`$ which occur in $`E_{๐
_2}W`$. Each of the $`q^21`$ nontrivial characters occurs once. The group of matrices of the form $`\mathrm{diag}(\lambda ,\mu ,\mu ^1,\lambda ^1)`$, with $`\lambda `$, $`\mu ๐
_q\{0\}`$, lies in $`G_{p_0}`$ and acts transitively on the $`q1`$ nontrivial elements, hence also on the $`q1`$ nontrivial characters, of each $`Q/Q[x:y]`$. Then, since $`G_{p_0}`$ is transitive on the $`q+1`$ lines through $`p_0`$, it follows that the $`q^21`$ nontrivial characters of $`Q/C`$ form a single $`G_{p_0}`$-orbit. By Cliffordโs Theorem ((11.1) in ) it follows that $`E_{๐
_2}W`$ is a simple $`EG_{p_0}`$-module. Hence $`W`$ is a simple $`๐
_2G_{p_0}`$-module.
We are now ready to consider the intersection
$$\mathrm{ker}\pi _{P_1}C(P,L)=๐
_2[p_0^{}]C(P,L),$$
(11)
which is an $`๐
_2G_{p_0}`$-submodule of $`๐
_2[p_0^{}]`$. Clearly, $`X`$ is a linearly independent subset of this intersection. Moreover, each element of $`X`$ is a fixed point of $`Q`$. We must prove that the intersection is no bigger than the span of $`X`$. If it were, then by what we know of the $`๐
_2G_{p_0}`$-submodules of $`๐
_2[p_0^{}]`$, we see that either $`๐
_2[p_0^{}]C(P,L)`$ must contain all the $`Q`$-fixed points of $`๐
_2[p_0^{}]`$ or else it must contain $`W`$. The first possibilty is ruled out because it implies that $`C(P,L)`$ contains the characteristic function of the point $`p_0`$, which is absurd since the number of points on a line is even. In the second case, we would have that $`๐
_2[p_0^{}]C(P,L)`$ is of codimension one in $`๐
_2[p_0^{}]`$. Then, for any point $`pp_0^{}`$, since neither $`\chi _p`$ nor $`\chi _{p_0}`$ is in $`C(P,L)`$, we would have $`\chi _p\chi _{p_0}C(P,L)`$. Then, by transitivity of $`\mathrm{Sp}(V)`$ on $`P`$ and the connectedness of the adjacency graph of $`P`$, we would have that $`\chi _p\chi _{p_0}C(P,L)`$ for all points $`pP`$, leading to the conclusion that $`C(P,L)`$ has codimension one in $`๐
_2[P]`$, contrary to known fact. Thus, the intersection is as claimed. โ
###### Lemma III.5
$`\mathrm{ker}\pi _{P_1}C(P,L_1)`$ has dimension $`q1`$, and basis the set of functions $`\chi _{\mathrm{}}\chi _{\mathrm{}^{}}`$, where $`\mathrm{}\mathrm{}_0`$ is an arbitrary but fixed line through $`p_0`$ and $`\mathrm{}^{}`$ varies over the $`q1`$ lines through $`p_0`$ different from $`\mathrm{}_0`$ and $`\mathrm{}`$.
###### Proof:
By Corollary III.3 applied to $`p_0`$, we see that if $`\mathrm{}`$ and $`\mathrm{}^{}`$ are any two of the $`q`$ lines through $`p_0`$ other than $`\mathrm{}_0`$, the function $`\chi _{\mathrm{}}\chi _{\mathrm{}^{}}`$ lies in $`C(P,L_1)`$. It is obviously in $`\mathrm{ker}\pi _{P_1}`$. Clearly, we can find $`q1`$ linearly independent functions of this kind as described in the statement. Thus $`\mathrm{ker}\pi _{P_1}C(P,L_1)`$ has dimension $`q1`$. On the other hand $`C(P,L_1)`$ is in the kernel of the restriction map to $`\mathrm{}_0`$, while the image of the restriction of $`\mathrm{ker}\pi _{P_1}C(P,L)`$ in $`\mathrm{}_0`$ has dimension 2, spanned by the images of $`\chi _\mathrm{}_0`$ and $`\chi _{p_0}`$. Thus $`\mathrm{ker}\pi _{P_1}C(P,L_1)`$ has codimension at least 2 in $`\mathrm{ker}\pi _{P_1}C(P,L)`$, which has dimension $`q+1`$, by Lemma III.4. โ
Our final lemma completes the proof of Theorem I.1.
###### Lemma III.6
$`ZX_0Y`$ spans $`C(P,L)`$ as a vector space over $`๐
_2`$.
###### Proof:
By Lemma III.5, the span of $`X_0Z`$ is equal to the subspace spanned by $`X_0`$ and $`L_1`$, since $`\mathrm{ker}\pi _{P_1}C(P,L_1)`$ is contained in the span of $`X_0`$. We must show that the subspace spanned by $`X_0Y`$ and $`L_1`$ contains the characteristic functions of all lines intersecting $`\mathrm{}_0`$, including $`\mathrm{}_0`$. First, consider a line $`\mathrm{}\mathrm{}_0`$ meeting $`\mathrm{}_0`$. We can assume that $`\mathrm{}`$ meets $`\mathrm{}_0`$ at a point other than $`p_0`$, since otherwise $`\mathrm{}X_0`$. Therefore $`\mathrm{}`$ meets $`\mathrm{}_0`$ in the same point $`p`$ as some element $`\mathrm{}^{}Y`$. Then Corollary III.3 shows that $`\chi _{\mathrm{}}`$ lies in the subspace spanned by $`Y`$ and $`L_1`$. The only line still missing is $`\mathrm{}_0`$, so our last task is to show that $`\chi _\mathrm{}_0`$ lies in the span of the characteristic functions of all other lines. First, by Lemma III.1 applied to $`\mathrm{}_0`$, we see that the constant function 1 is in the span. Finally, we see from Lemma III.2 that
$$\underset{\mathrm{}X_0}{}\mathrm{\Phi }_{\mathrm{}}=1\chi _\mathrm{}_0,$$
(12)
so we are done. โ
###### Remark III.7
One can also consider the binary code $`\mathrm{LU}(3,q)`$ when $`q=2^t`$, $`t1`$. The exact dimension is not known yet, but Corollary II.2 provides a lower bound, since by we have
$$dim_{๐
_2}C(P,L)=1+\left(\frac{1+\sqrt{17}}{2}\right)^{2t}+\left(\frac{1\sqrt{17}}{2}\right)^{2t}.$$
(13)
This formula is quite different from the one for odd $`q`$. Nevertheless, it may well be that the inequality (3) is an equality for even $`q`$, just as it is for odd $`q`$, despite the difference in the $`dim_{๐
_2}C(P,L)`$ term. Computer calculations of J.-L. Kim verify this up to $`q=16`$.
In this appendix $`q`$ is an arbitrary prime power. Here we explain why our incidence system $`(P_1,L_1)`$ is equivalent to the incidence system $`(P^{},L^{})`$ defined by the equations (1). The explanation is given by the classical Klein correspondence.
We first look at $`(P_1,L_1)`$ in coordinates. Let $`x_0`$, $`x_1`$, $`x_2`$, $`x_3`$ be homogeneous coordinates of $`P`$ corresponding to our symplectic basis. Recalling that $`p_0=e_0`$, we see that $`P_1`$ is the set of points such that $`x_30`$. If we represent such a point as $`(a:b:c:1)`$ we have a bijection of $`P_1`$ with $`๐
_{q}^{}{}_{}{}^{3}`$.
Our choice of basis of $`V`$ yields the basis $`e_ie_j`$, for $`0i<j3`$, of the exterior square $`^2(V)`$. Denote the corresponding homogeneous coordinates of the projective space $`๐(^2(V))`$ by $`p_{01}`$, $`p_{02}`$, $`p_{03}`$, $`p_{12}`$, $`p_{13}`$ and $`p_{23}`$. A $`2`$-dimensional subspace of $`V`$ spanned by vectors $`_{i=0}^3a_ie_i`$ and $`_{i=0}^3b_ie_i`$ defines, by taking its exterior square, a point of $`๐(^2(V))`$ with coordinates $`p_{ij}=a_ib_ja_jb_i`$, known as the Plรผcker or Grassmann coordinates of the subspace. The totality of points of $`๐(^2(V))`$ obtained in this way from lines of $`๐(V)`$ forms the set with equation $`p_{01}p_{23}p_{02}p_{13}+p_{03}p_{12}=0`$, called the Klein Quadric. The totally isotropic $`2`$-dimensional subspaces of $`V`$, namely the lines of $`L`$, correspond to those points of the Klein quadric which satisfy the additional linear equation $`p_{03}=p_{12}`$. Recalling that $`\mathrm{}_0=e_0,e_1`$, the set $`L_1`$ is the subset of $`L`$ given by $`p_{23}0`$, so taking into consideration the quadratic relation, we see that $`L_1`$ consists of the points of $`๐(^2(V))`$ which have Plรผcker coordinates $`(z^2+xy:x:z:z:y:1)`$, hence is in bijection with $`๐
_{q}^{}{}_{}{}^{3}`$. Next we consider when $`(a:b:c:1)P_1`$ is contained in $`(z^2+xy:x:z:z:y:1)L_1`$. Suppose the latter is spanned by points with homogeneous coordinates $`(a_0:a_1:a_2:a_3)`$ and $`(b_0:b_1:b_2:b_3)`$. The given point and line are incident if and only if all $`3\times 3`$ minors of the matrix
$$\left(\begin{array}{cccc}a& b& c& 1\\ a_0& a_1& a_2& a_3\\ b_0& b_1& b_2& b_3\end{array}\right)$$
(14)
are zero. The four equations which result reduce to the two equations
$$z=cy+b,x=cza.$$
(15)
By a simple change of coordinates, these equations transform to (1). This shows that $`(P_1,L_1)`$ and $`(P^{},L^{})`$ are equivalent. |
warning/0506/hep-ph0506309.html | ar5iv | text | # Off-shell scattering amplitudes in the double-logarithmic approximation
## I Introduction
Because off-shell scattering amplitudes are not gauge-invariant quantities, they have not been studied as extensively as their on-shell counterparts. Nevertheless, off-shell scattering amplitudes are important to calculate as they arise in many well-known cases, such as when the process of interest can be divided into sub-processes. One example is the $`e^+e^{}`$ annihilation into off-shell unstable particles (a quark pair, or a $`W^+W^{}`$\- pair) followed by the decay of the off-shell states. Other examples include the Bethe-Salpeter equation which relates off-shell amplitudes to on-shell amplitudes, and the DGLAPdglap evolution equations which operate with off-shell amplitudes. The Deep-Inelastic Scattering (DIS) structure functions provide yet another example when the initial parton is not assumed to be on the mass shell.
In general, the expressions for the on-shell and off-shell scattering amplitudes are quite different and should be calculated independently. There exist certain cases where it is possible to relate off-shell and on-shell amplitudes; however, even if we limit our scope to the the double logarithmic approximation (DLA), there are many cases where the on-shell amplitude cannot be obtained by a simple limit of the off-shell amplitude.
In order to demonstrate the relation between the off-shell and on-shell amplitudes, let us consider the well-known example of the Sudakov form factor in QED.sud When the virtualities of the initial $`(p_1^2)`$ and the final $`(p_2^2)`$ electrons are much greater than the electron mass squared $`(|p_{1,2}^2|m^2)`$, the expression for the off-shell Sudakov form-factor $`S(q^2)`$ of electron in QED is:
$$S(q^2)_{off}=\mathrm{exp}[(\alpha /2\pi )\mathrm{ln}(q^2/|p_1^2|)\mathrm{ln}(q^2/|p_2^2|)]$$
(1)
where $`q`$ is the momentum transfer: $`q=p_2p_1`$. We note that Eq. (1) accounts for all DL contributions to the first part $`(\gamma _\mu )`$ of the photon-electron vertex<sup>1</sup><sup>1</sup>1The expressions for the DL asymptotics of the second part ($`\sigma _{\mu \nu }`$) of the photon-fermion vertex are obtained in Ref. et in the limit $`|q^2||p_{1,2}^2|`$. Conversely, when the electron is on-shell before and after the scattering process, the Sudakov form-factor is
$$S(q^2)_{on}=\mathrm{exp}[(\alpha /4\pi )\mathrm{ln}^2(q^2/m^2)].$$
(2)
The minus sign in the logarithmic terms $`\mathrm{ln}(q^2)`$ of Eqs. (1,2) is related to the analyticity in the $`q^2`$-plane; indeed, the Sudakov form factor should not have an imaginary part when $`q^2<0`$. To simplify the notation, we will drop this minus sign, together with the modulus for $`|p_{1,2}^2|`$, throughout the paper.
Obviously, Eq. (2) cannot be obtained from Eq. (1) by taking the limit $`p_1^2=p_2^2=m^2`$. Hence, the on-shell and off-shell Sudakov form-factors cannot be related by a simple analytic continuation of the mass scales. This result stems from the fact that the DL contributions to the Sudakov form-factors come from infrared singularities in the integration over the virtual photon momenta; the separate mass scales $`m^2`$ and $`p_{1,2}^2`$ regulate these singularities.
In the one-loop approximation, the general expression for the DL contribution to the Sudakov form-factors can be written as:
$$W=(\alpha /4\pi )\left[\mathrm{ln}^2(q^2/m^2)\mathrm{ln}^2(p_1^2/m^2)\mathrm{ln}^2(p_2^2/m^2)+\mathrm{\Theta }(p_1^2p_2^2m^2q^2)\mathrm{ln}^2(m^2q^2/p_1^2p_2^2)\right]$$
(3)
where the value of the infrared cut-off $`\mu `$ is chosen as ($`\mu m`$). Summing the higher loop contributions for $`W`$ leads to the exponentiation of the one-loop result. The individual terms in the squared brackets in Eq. (3) are obtained from the integration over separate kinematic regions.
The calculation in DLA is especially simple when the Sudakov parametrization is used for the soft momenta $`k_\mu `$:
$$k_\mu =\alpha p_{}^{}{}_{2\mu }{}^{}+\beta p_{}^{}{}_{1\mu }{}^{}+k_\mu ,$$
(4)
with $`p_{}^{}{}_{1\mu }{}^{}=p_{1\mu }(p_1^2/q^2)p_{2\mu }`$ and $`p_{}^{}{}_{2\mu }{}^{}=p_{2\mu }(p_2^2/q^2)p_{1\mu }`$ so that $`p_{}^{}{}_{1}{}^{2}=p_{}^{}{}_{2}{}^{2}=0,2p_{}^{}{}_{1}{}^{}p_{}^{}{}_{2}{}^{}=q^2`$. Additionally, $`\alpha `$ and $`\beta `$ are parameters, and $`k_\mu `$ is the component of $`k_\mu `$ orthogonal to both $`p_1`$ and $`p_2`$.
Making use of the Sudakov parameterization, we can trace the origin of the separate terms for the one-loop form factor of Eq. (3). The first term \[$`\mathrm{ln}^2(q^2/m^2)`$\] arises from the integration over the region $`1\alpha ,\beta m^2/q^2`$ and $`\alpha \beta m^2/q^2`$ for the case of an on-shell electron. When the initial (final) electron is off-shell, the condition $`\alpha p_1^2/q^2`$ ($`\beta p_2^2/q^2`$) applies (provided the virtualities of the electron are so small that $`p_1^2p_2^2<m^2q^2`$) and we obtain the second \[$`\mathrm{ln}^2(p_1^2/m^2)`$\] and the third \[$`\mathrm{ln}^2(p_2^2/m^2)`$\] DL terms of Eq. (3). Finally, when $`p_1^2p_2^2>m^2q^2`$, the integration region is $`1\alpha p_1^2/q^2,1\beta p_2^2/q^2`$, and this leads to the last term \[$`\mathrm{\Theta }(p_1^2p_2^2m^2q^2)\mathrm{ln}^2(m^2q^2/p_1^2p_2^2)`$\] of Eq. (3).
When considering only the first three terms of Eq. (3), the off-shell $`W`$ and on-shell $`W`$ are related analytically. One can easily take the limit $`p_1^2,p_2^2m^2`$ in the off-shell $`W`$ to obtain the on-shell $`W`$. However, the last term \[$`\mathrm{\Theta }(p_1^2p_2^2m^2q^2)\mathrm{ln}^2(m^2q^2/p_1^2p_2^2)`$\] of Eq. (3) changes this situation and results in a non-trivial relation between the on-shell and off-shell form factors.
Beyond QED, problems relating $`S(q^2)_{on}`$ and $`S(q^2)_{off}`$ also exist in QCD and in the EW Standard Model. Additionally, when calculating DL scattering amplitudes in the EW Standard Model at energies $``$ 100 GeV<sup>2</sup>, the situation becomes more involved because, in addition to the virtual photon exchanges, one must include $`W`$ and $`Z`$ exchanges. In contrast to QED, the logarithmic contributions from the $`W`$ and $`Z`$ bosons are regulated by their masses ($`M_W`$ and $`M_Z`$); therefore, one should not introduce an infrared cut-off for them.
To simplify our notation, we introduce the mass scale $`M`$
$$MM_ZM_W,$$
(5)
which we will use to treat the $`W`$ and $`Z`$ exchange processes in a more symmetric manner. We stress that introducing $`M`$ is purely a technical point, which is useful for performing an all-order summations of DL contributions in the Standard Model.
Most of the available DL calculations of EW on-shell amplitudes have been performed in the one-loop and two-loop approximation in the EW couplings,e ; bw An all-orders DL summation for the EW on-shell amplitudes is also well-studied for the case of the hard kinematics where the total resummation of DL contributions exponentiatesflmm ; k . The more complicated case of the Regge kinematics was studied in detail in Refs. b .
In the present paper, we relate the off-shell and on-shell double-logarithmic scattering amplitudes in QED, QCD, and the EW Standard Model. We obtain explicit expressions for the off-shell double-logarithmic scattering amplitudes in the Feynman gauge (as the off-shell amplitudes are gauge-dependent), and we also discuss the Coulomb gauge.
The organization of the paper is as follows. In Sect. 2 we relate the double-logarithmic off-shell and on-shell amplitudes in QED in the hard kinematic region. In Sec. 3 we consider more involved examples of the off-shell $`22`$ scattering amplitudes where both virtual photons and $`W,Z`$ bosons are accounted to all orders in the electroweak couplings. In Sect. 4 we generalize the technique of Sect. 2 to the case of the forward Regge kinematics, and obtain explicit expressions for the DIS non-singlet structure function $`g_1`$ with the initial quark being off-shell. Specifically, we examine the effect of the initial quark virtuality on the small-$`x`$ asymptotics of $`g_1`$. Finally, in Sect. 5 we present our conclusions.
## II Relations between on-shell and off-shell amplitudes in QED
In the present section, we obtain a general prescription for relating off-shell and on-shell scattering amplitudes in the DLA for the (simplest) case of QED. We consider a $`22`$ process in QED, where the initial and final particles are fermions. Although we derive our results for a general process, we will use the $`e^+e^{}`$ annihilation into quark or lepton pairs \[e.g., $`e^{}(p_1)e^+(p_2)q^{}(p_3)q^+(p_4)`$\] as our โmodelโ process; the final formula however are valid for any channel. For this $`22`$ process, we work in the hard kinematic limit where:
$$s=(p_1+p_2)^2t=(p_3p_1)^2u=(p_1p_4)^2|p_{1,2}^2|,|p_{3,4}^2|.$$
(6)
We will use the Infrared Evolution Equations (IREE) to resum the DL contributions.<sup>2</sup><sup>2</sup>2For a recent alternate approach, see Refs. k ; ciaf ; kur . The IREE method has successfully been used to calculate the scattering amplitudes of various processes in QED, QCD and the EW Standard Model.<sup>3</sup><sup>3</sup>3See Refs. et ; flmm ; b ; kl ; efl ; egt1 , and the references therein. The first step is to use the IREE to factorize the DL contributions from the โsoftestโ virtual particle, which we define to be the particle with the minimal transverse momentum $`k_{}`$.
For QED scattering amplitudes in the hard kinematic limit Eq. (6), the softest particles must be photons. Their DL contributions can be factorized according to the Gribov bremsstrahlung theorem.<sup>4</sup><sup>4</sup>4The generalization of Gribov bremsstrahlung theoremg to QCD discussed in Refs. efl ; ce ; kl . The integration region over the softest photon momentum contains an infrared singularity, and therefore we must introduce an infrared cut-off, $`\mu `$. We impose this cut-off in the transverse momentum space by restricting the transverse momenta of all virtual particles to be greater than $`\mu `$. This constraint modifies the boundaries of the integration region, and results in relations between the cut-off scale $`\mu `$ and the masses (virtualities) of the external particles.
### II.1 The case of the on-shell amplitude $`A(s,\mu ^2)`$ with massless external particles
We begin with the well-know case where all the external particles are on-shell:
$$p_{1,2}^2=m_{1,2}^2,p_{3,4}^2=m_{3,4}^2.$$
(7)
For simplicity, we take the case where the fermion masses are negligible:
$$\mu ^2>m_{1,2}^2,m_{3,4}^2.$$
(8)
For this case, in the hard kinematic region the amplitude $`A`$ only depends on $`s`$ and $`\mu ^2`$, and we can easily obtain the IREE for this process. The RHS of the IREE includes the Born term, $`A^{Born}`$, and the contributions of the graphs depicted in Fig. 1. These graphs represent all the possible ways to factorize the softest photon, providing the Feynman gauge is used. Using the standard Feynman rules, we obtain at the following IREE:
$$A(s,\mu ^2)=A^{Born}2\lambda Q^2_{\mu ^2}^s\frac{dk_{}^2}{k_{}^2}\mathrm{ln}(s/k_{}^2)A(s,k_{}^2),$$
(9)
where $`\lambda =\alpha /8\pi `$, and $`Q`$ is the modulus of the electric charge of the fermions. In Eq. (9) recall that $`k_{}`$ plays the role of a new cut-off for any other virtual particles inside the blobs in Fig. 1,
The logarithmic factor in the integrand of Eq. (9) arises from the integration over the longitudinal Sudakov variable $`\beta `$ \[see Eq. (4)\] of the factorized photon in the region $`\mu ^2k_{}^2s\beta s`$. The integration region over $`\beta `$ and $`k_{}^2`$ is especially simple in terms of the variables $`\beta `$ and $`\alpha k_{}^2/s\beta `$, and this region is depicted in Fig. 2a. Differentiation with respect to $`\mathrm{ln}\mu ^2`$ converts the integral equation Eq. (9) into a differential one:
$$A(x)/x=2\lambda Q^2xA(x)$$
(10)
with the solution
$$A(s,\mu ^2)=A^B\mathrm{exp}(\lambda Q^2x^2).$$
(11)
Here, we have introduced the following notations: $`x=\mathrm{ln}(s/\mu ^2)`$, $`\lambda =\alpha /(8\pi )`$, and $`Q^2`$ is defined to be:
$$Q^2=Q_1^2+Q_2^2+Q_3^2+Q_4^2,$$
(12)
with $`Q_i(i=1,2,3,4)`$ being the modulus of the electric charge of the initial and final particles in units of $`e=\sqrt{4\pi \alpha }`$. As we are using the Feynman gauge, the DL contributions from the graphs depicted in Fig. 1 yield another constraint for $`Q^2`$:
$$Q^2=Z_a+Z_b+Z_c+Z_dZ_eZ_f$$
(13)
where the subscripts $`a,..,f`$ refer to the graphs $`a,..,f`$ in Fig. 1. Furthermore, we have the additional relations
$`Z_a`$ $``$ $`2Q_1Q_2`$
$`Z_b`$ $``$ $`2Q_3Q_4`$
$`Z_c`$ $``$ $`2Q_1Q_3`$
$`Z_d`$ $``$ $`2Q_2Q_4`$ (14)
$`Z_e`$ $``$ $`2Q_1Q_4`$
$`Z_f`$ $``$ $`2Q_2Q_3.`$
Finally, making use of electric charge conservation<sup>5</sup><sup>5</sup>5In the explicit formula of this paper we refer to the annihilation channel, although the final results are true in any channel.
$$Q_1Q_2Q_3+Q_4=0.$$
(15)
we can rewrite $`Q^2`$ in the form of Eq. (12).
### II.2 The case of the on-shell amplitude $`\stackrel{~}{A}(x,y_1,y_2)`$ with massive external particles
When $`\mu `$ is less than some of the masses of the external particles, the situation becomes more complex. We focus here on the specific case where:
$$m_{1,2}<\mu <m_{3,4}.$$
(16)
In this case, the scattering amplitude $`\stackrel{~}{A}`$ depends on $`m_{3,4}^2`$ in addition to $`s`$ and $`\mu ^2`$. The integration region in the $`\{\alpha ,\beta \}`$ plane<sup>6</sup><sup>6</sup>6Recall $`\alpha k_{}^2/s\beta `$, and $`s\beta `$ is the longitudinal momentum of the softest photon. is depicted in Fig. 2b; note how this differs from the integration region for Eq. (9) shown in Fig. 2a. Performing the $`\beta `$ integration over the region of Fig. 2b, and differentiating the remaining integral with respect to $`\mu `$, we arrive at the new IREE in the differential form:
$$\stackrel{~}{A}/x+\stackrel{~}{A}/y_1+\stackrel{~}{A}/y_2=2\lambda (Q^2xZ_Ay_1/2Z_By_2/2)\stackrel{~}{A},$$
(17)
with $`y_{1,2}=\mathrm{ln}(m_{3,4}^2/\mu ^2)`$, $`Z_A=(Z_b+Z_cZ_f)=2Q_3^2`$, and $`Z_B=Z_b+Z_dZ_e=2Q_4^2`$. The general solution to Eq. (17) is
$$\stackrel{~}{A}(x,y_1,y_2)=\stackrel{~}{\mathrm{\Phi }}(xy_2,y_1y_2)\mathrm{exp}\left[\lambda y_2\left(2Q^2x2Z_Ay_1+(Q^2+Z_AZ_B)y_2\right)\right].$$
(18)
For $`y_2=\mathrm{ln}(m_4^2/\mu ^2)=0`$ (i.e., $`m_4^2=\mu ^2`$) there should be an obvious matching condition
$$\stackrel{~}{A}(x,y_1,y_2)|_{y_2=0}=\stackrel{~}{A^{}}(x,y_1),$$
(19)
where $`\stackrel{~}{A^{}}(x,y_1)`$ is the amplitude of the same process, but for the case where $`m_{1,2}<m_4<\mu <m_3`$. This matching allows us to fix the unknown function $`\stackrel{~}{\mathrm{\Phi }}(xy_2,y_1y_2)`$, providing the amplitude $`\stackrel{~}{A^{}}(x,y_1)`$ is known. Taking the $`y_2=0`$ limit of Eq. (17), it is easy to see that $`\stackrel{~}{A^{}}(x,y_1)`$ obeys the following IREE:
$$\stackrel{~}{A^{}}/x+\stackrel{~}{A^{}}/y_1=2\lambda (Q^2xZ_Ay_1/2)\stackrel{~}{A}$$
(20)
with the general solution
$$\stackrel{~}{A^{}}(x,y_1)=\stackrel{~}{\mathrm{\Phi }^{}}(xy_1)\mathrm{exp}\left[\lambda y_1\left(2Q^2x2Z_Ay_1+(Q^2Z_A)y_1\right)\right].$$
(21)
In order to determine the function $`\stackrel{~}{\mathrm{\Phi }^{}}(xy_1)`$, we again apply the matching condition for $`y_1=\mathrm{ln}(m_3^2/\mu ^2)=0`$ ($`m_3^2=\mu ^2`$)
$$\stackrel{~}{A^{}}(x,y_1)|_{y_1=0}=A(x),$$
(22)
where the amplitude $`A(x)`$ is defined in Eq. (11). The matching condition of Eq. (22) immediately leads to the solution for $`\stackrel{~}{A^{}}(x,y_1)`$:
$$\stackrel{~}{A^{}}(x,y_1)=A(xy_1)\mathrm{exp}\left[\lambda y_1\left(2Q^2x2Z_Ay_1+(Q^2Z_A)y_1\right)\right].$$
(23)
Inserting this result in Eq. (19), we can obtain $`\stackrel{~}{A}(x,y_1,y_2)`$ in terms of $`\stackrel{~}{A^{}}(x,y_1)`$
$$\stackrel{~}{A}(x,y_1,y_2)=\stackrel{~}{A^{}}(xy_2,y_1y_2)\mathrm{exp}\left[\lambda y_2\left(2Q^2x2Z_Ay_1+(Q^2+Z_AZ_B)y_2\right)\right].$$
(24)
Finally, after performing all the substitutions, we find:
$$\stackrel{~}{A}(x,y_1,y_2)=A^{Born}\mathrm{exp}\left[\lambda \left(Q^2x^2Z_Ay_1^2/2Z_By_2^2/2\right)\right].$$
(25)
In general, one can use Eq. (19) when $`m_3>m_4`$.
### II.3 The case of $`A_1(s,\mu ^2,p_3^2)`$ with one off-shell final-state particle
Let us consider now the situation when one of the final-state particles is off-shell:
$$sp_3^2\mu ^2.$$
(26)
We denote the amplitude for this case as $`A_1(s,\mu ^2,p_3^2)`$. Let us replace this by: In the DLA, this amplitude does not depend on $`m_3`$; therefore $`m_3`$ can be dropped from Eq. (8) and Eq. (16) relating the value of the cut-off and the fermion masses. However, in order to use the matching between the off-shell and on-shell amplitudes in the simplest manner, we will assume that $`\mu `$ obeys Eq. (8). In contrast to the case of the on-shell amplitudes, the softest photon can now be factorized out of $`A_1(s,\mu ^2,p_3^2)`$ when $`\mu ^2k_{}^2p_3^2`$. As was observed in the previous case, the $`\beta `$integration region of the longitudinal momentum $`s\beta `$ of the softest photon momentum is now different for each individual graph in Fig. 1. Specifically, when the propagator of the photon connects to an on-shell external line, the integration region (depicted in Fig. 2a) yields the factor $`\mathrm{ln}(s/k_{}^2)`$. Conversely, when the off-shell line with momentum $`p_3`$ is involved, the integration region (depicted in Fig. 3a) yields instead the factor $`\mathrm{ln}(s/p_3^2)`$. This leads to the following integral form for the IREE:
$`A_1(s/\mu ^2,p_{}^{}{}_{}{}^{2}/\mu ^2)`$ $`=`$ $`A_1^{Born}2\lambda [Z_C{\displaystyle _{\mu ^2}^{p_{}^{}{}_{}{}^{2}}}{\displaystyle \frac{dk_{}^2}{k_{}^2}}\mathrm{ln}(s/k_{}^2)A_1(s/k_{}^2,p_3^2/k_{}^2)`$
$`=`$ $`+Z_A{\displaystyle _{\mu ^2}^{p_{}^{}{}_{}{}^{2}}}{\displaystyle \frac{dk_{}^2}{k_{}^2}}\mathrm{ln}(s/p_3^2)A_1(s/k_{}^2,p_3^2/k_{}^2)]`$
Here, we have defined $`Z_C=Q^2Z_A`$ where $`Z_A`$ was introduced in Eq. (17). Differentiating Eq. (II.3) with respect to $`\mu ^2`$, we obtain a partial differential equation for $`A_1(s,\mu ^2,p_3^2)`$
$$A_1/x+A_1/z_1=2\lambda (Q^2xZ_Az_1)A_1$$
(28)
where we have introduced the new variable $`z_1=\mathrm{ln}(p_3^2/\mu ^2)`$. The solution to Eq. (28) is given by
$$A_1(x,z_1)=\mathrm{\Phi }(xz_1)\mathrm{exp}[2\lambda Q^2(xz_1)z_1+\lambda Z_Az_1^2)]$$
(29)
where the function $`\mathrm{\Phi }`$ is, as yet, unknown. In order to determine $`\mathrm{\Phi }`$, we use the matching conditions at $`z_1=\mathrm{ln}(p_3^2/\mu ^2)=0`$ (i.e., $`p_3^2=\mu ^2`$)
$$A_1(x,z_1)|_{z_1=0}=A(x).$$
(30)
With this relation, we then obtain
$$A_1(x,z_1)=A(xz_1)\mathrm{exp}[2\lambda Q^2(xz_1)z_1\lambda Z_Cz_1^2)]$$
(31)
which relates the off-shell amplitude $`A_1`$ to the on-shell amplitude $`A`$. Using Eq. (11), we obtain an explicit expression for $`A_1(x,z_1)`$:
$$A_1(x,z_1)=A^{Born}\mathrm{exp}\left[\lambda (Q^2x^2Z_Az_1^2)\right].$$
(32)
### II.4 The case of $`A_2(s,p_3^2,p_4^2,\mu ^2)`$ and $`\stackrel{~}{A}_2(x,y_1,y_2)`$ with two off-shell final-state particles
Let us consider now the amplitude $`A_2(s,p_3^2,p_4^2,\mu ^2)`$ where both the final-state particles are off-shell. The amplitude $`A_2(s,p_3^2,p_4^2,\mu ^2)`$ will obey different IREE depending on the value of the virtualities; we will address these in turn.
#### II.4.1 Moderately-Virtual Kinematics: $`A_2(s,p_3^2,p_4^2,\mu ^2)`$
We first consider the case where $`p_{3,4}^2>\mu ^2`$ but with the condition
$$p_3^2p_4^2<s\mu ^2.$$
(33)
We will refer to this as the moderately-virtual case. The IREE for $`A_2(x,z_1,z_2)`$ is similar to Eq. (28):
$$A_2/x+A_2/z_1+A_2/z_2=2\lambda (Q^2xZ_Az_1Z_Bz_2)A_2,$$
(34)
where we define $`z_{1,2}\mathrm{ln}(p_{3,4}^2/\mu ^2)`$, and $`Z_{A,B}`$ is introduced in Eq. (17). Again, we solve this equation and using the matching conditions at $`z_2=0`$ (i.e., $`p_3^2=\mu ^2`$)
$$A_2(x,z_1,z_2)|_{z_2=0}=A_1(x,z_1),$$
(35)
and eventually obtain
$$A_2(x,z_1,z_2)=A^{Born}\mathrm{exp}\left[\lambda \left(Q^2x^2Z_Az_1^2Z_Bz_2^2\right)\right].$$
(36)
#### II.4.2 Deeply-Virtual Kinematics: $`\stackrel{~}{A}_2(x,y_1,y_2)`$
Finally, let us consider the case of the deeply-virtual kinematics
$$p_3^2p_4^2>s\mu ^2.$$
(37)
We denote the deeply-virtual amplitude as $`\stackrel{~}{A}_2(x,y_1,y_2)`$. As the integration in graph (b) of Fig. 1 does not depend on $`\mu `$, the IREE is similar to that of Eq. (34), with the exception of the contribution proportional to $`Z_b`$:
$$\stackrel{~}{A}_2/x+\stackrel{~}{A}_2/z_1+\stackrel{~}{A}_2/z_2=2\lambda (\stackrel{~}{Q}^2x\stackrel{~}{Z}_Az_1\stackrel{~}{Z}_Bz_2)\stackrel{~}{A}_2.$$
(38)
Here, the replacement $`\stackrel{~}{Q}^2=Q^2Z_b,\stackrel{~}{Z}_{A,B}=Z_{A,B}Z_b`$ should be made because the integration of graph (b) in Fig. 1 is performed over the region depicted in Fig. 3b. The solution of Eq. (38) is
$$\stackrel{~}{A}_2(x,y_1,y_2)=\stackrel{~}{\mathrm{\Phi }}_2(xz_1,xz_2)\mathrm{exp}\left[\lambda \left(\stackrel{~}{Q}^2x^2\stackrel{~}{Z}_Az_1^2\stackrel{~}{Z}_Bz_2^2\right)\right].$$
(39)
We determine the unknown function $`\stackrel{~}{\mathrm{\Phi }}_2`$ using the matching condition at $`p_3^2p_4^2=s\mu ^2`$ (i.e., when $`x=z_1+z_2`$), and obtain:
$$\stackrel{~}{A}_2(x,y_1,y_2)|_{x=z_1+z_2}=A_2(x,z_1,z_2).$$
(40)
For the function $`\stackrel{~}{\mathrm{\Phi }}_2`$ we obtain
$$\stackrel{~}{\mathrm{\Phi }}(xz_1,xz_2)=\mathrm{exp}\left[2\lambda Z_b(xz_1)(xz_2)\right],$$
(41)
and therefore the deeply-virtual amplitude $`\stackrel{~}{A}_2(x,z_1,z_2)`$ is given by:
$`\stackrel{~}{A}_2(x,z_1,z_2)`$ $`=`$ $`A\mathrm{exp}\left[\lambda [\stackrel{~}{Q}^2x^2\stackrel{~}{Z}_Az_1^2\stackrel{~}{Z}_Bz_2^2+2Z_b(xz_1)(xz_2)]\right]`$
$`=`$ $`A^{Born}\mathrm{exp}\left[\lambda [Q^2x^2Z_Az_1^2Z_Bz_2^2+Z_b(xz_1z_2)^2]\right].`$
The last term in the exponent of Eq. (II.4.2) contains the factor $`Z_b`$ defined by:
$$Z_b=2Q_3Q_4=Q_3^2+Q_4^2(Q_3Q_4)^2.$$
(43)
Obviously, this factor cannot be rewritten in terms of the individual charges $`Q_{3,4}^2`$ of the external particles.
When the Coulomb gauge is used for the on-shell amplitudes, DL contributions are always proportional to $`Q_k^2`$ because they come from only the self-energy graphs.sv ; bw Conversely, the interference graphs shown in Fig. 1 (i.e., the graphs where the boson propagators connect different external lines) do not yield DL contributions. This feature makes the Coulomb gauge especially convenient for the calculations of non-Abelian gauge theories; the nontrivial group factors are nothing but the standard Casimir factors. On the contrary, when the Feynman gauge is used, the DL contributions come from the interference graphs shown in Fig. 1, whereas the self-energy graphs do not yield DL contributions.
While this simple pattern holds true for on-shell amplitudes, it becomes more complex in the case of off-shell amplitudes. For off-shell amplitudes, both the Feynman and Coulomb gauges received DL contributions from the interference graphs. One can verify (using $`\stackrel{~}{Q}^2=Z_A=Z_B=Z_b=2`$) that Eqs. (11) and (II.4.2) reproduce, as a particular case, Eqs. (2) and (1) for the on-shell and off-shell Sudakov form factors of the electron. Finally, the generalization of these results to the case with more than 4 external particles can be calculated in the DLA in a similar manner.
## III Generalization to the EW Standard Model
When the $`22`$ scattering amplitudes of Sect. 2 are considered in the framework of the ElectroWeak (EW) Standard Model, the situation becomes more involved due to introduction of two cut-off scales, $`\mu `$ and $`M`$. We will only comment briefly on the on-shell EW amplitudes in DLA, as these amplitudes have been computed in Refs. flmm ; b using the appropriate IREEโs. We will examine in detail the $`22`$ EW amplitudes in the hard kinematic limit.
In order to obtain the IREE for both photon and $`W,Z`$ exchanges, one must factorize the DL contributions where the EW bosons have minimal $`k_{}`$. If the softest boson is a photon, then the lower limit of integration over $`k_{}`$ is $`\mu `$, and the result resembles the QED case discussed in the previous section. Conversely, when the softest boson is a $`W`$ or $`Z`$, then the lower limit of integration over $`k_{}`$ is $`M`$. Therefore, the integral on the RHS of Eq. (II.3) corresponding to the factorization of the softest photon should be replaced by the DL contributions of the factorization of the softest $`W`$ and $`Z`$ bosons. The contributions of the softest $`Z`$ boson looks quite similar to the photon case shown in Fig. 1; in contrast, the $`W`$ boson exchange changes the flavors of the interacting fermions.
Considerable technical difficulties arise when using either the Feynman or the Coulomb gauges.<sup>7</sup><sup>7</sup>7For DL calculations with the Coulomb gauge, see e.g. Refs. sv ; bw . In the Coulomb gauge, the DL contributions of the softest bosons come from the boson self-energy graphs, and this means that the amplitudes are the exponential form obtained in Ref. flmm . On the contrary, when the Feynman gauge is used, the factorization of the softest $`W`$-boson leads to a system of four IREEโs, the general solution of which consists of four exponentials. If one of these four exponentials were negative, this would lead to DL contributions that grow with $`s`$, and therefore violate unitarity. However, the Born contributions (which are the inhomogeneous terms in the integral IREEโs) lead to the cancellation of three of the exponentials. The fourth exponential then coincides with the solution obtained using the Coulomb gauge. Consequently, the Coulomb gauge is more convenient for calculations in the hard kinematic limit of Eq. (6), but becomes inappropriate when working in the Regge kinematic region where the Feynman gauge is more convenient. Indeed, in the Feynman gauge one can avoid the above technical problems by expanding the EW amplitude into invariant ones, as was shown in Ref. b .
### III.1 The Generalized IREE
Let us consider a $`22`$ electroweak scattering amplitude $`A(s,p_3^2,p_4^2,\mu ^2,M^2)`$ in the hard kinematic limit of Eq. (6). We use the notation of the previous Section for the momenta of the external particles. We consider several cases assuming again that the initial particles are always on-shell, whereas the final particles can be either on-shell or off-shell. For any of these cases, the amplitude obeys the following integral IREE:
$`A(s,p_3^2,p_4^2,\mu ^2,M^2)`$ $`=`$ $`A^{Born}{\displaystyle \frac{U}{16\pi ^2}}{\displaystyle _D}{\displaystyle \frac{dk_{}^2}{k_{}^2}}{\displaystyle \frac{d\beta }{\beta }}A(s,p_3^2,p_4^2,k_{}^2,M^2)`$ (44)
$`{\displaystyle \frac{g^2V}{16\pi ^2}}{\displaystyle _D^{}}{\displaystyle \frac{dk_{}^2}{k_{}^2}}{\displaystyle \frac{d\beta }{\beta }}A(s,p_3^2,p_4^2,k_{}^2,k_{}^2)`$
where $`A^{Born}`$ is the Born approximation for $`A`$, and $`g=\sqrt{4\pi \alpha }/\mathrm{sin}\theta _W`$ is the $`SU(2)`$-coupling of the EW Standard Model. The factors $`U`$ and $`V`$ depend on the particular relations between the cut-offs and virtualities $`p_{1,2}^2`$; we will specify them below. The first (second) integral in Eq. (44) corresponds to the factorization of the softest photon (weak boson). The integration regions $`D`$ and $`D^{}`$ also depend on the relations between the parameters $`p_3^2,p_4^2,\mu ^2,`$ and $`M^2`$. Generally, the solution to Eq. (44) can be written as
$$A=A^{Born}\mathrm{exp}(\psi )$$
(45)
with appropriate exponents $`\psi `$. When both final-state particles are on-shell, the integration over $`\beta `$ yields the common factor $`\mathrm{ln}(s/k_{}^2)`$ in both the $`U`$ and $`V`$ integrals of Eq. (44) even through the integration over $`k_{}^2`$ runs from $`\mu ^2`$ to $`s`$ in the first integral and from $`M^2`$ to $`s`$ in the second one. In a manner similar to Eq. (10), one obtains the on-shell factor $`U_{on}=32\pi ^2\lambda Q^2`$. The on-shell weak factor, $`C_{WZ}`$, is the sum of the factors $`V_i`$:
$$C_{WZ}=V_1+V_2+V_3+V_4,$$
(46)
where the subscripts 1 and 2 (3 and 4) refer to the initial-state particles with momenta $`p_{1,2}`$ (the final-state particles with momenta $`p_{3,4}`$). The factors $`V_i`$ are expressed in terms of the Weinberg angle $`\theta _W`$, the weak isospins $`T_i^2`$ ($`T_i^2=3/4`$ for a fermion), the hypercharges $`Y_i`$, and the electric charges $`Q_i`$ of the initial-state and final-state particles:
$$V_i=[T_i^2+(Y_i^2/4)\mathrm{tan}^2\theta _WQ_i^2sin^2\theta _W].$$
(47)
Differentiating Eq. (44) with respect to the cut-offs, we obtain the on-shell exponent $`\psi \psi _{on}`$: flmm
$$\psi _{on}=\alpha Q^2/(8\pi ^2)\mathrm{ln}^2(s/\mu ^2)+g^2/(32\pi ^2)C_{WZ}\mathrm{ln}^2(s/M^2).$$
(48)
We recall that the photon double-logarithmic term $`\mathrm{ln}^2(s/\mu ^2)`$ in Eqs. (11,48) is obtained under the assumption that the cut-off $`\mu `$ is greater than any mass involved.
In the case where the initial-state particles are light but the final-state particles are so heavy that Eq. (16) is satisfied, the term $`Q^2\mathrm{ln}^2(s/\mu ^2)`$ should be replaced by
$$Q^2\mathrm{ln}^2(s/\mu ^2)(Z_A/2)\mathrm{ln}^2(m_1^2/\mu ^2)(Z_B/2)\mathrm{ln}^2(m_2^2/\mu ^2)$$
(49)
in agreement with Eq. (25).
### III.2 The off-shell case:
When the final-state particles are off-shell, the regions $`D`$ and $`D^{}`$ in Eq. (44) are more complicated. We study below several interesting situations, denoting these as $`R_{1,2,3,4}`$.
#### III.2.1 $`R_1`$ Case: final-state particles of small virtuality
We call $`R_1`$ the simplest case when the virtualities of the final particles are relatively small:
$`\mu ^2`$ $`<`$ $`p_{3,4}^2<M^2s,`$
$`p_3^2p_4^2`$ $`<`$ $`s\mu ^2.`$ (50)
This case is quite similar to Eq. (33), so we immediately conclude that the exponent $`\psi _1`$ of Eq. (45) for this case is
$$\psi _1=\frac{\alpha }{8\pi }[Q^2\mathrm{ln}^2(s/\mu ^2)Z_A\mathrm{ln}^2(p_3^2/\mu ^2)Z_B\mathrm{ln}^2(p_4^2/\mu ^2)]+\frac{g^2}{32\pi ^2}C_{WZ}\mathrm{ln}^2(s/M^2).$$
(51)
#### III.2.2 $`R_2`$ Case: Deeply virtual for photon, and on-shell for $`W/Z`$ boson
In the second case, $`R_2`$, the virtualities $`p_{3,4}^2`$ are larger and satisfy
$`\mu ^2`$ $`<`$ $`p_{3,4}^2<M^2,`$
$`s\mu ^2`$ $`<`$ $`p_3^2p_4^2<sM^2.`$ (52)
This kinematics are deeply-virtual for the softest photon (with the solution given by Eq. (II.4.2)), but at the same time it is on-shell for the softest $`W,Z`$-bosons (with the solution given by Eq. (11)). Therefore the exponent $`\psi _2`$ for this case is
$$\psi _2=\lambda [\stackrel{~}{Q}^2x^2\stackrel{~}{Z}_Az_1^2\stackrel{~}{Z}_Bz_2^2+2Z_b(xz_1)(xz_2)]+\frac{g^2}{32\pi ^2}C_{WZ}\mathrm{ln}^2(s/M^2)$$
(53)
#### III.2.3 $`R_3`$ Case: Deeply virtual for photon, and moderately virtual for $`W/Z`$ boson
The case $`R_3`$, defined as
$`p_{3,4}^2`$ $`>`$ $`M^2`$
$`s`$ $`<`$ $`p_3^2p_4^2<sM^2,`$ (54)
describes a situation which is deeply-virtual for the photon, and moderately-virtual for the $`W,Z`$-bosons. Combining the results of Eq. (36) and Eq. (II.4.2), we find the solution:
$`\psi _3`$ $`=`$ $`\lambda [\stackrel{~}{Q}^2x^2\stackrel{~}{Z}_Az_1^2\stackrel{~}{Z}_Bz_2^2+2Z_b(xz_1)(xz_2)]`$
$`+{\displaystyle \frac{g^2}{32\pi ^2}}\left[C_{WZ}\mathrm{ln}^2(s/M^2)2V_3\mathrm{ln}^2(p_3^2/M^2)2V_4\mathrm{ln}^2(p_4^2/M^2)\right].`$
#### III.2.4 $`R_4`$ Case: Deeply virtual for photon $`W/Z`$ boson
Finally, when the momentum $`p_{3,4}^2`$ are so large that
$`p_{3,4}^2`$ $`>`$ $`M^2`$
$`p_3^2p_4^2`$ $`>`$ $`sM^2,`$ (56)
the situation is deeply-virtual for all of the softest electroweak bosons, and this case is equivalent to the scheme with a single cut-off $`M`$. The EW exponent $`\psi _4`$ for this case corresponds to the Eq. (II.4.2):
$`\psi _4`$ $`=`$ $`\left[\lambda \stackrel{~}{Q}^2\mathrm{ln}^2(s/\mu ^2)+g^2\stackrel{~}{C}_{WZ}/(32\pi ^2)\mathrm{ln}^2(s/M^2)\right]`$ (57)
$``$ $`\left[\lambda \stackrel{~}{Z}_A\mathrm{ln}^2(p_3^2/\mu ^2)+g^2\stackrel{~}{X}_3/(32\pi ^2)\mathrm{ln}^2(p_3^2/M^2)\right]`$
$``$ $`\left[\lambda \stackrel{~}{Z}_B\mathrm{ln}^2(p_4^2/\mu ^2)+g^2\stackrel{~}{X}_4/(32\pi ^2)\mathrm{ln}^2(p_4^2/M^2)\right]`$
$`+`$ $`2\left[\lambda Z_b+g^2X_b/(32\pi ^2)\right]\mathrm{ln}^2(s/p_3^2)\mathrm{ln}^2(s/p_4^2).`$
Here, $`X_{3,4}=2V_{3,4}X_b`$, and $`X_b=1/2+(1/\mathrm{cos}^2\theta _W)[t_3^{(3)}Q_3\mathrm{sin}^2\theta _W][t_3^{(4)}Q_4\mathrm{sin}^2\theta _W]`$, where $`Q_{3,4}`$ are the electric charges and $`t_3^{3,4}`$ are the eigenvalues of the $`SU(2)`$ -generator $`T_3`$ acting on the final-state particles 3 and 4, respectively. It is convenient to rewrite $`\psi _4`$ in the following form:
$`\psi _4`$ $`=`$ $`\left[\lambda Q^2\mathrm{ln}^2(s/\mu ^2)+g^2C_{WZ}/(32\pi ^2)\mathrm{ln}^2(s/M^2)\right]`$ (58)
$``$ $`\left[\lambda Z_A\mathrm{ln}^2(p_3^2/\mu ^2)+g^2X_3/(32\pi ^2)\mathrm{ln}^2(p_3^2/M^2)\right]`$
$``$ $`\left[\lambda \stackrel{~}{Z}_B\mathrm{ln}^2(p_4^2/\mu ^2)+g^2X_4/(32\pi ^2)\mathrm{ln}^2(p_4^2/M^2)\right]`$
$`+`$ $`\left[\lambda Z_b+g^2X_b/(32\pi ^2)\right]\mathrm{ln}^2(sM^2/p_3^2p_4^2).`$
Eq. (45), together with the phases specified in Eqs. (51,53,III.2.3,57,58), describe the $`22`$ off-shell electroweak scattering amplitudes in the hard kinematic limit for several virtualities of the final-state particles. The above techniques can be generalized for other particular cases in a similar manner.
## IV Off-shell amplitudes in the Regge kinematics
The evaluation of the DL contributions to the on-shell and off-shell scattering amplitudes in the hard kinematic limit (6) is relatively straightforward, and leads to the exponentiation of the one-loop DL contributions. When the scattering amplitudes are evaluated in the forward Regge limit
$$sut,$$
(59)
or backward Regge limit
$$stu,$$
(60)
the DL expressions are more complicated.
The $`22`$ EW Regge amplitudes at TeV energies for electron-positron colliders have been discussed previously in Refs. b . In the following, we will consider one example of the off-shell Regge amplitudes in the QCD theory. We obtain a relation between on-shell and off-shell amplitudes for quark scattering, and use this result to calculate the non-singlet structure functions for Deep Inelastic Scattering when the initial quark is off-shell. In doing so, we make use of the expressions for the on-shell non-singlet structure functions obtained in Ref. egt1 at small $`x`$. These expressions account for the DL and single-logarithmic (SL) contributions to all orders in the QCD coupling $`\alpha _s`$. Explicit expressions combining the DGLAPdglap expressions and our total resummation (DL and SL contributions) are obtained in Ref. egt3 . These combined results are especially important in the region of small $`x`$, as the combination leads to the power-like asymptotic behavior $`(x^{\omega _0})`$ of the structure functions when $`x0`$, with intercept $`\omega _00.4`$. We briefly summarize these on-shell results.
According to the Optical Theorem, it is possible to relate the on-shell non-singlet structure function $`F^{(\pm )}(x,Q^2)`$ to the imaginary part of the forward Compton amplitude $`A^{(\pm )}(x,Q^2)`$ with $`t\mu ^2`$. It is convenient to present $`A^{(\pm )}(x,Q^2)`$ in the form of the Mellin integral:
$$A^{(\pm )}(x,Q^2)=_ฤฑ\mathrm{}^{+ฤฑ\mathrm{}}\frac{d\omega }{2\pi ฤฑ}\left(\frac{s}{\mu ^2}\right)^\omega \xi ^{(\pm )}(\omega )T^{(\pm )}(\omega ,Q^2),$$
(61)
with the signature factor $`\xi ^{(\pm )}(\omega )=(e^{ฤฑ\pi \omega }\pm 1)/2`$, total energy $`s=Q^2/x`$, and the infrared cut-off $`\mu `$. The Mellin amplitude $`T^{(\pm )}(\omega ,Q^2)`$ obeys the following IREE:
$$T^{(\pm )}/y+\omega T^{(\pm )}=T^{(\pm )}H^{(\pm )}(\omega ),$$
(62)
with anomalous dimensions $`H^{(\pm )}`$, and $`y=\mathrm{ln}(Q^2/\mu ^2)`$ so that $`\mu ^2`$ is the starting point of the $`Q^2`$-evolution. $`H^{(\pm )}`$ include both double-logarithmic and the most important part of the single-logarithmic contributions. From Refs. egt1 we have
$$H^{(\pm )}=(1/2)\left[\omega \sqrt{\omega ^2B^\pm (\omega )}\right]$$
(63)
where
$$B^{(\pm )}(\omega )=[(1+\omega /2)4\pi C_FA(\omega )+D^{(\pm )}(\omega )]/(2\pi ^2).$$
(64)
In Eq. (64), $`D^{(\pm )}(\omega )`$ and $`A(\omega )`$ can be expressed in terms of $`\rho =\mathrm{ln}(1/x)`$, $`b=(332n_f)/12\pi `$, and the color factors $`C_F=4/3`$, $`N=3`$:
$$D^{(\pm )}(\omega )=\frac{2C_F}{b^2N}_0^{\mathrm{}}๐\eta e^{\omega \eta }\mathrm{ln}\left(\frac{\rho +\eta }{\eta }\right)\left[\frac{\rho +\eta }{(\rho +\eta )^2+\pi ^2}\frac{1}{\eta }\right],$$
(65)
$$A(\omega )=\frac{1}{b}\left[\frac{\eta }{\eta ^2+\pi ^2}_0^{\mathrm{}}\frac{d\rho e^{\omega \rho }}{(\rho +\eta )^2+\pi ^2}\right].$$
(66)
In Eq. (66), $`A(\omega )`$ is the Mellin transform of $`\alpha _s(k^2)=1/(b\mathrm{ln}(k^2/\mathrm{\Lambda }^2))`$ with time-like argument $`k^2`$. Solving Eq. (62) and introducing $`F^{(\pm )}(x,Q^2)`$,
$$F^{(\pm )}(x,Q^2)=\frac{1}{\pi }\mathrm{}A^{(\pm )}(x,Q^2),$$
(67)
we obtain
$$F^{(\pm )}(x,Q^2)=(e_q^2/2)_ฤฑ\mathrm{}^{+ฤฑ\mathrm{}}\frac{d\omega }{2\pi ฤฑ}(1/x)^\omega C^{(\pm )}(\omega )\delta q(\omega )\mathrm{exp}\left(H^{(\pm )}(\omega )y\right).$$
(68)
Eq. (67) relates $`F^{(+)}`$ and the forward amplitude $`A^{(+)}`$ via the Optical theorem. In fact, $`F^{(+)}`$ is the non-singlet contribution to the structure function $`F_1`$, and $`F^{()}`$ is the non-singlet contribution to the polarized structure function $`g_1`$. $`\delta q`$ is the initial quark density, which is commonly determined from fitting the experimental data. The coefficient functions $`C^{(\pm )}`$ can be expressed in terms of the anomalous dimensions $`H^{(\pm )}`$ as:
$$C^{(\pm )}=\frac{\omega }{\omega H^{(\pm )}(\omega )}.$$
(69)
We now consider the case where the initial quark is off-shell with virtuality $`p^2`$,<sup>8</sup><sup>8</sup>8As in the previous Sections, we drop the modulus sign in Eq. (70), keeping the notation $`p^2`$ instead of $`|p^2|`$.
$$\mu ^2<|p^2|Q^2.$$
(70)
When the initial quark is off-shell, the forward Compton amplitude $`\stackrel{~}{A}^{(\pm )}`$ depends on $`p^2`$ as well: $`\stackrel{~}{A}^{(\pm )}=\stackrel{~}{A}^{(\pm )}(s,Q^2,p^2,\mu ^2)`$. As before, it is convenient to work with the Mellin off-shell amplitude $`G(\omega ,Q^2,p^2,\mu ^2)`$ which is related to $`\stackrel{~}{A}^{(\pm )}(s,Q^2,p^2,\mu ^2)`$ through the Mellin transform (61). The IREE for the off-shell amplitude $`G(\omega ,Q^2,p^2,\mu ^2)`$ is similar to the one for the on-shell amplitude $`T^{(\pm )}`$:
$$G^{(\pm )}(\omega ,y,z)/y+G^{(\pm )}(\omega ,y,z)/z+\omega G^{(\pm )}(\omega ,y,z)=T^{(\pm )}(\omega ,y)h^{(\pm )}(\omega ,z)$$
(71)
where $`z=\mathrm{ln}(p^2/\mu ^2)`$. $`h^{(\pm )}(\omega ,z)`$ are new anomalous dimensions which can be found from the following IREE:
$$h^{(\pm )}/z+\omega h^{(\pm )}=D^{(\pm )}+h^{(\pm )}H^{(\pm )}(\omega ),$$
(72)
with $`D^{(\pm )}`$ defined in Eq. (64).
Solving Eq. (72), and using the matching conditions for $`z=0`$
$$h^{(\pm )}(\omega ,z)|_{z=0}=H^{(\pm )}(\omega )$$
(73)
we obtain
$$h^{(\pm )}=\left[H^{(\pm )}+_0^z๐ue^{(\omega H^{(\pm )})u}D^{(\pm )}(u)\right]e^{(\omega H^{(\pm )})z}.$$
(74)
Once the $`h^{(\pm )}`$ are known, we can calculate $`G^{(\pm )}`$. Substituting $`h^{(\pm )}`$ and $`T^{(\pm )}`$ into Eq. (71) and solving, we obtain
$$G^{(\pm )}(\omega ,y,z)=T^{(\pm )}(\omega ,2\eta )e^{\omega z}+e^{\omega l}_\eta ^l๐ue^{\omega u}T^{(\pm )}(\omega ,u+\eta )h^{(\pm )}(u\eta )$$
(75)
where $`l=(y+z)/2,\eta =(yz)/2`$. Eq. (75) gives $`G^{(\pm )}(\omega ,y,z)=T^{(\pm )}(\omega ,y)`$ when $`z=0`$.
To demonstrate a simple application of Eq. (75), let us estimate the effect of the virtuality $`p^2`$ of the initial quark on the small-$`x`$ asymptotics of the off-shell non-singlet structure functions $`\stackrel{~}{F}^{(\pm )}`$. Given the results of Refs. egt1 , when $`x0`$ the small-$`x`$ asymptotics of the on-shell structure functions $`F^{(\pm )}`$ are Regge-like:
$$F^{(\pm )}(x,Q^2)e_q^2\delta q(\omega _0^{(\pm )})c^{(\pm )}\xi ^{\omega _0^{(\pm )}}/\mathrm{ln}^{3/2}\xi $$
(76)
where $`\xi =\sqrt{Q^2/(x^2\mu ^2)}`$, $`c^{(\pm )}=\left[2(1q)/(\pi \sqrt{B^{(\pm )}})\right]^{1/2}`$, $`q^{(\pm )}=d\sqrt{B^{(\pm )}}/d\omega |_{\omega =\omega _0^{(\pm )}}`$, and the intercepts are $`\omega _0^{(+)}=0.38`$ and $`\omega _0^{()}=0.43`$. To estimate the asymptotics of $`\stackrel{~}{F}^{(\pm )}(x,Q^2,p^2)`$, we neglect $`D^{(\pm )}`$ in Eq. (74) since the effect of $`D^{(\pm )}`$ on the small-$`x`$ behavior of $`g_1`$ is not large.egt1 Performing the integration in Eq. (75) using the saddle-point method, we obtain at the following expression for the small-$`x`$ asymptotics of the off-shell structure functions $`\stackrel{~}{F}^{(\pm )}`$:
$$\stackrel{~}{F}^{(\pm )}(x,Q^2,p^2)F^{(\pm )}(x,Q^2)e^{\omega _0z/2}\left[1+\omega _0^{(\pm )}z/2\right],$$
(77)
with $`\omega _0^{(+)}=0.38`$ and $`\omega _0^{()}=0.43`$.<sup>9</sup><sup>9</sup>9See Refs. egt1 for details. In order to estimate the difference between the on-shell and the off-shell non-singlet structure functions, we define the deviation $`R`$:
$$R=\left(F^{(\pm )}\stackrel{~}{F}^{(\pm )}\right)/F^{(\pm )}1e^{0.2\mathrm{ln}(p^2/\mu ^2)}\left[1+0.2\mathrm{ln}(p^2/\mu ^2)\right],$$
(78)
where we estimate $`\mu `$ to be $`1`$ GeV.egt1 We observe from Eq. (78) that $`R`$ increases with increasing $`p^2`$. By definition, $`R=0`$ when $`p^2=1`$ GeV<sup>2</sup> and therefore it grows when $`p^2`$ becomes greater than $`\mu ^2`$ (we remind that the notation $`p^2`$ actually denotes $`|p^2|`$ ). In particular, $`R0.02`$ for $`|p^2|=3`$ GeV<sup>2</sup> and $`R0.08`$ for $`|p^2|=10`$ GeV<sup>2</sup>.
## V Conclusions
In this paper, we have presented an approach for calculating the off-shell scattering amplitudes in the double-logarithmic (DL) approximation. We have considered, in particular, the case where the initial particles are on-shell, but some of the final particles can be off-shell. Technically, this approach is based on constructing the appropriate infrared evolution equations (IREE) for the off-shell amplitudes. The specific form of the solutions depends on the size of the virtualities of the final particles. The general strategy for obtaining the off-shell solutions was to use the matching conditions to relate amplitudes of the same process, but with smaller virtualities of the final particles. Such a procedure allows us to relate the amplitudes with $`N`$ off-shell particles to those with $`N1`$ off-shell particles. Using this recursive relation, we can eventually relate the off-shell amplitudes to the on-shell amplitudes.
To demonstrate the utility of this approach, we have computed the scattering amplitudes in the hard kinematic limit of Eq. (6), where the DL corrections simply exponentiation the one-loop DL contribution. We have reproduced the well-known results for the $`22`$ on-shell QED amplitudes of Eqs. (11,25), and used these results to obtain explicit expressions for the amplitudes with one (Eq. (32)) and two (Eqs. (36,II.4.2)) final particles off-shell. We have considered both cases of moderate and large virtualities for the final particles, and compared the results in the Feynman and Coulomb gauges. In particular, we have shown that both the self-energy and interference graphs yield DL contributions in the Coulomb gauge; this result is contrary to the on-shell case. This result leads to the conclusion that the Coulomb gauge does not have provide any technical advantage, compared to the Feynman gauge, for calculating the off-shell amplitudes.
In Sect. 3 we have computed the more complex case of the off-shell amplitudes for the $`22`$ electroweak processes in the hard kinematic limit at energies $`100`$ GeV. In addition to the virtual photon exchanges, we also took the exchanges of virtual $`W`$ and $`Z`$ bosons into account. The results of the DL contributions for these scattering amplitudes depends on the virtualities of the final particles; the results are presented in Eqs. (51,53,III.2.3,57).
Finally, we have studied the QCD forward Compton scattering amplitude in the Regge kinematic limit (Eqs. (59,60)) with both the photons and the quarks being off-shell. This result was used to calculate the DIS non-singlet structure functions with off-shell initial quarks (Eq. (75)). To estimate the difference between the off-shell and the on-shell structure functions, we compared their small-$`x`$ asymptotics (Eqs. (76,77)), and the result is presented in Eq. (78). The on-shell and off-shell asymptotic amplitudes differ at larger values of $`p^2`$; for example, this difference is $`8\%`$ for $`p^2=10GeV^2`$. A more accurate estimate of the effect of the virtuality on the structure functions could be directly obtained from the explicit expression of Eq. (75), instead of using their asymptotic behavior of Eq. (77); this topic will be reserved for future study.
## VI Acknowledgment
We are grateful to A. Barroso for useful discussions. The work is supported in part by grant RSGSS-1124.2003.2, U.S. DoE grant DE-FG03-95ER40908, and the Lightner-Sams Foundation. |
warning/0506/astro-ph0506159.html | ar5iv | text | # Wide-Angle Wind Driven Bipolar Outflows: High Resolution Models with Application to Source I of the Becklin-Neugebauer / Kleinmann-Low OMC-I Region
## 1 Introduction
Bipolar outflows are recognized as a fundamental component of the star formation process for low and intermediate mass stars. The precise nature of these outflows has received considerable attention and their observational properties have been well characterized (Richer et al. (2000)). The nature of the mechanisms driving the outflows has not, however, been determined. While there is now a consensus that the outflow is made up of swept-up ambient material driven by a โwindโ from the central source (a proto-star or young star) there remains a debate over the form such a wind will take. In some models the star produces a well collimated jet on small scales ($`<10AU`$) due (most likely) to the action of an accretion disk (Konigl & Pudritz (2000)). Other models assume that the proto-stellar wind is not as strongly collimated taking the form of a โwide-angle windโ (WAW). In many models the WAW will have a momentum distribution such that wind has dense core in the polar direction surrounded by less dense flow at lower latitudes (Shu et al. (2000)). Thus WAWs are often envisioned as having a jet-like component. The distinction between these models plays an important role in the debate about star formation because the nature of the WAW/jet is tied to the nature of the wind launching mechanism at the base. Understanding how winds are driven from an accretion disk (most likely via a coupling of magnetic fields and rotation) is a problem of great importance both because jets/winds are ubiquitous in astrophysical environments and because these flows may carry away a significant fraction of the diskโs angular momentum (Konigl & Pudritz (2000)). Thus the distinction between tightly collimated jets and wide-angle winds informs the debate about the nature of jet launching and accretion disk physics. Since it is difficult to directly observe the wind launching regions (see Woitas et al. (2005) ), the properties of large scale bipolar outflows can aid in understanding the nature of their initial conditions.
In addition to questions concerning the nature of the driving winds, bipolar outflows may play an important role in modifying their natal environments. The energy and momentum budgets associated with the aggregate of outflows in a young cluster can be large enough to either power turbulence in the cloud from which the cluster was born or, in some cases, unbind some fraction of the cloud material (Arce (2003)). Observational and theoretical studies have yet to address this issue in its full complexity and so the role of outflows as environmental factors determining cloud properties remains unclear. Recent studies of the NGC 1333 region have shown that fossil bipolar outflows or โcavitiesโ which remain after the central source has turned off provide a significant coupling agent linking wind momenta to the cloud (Quillen et al. (2004)). This work emphasizes the importance of understanding outflow properties in terms of momentum transfer processes. Understanding how the these properties operate over the entire outflow history will be critical to addressing their impact on the cloud both on intermediate scale ($`10^{16}10^{17}cm`$) and large scales ($`.1pc`$).
We note that while considerable progress has been made in understanding star formation for isolated low mass stars, the formation of high-mass stars remains less clear (see the excellent review by Shepherd (2003)). There is increasing observational work characterizing outflows in high mass star forming environments however fundamental questions such as the role of accretion disks, magnetic fields remain to be definitively answered. Thus there remains considerable work to be done in the study of bipolar outflows in the context of more massive stars. It is noteworthy that the distinction between jet driven and WAW driven outflows is even less clear in high mass stars since radiation pressure is capable of producing a significant stellar wind in many cases.
In this context observations of an organized distribution of $`Si0`$ and $`H_2O`$ maser emitting spots near Source I in the Orion BN/KL nebula are of particular interest. The exact nature of Source I is unclear, because it is so heavily embedded (Greenhill et al., 2003; Chandler & Greenhill, 2002). $`SiO`$ masers are distributed along an X-shaped locus of clumps extending 20 to 70 AU from the continuum radio source (Snyder & Buhl, 1974; Wright & Plambeck, 1983; Lane, 1982; Plambeck, Wright, & Carlstrom, 1990). Subsequent observations by Menten & Reid (1995); Greenhill et al. (1998); Doeleman, Lonsdale, & Pelkey (1999) have established that the maser sources appear to part of an outflow from source I. The maser proper motions lie in the range of 10 to 23 $`kms^1`$, oriented primarily along the limbs of the X, with systematically red-shifted and blue-shifted lobes about a southeast-northwest symmetry axis.
A number of models have been proposed to explain the pattern of maser proper motion and line of sight velocity (Plambeck, Wright, & Carlstrom (1990); Greenhill et al. (1998)). Recently Greenhill et al. (2003, 2005) have presented a model based on high resolution radio observations where the outflow is northeast-southwest so that the masers are situated along a biconical outflow with a wide opening angle. Greenhill et al. (2003) note a bridge of maser emission connecting the southern and western arms of the bicone, with a clear velocity gradient that is consistent with the edge of a rotating disk. Orienting the outflow northeast to southwest makes the ejection perpendicular to the disk, as seen in many low-mass young stellar objects, and suggests that the $`H_2O`$ masers seen on larger scales are also oriented with, and thus probably produced by, the same outflow as the SiO masers (Greenhill et al., 1998).
The radial velocity difference between the arms suggests rotation. By requiring that this rotational motion be consistent with gravitational binding, Greenhill et al. (2003) estimated the dynamical mass of source I to be $`6M_{\mathrm{}}`$.
What is the origin of this rotational motion? One possibility is that there are streams of slow-moving molecular ejecta from the outer edges of the disk (L. Greenhill, personal communication). The other possibility is that this rotation reflects the angular momentum in the ambient medium. The presence of rotation and outflow in the Source I system makes it an interesting test case for models of proto-stellar winds interacting with the surrounding media on relatively small scales ($`<70AU`$).
Motivated by these observations of source I, in this paper we present a numerical study of the interaction between a fast irrotational wind from a central source with an infalling, rotating protostellar envelope. This work is a continuation of an going study of outflow properties formed via wind/infalling envelope interactions. In previous works we have explored the roles of inflow ram pressure on outflow collimation (Delamarter, Frank, & Hartmann (2000)) as well as the role of toroidal magnetic fields in shaping the outflow (Gardiner, Frank, & Hartmann (2003)). In this work we focus on the walls of the outflow where the development of a shear layer between the infalling ambient material and the outflowing wind material is the site of strong mixing and could be the site of maser emission. We explore the dynamics of the shell walls and show that the velocities in this region are reasonably consistent with the observations of Source I.
Our use of an adaptive mesh refinement allows us to achieve high levels of resolution within the walls of the bipolar outflow. By capturing the hydrodynamics of these strongly cooling flows we can address some open issues in the physics of molecular outflows. In particular we have been able to observe instabilities occurring at the top, or head of the outflow, as well as partially resolve the flow pattern which occurs as shocked infalling material flows past shocked outflowing wind. The latter question is some importance as there has been debate in previous works on the subject as to the nature of the mixing in this region and the direction of the final bulk momentum (Delamarter, Frank, & Hartmann, 2000; Lee, Stone, Ostriker, & Mundy, 2001; Shu, Ruden, Lada, & Lizano, 1991; Wilkin & Stahler, 1998). Thus our paper addresses generic issues related to strongly cooling wind blown bubbles as well as some specific issues related directly to nature of the outflow from Source I.
In section II we discuss the numerical model and initial conditions used for the simulations as well as the assumptions and simplifications which support the model. In appendix A more detail of our micro-physics is provided. In section III we present our results. In section IV we discuss the results in light of observations of Source I and general considerations of molecular outflows. In section V we present our conclusions.
## 2 Method and Model
### 2.1 Numerical Code
We have carried out a series of radiative hydrodynamic simulations of an isotropic wind interacting with a rotating, infalling envelope. The micro-physics of H and He ionization, $`H_2`$ chemistry and optically thin cooling have been included. Our simulations are carried out in 2.5D (i.e. cylindrical symmetry) using the AstroBEAR adaptive mesh refinement (AMR) code. AMR allows high resolution to be achieved only in those regions which require it due to the presence of steep gradients in critical quantities such as gas density. The hydrodynamic version of AstroBEAR has been well tested on variety of problem in 1, 2 and 2.5D (Poludnenko et. al., 2004; Varnie et. al., 2004) The system of equations integrated are:
$$d_tQ+F=S_{geom}+S_m+S_{grav}$$
where the vector of conserved quantities $`Q`$, the flux function $`F`$, the geometric source terms $`S_{geom}`$, the micro-physical source terms $`S_m`$, and the central gravitational source terms $`S_{grav}`$ are given as:
$$Q=\left[\begin{array}{c}\rho \\ \rho v_r\\ \rho v_z\\ \rho v_\theta \\ E\\ \rho _{H_2}\\ \rho _{HI}\\ \rho _{HII}\\ \rho _{HeI}\\ \rho _{HeII}\end{array}\right],F_{}=v_{}\left[\begin{array}{c}\rho \\ \rho v_r\\ \rho v_z\\ \rho v_\theta \\ E+P\\ \rho _{H_2}\\ \rho _{HI}\\ \rho _{HII}\\ \rho _{HeI}\\ \rho _{HeII}\end{array}\right],S_{geom}=\left[\begin{array}{c}\rho v_r\\ \rho v_r^2\rho v_\theta ^2\\ \rho v_rv_z\\ 2\rho v_rv_\theta \\ v_r(E+P)\\ \rho _{H_2}v_r\\ \rho _{HI}v_r\\ \rho _{HII}v_r\\ \rho _{HeI}v_r\\ \rho _{HeII}v_r\end{array}\right],$$
$$S_m=\left[\begin{array}{c}0\\ 0\\ 0\\ 0\\ \mathrm{\Lambda }\\ \mu _{H_2}(R_{H_2}D_{H_2})\\ \mu _{HI}(2(D_{H_2}R_{H_2})+D_{HII}R_{HII})\\ \mu _{HII}(R_{HII}D_{HII})\\ \mu _{HeI}(D_{HeII}R_{HeII})\\ \mu _{HeI}(R_{HeII}D_{HeII})\end{array}\right],S_{grav}=\frac{GM}{(r^2+z^2)^{3/2}}\left[\begin{array}{c}0\\ \rho r\\ \rho z\\ 0\\ \rho v_rr+\rho v_zz\\ 0\\ 0\\ 0\\ 0\\ 0\end{array}\right],$$
where $`\rho `$ is the gas density, $`v_r`$, $`v_z`$ and $`v_\theta `$ are the components of the velocity, $`E`$ is the total energy, $`P`$ is the gas pressure and $`\mu _{}`$ is the molecular weight of each species. The code tracks $`H_2`$, $`HI`$, $`HII`$, $`HeI`$, and $`HeII`$ densities separately using the self-consistent multifluid advection method of Plewa & Mรผller (1999). The $`R_{}`$ and $`D_{}`$ terms indicate recombination dissociation rates respectively. The cooling, dissociation and recombination rates included in the microphysical source term, $`S_m`$, are given in appendix A. The local value of the adiabatic index, $`\gamma `$ and the mean molecular weight of the gas is dependent on the local gas composition. In our code these values are taken as piecewise constant across grid cells. We have neglected the effects of $`H_2`$ recombination heating and doubly ionized He. The effect of these processes are small for the conditions considered here.
The code employs an exact hydrodynamic Riemann solver. A spatial and temporal second order accurate wave propagation scheme (Leveque (1997)) is used to advance the solution of the source-free Euler equations. Our code also maintains spatial and temporal second order accuracy in pressure using the method of Balsara (1999). The geometric and micro-physical source terms are handled separately from the hydrodynamic integration using an operator split approach. The source term $`S=S_{geom}+S_m+S_{grav}`$ is integrated using an implicit fourth-order Rosenbrock integration scheme for stiff ODEโs. We note that adaptive mesh refinement has been particularly useful in resolving the neighborhood of thin shock bounded cavity walls that are prevalent in wind blown bubble environments.
### 2.2 Model Parameters and Assumptions
High mass YSO outflows typically exhibit wider opening angles than their low mass counterparts (Kรถnigl, 1999). The wide opening angle subtended by the outflow limbs in the case of source I is indicative of a poorly collimated driving source. Thus our simulations begin with a spherical wind driven into the grid via an inflow boundary condition. This boundary condition is set by reestablishing wind conditions on a โwind sphereโ in the grid before every time step. The wind impinges on a collapsing, rotating molecular envelope. The density distribution in the infalling ambient envelope are prescribed by equations (8), (9) and (10) of the self-gravitating, rotating collapse model of Hartmann, Calvet, & Boss (1996). This infalling envelope is the same as that used in previous outflow models of Delamarter, Frank, & Hartmann (2000) and Gardiner, Frank, & Hartmann (2003) where it was shown that a combination of inertial and ram pressure confinement was sufficient to collimate the wind into a bipolar outflow.
The rotational velocity and infall speed of the envelope is chosen to be appropriate for a $`M=10M_{\mathrm{}}`$ central gravitational source (Ulrich, 1976). Note that previous works Delamarter, Frank, & Hartmann (2000); Gardiner, Frank, & Hartmann (2003) explored the interaction of infall and winds in the context of low mass young stellar objects. The simulation parameters used here are listed in table 1. The wind consists of ionized H and 20% atomic He by mass. The ambient envelope is assumed to be composed of $`H_2`$ and 20% atomic He by mass. We note that even with AMR methods the high cooling rates achieved behind the shocks on the small scales at which these simulations are run ($`50AU`$) present a significant numerical challenge. From the solutions to the rotating collapsing sheet problem we find typical ambient densities of order $`n10^7cm^3`$. Because the cooling rate increases as $`n^2`$, the cooling parameter in the flow, defined as $`\chi =t_{cool}/t_{hydro}`$ is very small with $`\chi <<1`$ for most shock conditions. Thus, in previous works many authors have chosen isothermal equations of state such that $`P\rho ^\gamma `$ with $`\gamma 1`$. Such a description can mimic certain aspects of cooling, such as strong compressions behind shocks, but will not correctly recover the dynamics in more complicated flow patterns. It is more desirable therefore to explicitly track cooling when possible.
Note that in wind blown bubbles three discontinuities will form: a โwind shockโ facing back into the freely expanding wind, an ambient shock facing outward into the envelope and a contact discontinuity between the two shocks delineating the interface between shocked wind and shocked ambient material (figure 2). The strength of cooling determines the distance between the shocks and the contact discontinuity. The principle difficulty in carrying forward simulations such as those described here is achieving adequate resolution to track the flow between the two shocks in the wind-ambient material interaction region. The smallest scales which can be captured with our runs is of the order $`\mathrm{\Delta }x0.2AU`$. Thus cooling scale lengths must be of this order or not significantly less than this if we are to capture the details of the post-shock flow patterns.
In our simulations we have achieved a balance between realism and numerical efficacy by modifying conditions in the wind and ambient medium. For example we have used an initial wind temperature of $`T_w10^4K`$ in the launch region which is likely too high. We use this value as it provides a mach number of $`M20`$ which is useful for launching the simulation. We allow the wind to cool as it expands. More importantly we have reduced the densities in the wind and infalling cloud such that cooling plays a strong role but the interaction regions can be resolved. Thus while outflow rates of $`10^7M_{\mathrm{}}yr^1`$ to $`10^6M_{\mathrm{}}yr^1`$ are representative of higher mass YSO sources (Kรถnigl, 1999) we have chosen an outflow rate of $`1.5\times 10^9M_{\mathrm{}}yr^1`$ to allow our code to adequately resolve strongly cooled shear layers present in the flow. If the wind were launched from an accretion disk this would give a wind mass loss rate that was between $`.1`$ to $`.01`$ lower if standard disk wind theory can be applied. The discrepancy in mass loss rates is clearly a significant difference between our models and the actual situation in source I. This difference should not effect our principle conclusions as these are not sensitive to the details of the cooling.
First we note that in our simulations the flow along the walls of the cavity, at its base, is strongly cooling. The principle difference in this region of the flow between our simulations and one with higher mass loss rates is the width of the wall and distance between the wind shock and the ambient shock. We note that we are interested primarily in the morphology of the outflow base of Source I where the masers define the X as well as the dynamics of rotation in the swept-up material. It has been shown that the morphologies of wind blown bubbles are principally determined by the ratio of specific momenta (or inertia) in the envelope to that in the wind (Icke, 1988). In our case, where we hold the stellar wind velocity and stellar mass constant, it is the infall to outflow mass loss rates ($`f=\dot{M}_i/\dot{M}_w`$) which determine the qualitative details of the flow (Shu, Ruden, Lada, & Lizano, 1991; Delamarter, Frank, & Hartmann, 2000). This is particularly true along the arms of the flow at its base where higher mass loss rates (and stronger cooling) will only narrow the opening angle by a few degrees. We have observed this in our tests in which we have run larger scale simulations with higher mass loss rates as well as cases with an isothermal equation of state. These have shown similar results in terms of the shape of the outflow arms. It is the details of smaller scale flow features, such as those associated with instabilities in the swept up shell, which differ when the interaction region is resolved. With higher mass loss rates the walls of the cavity become so thin that, even with AMR, they span only a few zones and we are not able to resolve their internal dynamics.
## 3 Results
The interaction of a spherical wind expanding into an asymmetric density environment has been well studied both analytically and numerically (Shu, Ruden, Lada, & Lizano, 1991; Icke, 1988; Mellema, Eulderink, & Icke, 1991; Frank & Mellema, 1996). For a review see Frank (1999). In Delamarter, Frank, & Hartmann (2000) the combination of inertial confinement and ram pressure confinement from a collapsing envelope was shown capable of producing a variety of bipolar outflow configurations ranging from well collimated jet-like outflows to wider butterfly shaped outflows with narrow waists. In Delamarter, Frank, & Hartmann (2000) the parameter $`f=\dot{M}_i/\dot{M}_w`$ was found to be critical in determining the shape of the outflow where $`\dot{M}_i`$ and $`\dot{M}_w`$ are the infall and outflow rates respectively. They found that $`f=10`$ models produced fairly wide bubble with opening angle of $`60^{}`$. Delamarter, Frank, & Hartmann (2000) also found that the post-shock cooling produced second order effects on the flow. In particular the spherical wind will impinge upon inward facing โwind shockโ, (which defines the inner walls of the outflow cavity), at an oblique angle. This material will retain much of its initial velocity and will be directed to stream along the contact discontinuity (CD) between the shocked wind and shocked ambient material. Thus the CD becomes a strong slip surface. Note also that the ambient material that passes through the bow shock would be directed to flow toward the equator while the shocked wind will flow toward the head of the outflow. There has been some debate as to the nature of the mixing which occurs in this region and which way the net flow of momentum will travel. Our simulations provide some answer to these questions and these answers are relevant to the nature of the Source I outflow.
Figure 1 shows a density map of the outflow created in the simulation. Globally we see the evacuation of a bipolar cavity where the rotating ambient molecular material is swept toward the perimeter of the cavity walls. The global shape of the outflow is similar to that seen in the studies of Delamarter, Frank, & Hartmann (2000) who used an entirely different code. This gives us confidence that the basic dynamics is being correctly modeled. The shell of swept-up material is thin due to the strong energy losses from molecular dissociation and cooling. The bulk of the cavityโs volume is occupied by freely expanding pre-shocked wind. Along the sides of the cavity we see that once wind material strikes the inward facing wind shock its flow is directed along the CD in a relatively thin shell. At higher latitudes which represent the โheadโ of the outflow, the freely expanding wind encounters the wind shock at a far less oblique angle and is more strongly decelerated. The pressure retained in this region even after cooling pushes the wind shock away from the CD (figure 2). We note that simulations with higher mass infall/outflow rates but calculated on larger scales only show differences at the head of the outflow as the post-wind shock material is able to cool more effectively and moves closer to the CD. We briefly address the dynamics at the head of the outflow below. For now we note that the long term evolution of the outflow (in terms of the cavity walls) relaxes to what appears to be a steady state as inflow and outflow pressures balance. Thus while the outflow head dynamics is of general interest, the only part of the Source I outflow which can be observed through the $`SiO`$ masers on $`<70AU`$ scales will be the arms at the base of the cavity.
The most important aspect of these simulations for the subject at hand is the fact that supersonic wind and ambient rotating molecular material are compressed at the wind and ambient shocks forming a thin layer around the contact discontinuity (figures 1 & 2). This layer then becomes a slip stream surface. The morphology of the outflow is also crucial and is the result of the ambient materialโs density and velocity distributions (i.e. inertial and ram pressure confinement). Figures 3 & 4 show the kinematics of the flow with vectors indicating poloidal flow direction overlaid on a map of velocity magnitude (we defer discussion of rotation until the next section). Here one can see that wind material strikes the inner shock at an oblique angle. Post-shock wind is focused into a flow that streams parallel to the contact discontinuity. This high speed shocked wind moves ahead of the rest of the outflow, forming a cusp at mid to high latitudes that pushes past of the head of the outflow at the poles. These cusps are transient features and will eventually be subsumed into the rest of the outflow. Dense limbs of wind swept ambient material delineate the low latitude edges of the cavity. The dense limbs of the cavity and the narrow waist formed close to the inflow boundary are essential morphological signatures present in both our simulations and the source I maser spot observations.
Velocity vectors of the flow pattern shown in figures 3 and 4 illuminate the formation of a slip stream between initially infalling ambient molecular material and outflowing wind material. The inner line in figure 3 delineates the contact discontinuity. This marks the transition between mostly wind material and mostly ambient material. Note that the change in direction of the velocity vectors as one moves from ambient to wind material. This flow reversal indicates the presence of a vortex across the interaction region. We identify the large density and velocity gradients present across the slip stream as susceptible to Kelvin-Helmholtz instabilities and will discuss these in more detail in the next section.
Thus to conclude this section we find our simulations show that a wind blown bubble with appropriate morphology for Source I will form via the interaction of a spherical wind with a collapsing, rotating envelope. In the next section we discuss both generic outflow issues raised by the simulations as well as specific connections to source I maser observations.
## 4 Discussion
We break our discussion into two sections. First we review our results in light of previous simulations of wind driven molecular outflows examining those features of the simulations which shed new light on unresolved issues. In the second section we compare our simulation results to the observations of Source I focusing particularly on the rotational patterns.
### 4.1 Generic Outflow Issues
Several authors (Shu, Ruden, Lada, & Lizano, 1991; Wilkin & Stahler, 1998; Lee, Stone, Ostriker, & Mundy, 2001) have constructed analytic models of the formation of bipolar molecular outflows. These models invoke the assumption that mixing between the shocked wind and shocked ambient gas occurs instantaneously. Such rapid and local mixing yields an outflow cavity that is delineated by a purely momentum driven thin shell and the swept up wind mass is taken to be negligible. The last assumption implies that the post-shock flow is dominated by shocked ambient gas. Thus gas will be carried along the shell downward toward the disk. Our simulation contradicts this model assumption. While there is much more mass in the ambient material, the momentum in the wind material is significant owing to the high velocity of the wind and the oblique nature of the wind shock. We find that the momentum of the wind material is dynamically important and must be incorporated into any model used to predict the kinematics inside the thin shell. This is true even in cases where turbulence dominates and the cavity walls are fully mixed. Delamarter, Frank, & Hartmann (2000) derive a condition for the mixed flow to be directed upward towards the head of the outflow:
$$\left[\frac{\dot{M}_w}{\dot{M}_i}\left(\frac{\stackrel{~}{M}_{sw}}{\stackrel{~}{M}_{si}}\right)^2\right]>1$$
where $`\stackrel{~}{M}_{sw}`$ is the Mach number of the shocked wind material and $`\stackrel{~}{M}_{si}`$ is the Mach number of the shocked infalling material.
For the parameters used in this simulation, $`\dot{M}_i/\dot{M}_w=10`$, $`\stackrel{~}{M}_{sw}=5`$, $`\stackrel{~}{M}_{si}=.1`$. Thus the left hand side of this equation evaluates to 250 and the condition for upward flow is satisfied. While the cavity walls in our simulations are not fully mixed we can compute the direction the flow would take if mixing occurred by averaging the momentum across the zones which comprise the shell. Figure 5 plots the effective velocity $`v_{eff}`$ that would result along the direction of the outflow cavity wall if the shocked layer were to become fully mixed. Here positive velocity indicates flow streaming upward towards the head of the cavity. Figure 5 shows the flow within the cavity walls to be positive at essentially all latitudes with a velocity of $`v8kms^1`$. This result indicates that if the shocked flow were to mix by any means, the momentum of wind material would overwhelm that of the slower, denser ambient material resulting in poleward movement of material along the shell. Thus we conclude that the assumption that mixing within the bipolar outflow shell leads to a net downward streaming of mass towards the disk is incorrect.
Note that while do not see full mixing, our simulations do capture one net turnover in flow direction. We interpret the vortical reversal of vectors between the ambient shock and the wind shock as due to the development of Kelvin-Helmholtz instabilities and related mixing processes. This mixing results in the entrainment of outwardly directed wind momentum into the shocked ambient gas along the inner cavity walls. While our simulation clearly shows the vortical motions in the interaction region we operate with limited resolution due to computational constraints and cannot resolve the expected multiple non-linear Kelvin-Helmholtz roll-ups in this region. We rely, therefore, on an analytic description to quantify the unresolved instability responsible for entraining momentum from the wind into the region of ambient gas. The characteristic growth rate of the Kelvin-Helmholtz instability for wave mode $`k`$ is given by:
$$\mathrm{\Gamma }_{kh}=|V_{sw}V_{sa}|k\frac{\sqrt{\rho _{sw}\rho _{sa}}}{\rho _{sw}+\rho _{sa}}$$
where $`\rho _{sw}`$ and $`V_{sw}`$ denote the density and velocity on the shocked wind side of the slip stream and $`\rho _{sa}`$ and $`V_{sa}`$ denote the density and velocity on the shocked infalling side of the slip stream. We have computed the time scale for the growth Kelvin-Helmholtz modes $`t_{kh}`$ of wavelength $`\lambda `$, relative to the dynamical time scale of the outflow source as $`1/\mathrm{\Gamma }_{kh}=t_{kh}0.05t_{dyn}\lambda /r_w`$ where $`r_w`$ is the distance from the central source to the wind shock. The analysis used here is for the Kelvin-Helmholtz growth across a thin shear layer. Since wavelengths less than the thickness of the slip stream layer are damped, we choose the wavelength of fastest growth equal to the width of the cavity wall or shell $`\lambda 0.25r_w`$. Using values of the parameters taken from the simulations we find timescale for the growth of this mode is much less than the dynamical timescale of the outflow (i.e. $`t_{kh}<<t_{dyn}=7.1years`$). This implies that unresolved Kelvin-Helmholtz instabilities will provide the mixing required to exchange momentum across the shear layer even though in our simulations we see a smooth (numerically) diffusive mixing in the shear layer. This result supports our interpretation that mixing has resulted in the entrainment of outwardly directed wind material into the layer of shocked ambient gas. With higher resolution we would expect to see the shear layer become more turbulent.
We now address the fragmentation of the head of the outflow in our simulations. While details of the polar caps of the outflow are not illuminated by the Source I $`Si0`$ maser observations, the fragmentation of the thin shell of swept-up ambient material in this region warrants some discussion. The presence of these fragments dominates the dynamics of the flow in the polar caps. We note that the dynamical age of the outflow in our model is very young; 7.1 years. A typical astrophysical outflow of much greater dynamical age would have time to produce even richer fragmentation and clumping spectrum than our the flow modelled in our simulation. The instability along the polar caps of the outflow and the formation of associated fragments and inhomogeneities is interesting in that it is characteristic of YSO molecular outflows (McCaughrean & Mac Low, 1997). The nature of โclumpyโ flows remains largely unexplored though Poludnenko, Frank, & Blackman (2002); Meliolo et al. (2005) have examined properties of inhomogeneities on the propagation of shocks and the role such phenomena play in astrophysical contexts. Our simulations imply that such clumpy flows are likely to be a generic feature of bipolar outflows and jets.
We can, perhaps, understand the fragmentation seen in our simulations through known unstable modes of shocks. Vishniac (1983) has examined the stability of a thin spherical shock of wavelengths $`>>`$ than the thickness of the thin shell. They find that shocks that are sufficiently radiative to produce a density contrast of $`10`$ are dynamically unstable. The growth rate associated with these modes are given by $`\mathrm{\Gamma }_{ts}c_s/h`$ where $`c_s`$ is the sound speed inside the thin shell and $`h`$ is the shell thickness. The characteristic temperature $`10^4K`$ and shell thickness of $`2AU=0.1r_w`$ along the thin dense shell at the rapidly expanding polar caps in the simulated outflow yield a timescale for the growth rate of this instability of $`1/\mathrm{\Gamma }_{ts}=t_{ts}0.04t_{dyn}`$. Thus the rapid growth of the fragments which develop in this region are consistent with Vishniacโs analytic prediction.
### 4.2 Comparison with Observations
In this section we make a qualitative comparison between the broad morphological and kinematic properties of our simulations and those observed in the $`SiO`$ maser spots of Source I. We note that more detailed quantitative comparisons will require high resolution 3-D simulations including radiation-transfer calculations beyond the scope of the current work.
First we note our previous works have also shown that the opening angle of the outflow depends on the degree of flattening of the ambient material (Delamarter, Frank, & Hartmann (2000). Flatter, more pancake-like density (and infall velocity) distributions lead to more spherical bubbles with a larger opening angle for the arms of X. Thus the morphology of the models can be smoothly adjusted to fit the conditions in Source I. Most importantly the wide opening angle seen in the Source I masers argues strongly for the presence of a significant wide-angle wind component to the driving wind. While a jet component may still exist the laterial expansion of the lobes at the base would be difficult to reproduce without a wind component with significant momenta expanding into low latitudes.
We have proposed that the line of sight maser velocities observed about Source I are due to rotation. Specifically it is rotation retained by infalling molecular material that has been intercepted by a biconical outflow. We now demonstrate this assertion. Figures 6 & 7 show a color map of the rotational component of velocity $`v_\theta `$ in the shell of swept up material. The non-rotating wind material is delineated by the interior line. Notice that the region of greatest shear occurs midway between the outer edge of the ambient shock and the slip surface (CD). This is where the flow undergoes direction reversal across the slip stream. As discussed above we can identify this region as adjacent to the line of greatest vorticity generation. This region is most susceptible Kelvin-Helmholtz instabilities and the creation of related density inhomogeneities. The formation of fragments and filaments in this region would provide the density and path length enhancement most likely to result in the amplification of master spots. Because the masing material most likely lies in this region, its kinematics would be likely be revealed by observations of maser spots.
Figure 8 plots a cut of rotational velocity taken along this surface midway between the ambient shock and the contact surface. The magnitude of the rotational velocities is $`3kms^1<v_{rot}<15kms^1`$ in rough agreement with the line of sight velocities seen in the maser data. These values reflect the rotational velocities of the pre-shock ambient material. Thus as ambient material spirals inward towards the star it intercepts the cavity wall, is shocked and eventually reverses its poloidal but not its toroidal velocity. Our results thus indicate that the observed line of sight velocity of the maser spots can be interpreted as the rotation of infalling material that has been intercepted by a poorly collimated biconical outflow. Note that we correct for a systematic $`5kms^1`$ redshift as is appropriate for the Orion region.
Note also that our results show a decrease in rotational velocity as one moves outward along the arms of the X defined by the maser spots. The presence of a velocity gradient is ambigious in the observations. For example, in the data of Greenhill et al. (1998), the redshifted lobes show a velocity gradient with values ranging from $`25kms^1`$ at the base of the X to $`15kms^1`$ at its farthest extent. The blueshifted lobes do not show such a clear gradient however they do show a similar range of velocities. In our model gradients in $`v_{rot}`$ in the maser spots will reflect, to first order, gradients in the rotational velocities of the ambient material. A better accounting of turbulent advection of material along the walls of the cavity could smooth out the observed rotational gradient. However since no torques act on the swept up material conservation of angular momentum will still reduce its rotational velocity. Consideration of the initial conditions shows that the gradient in $`v_{rot}`$ is stronger in the radial direction than in height. Thus for models with wider opening angles the we expect the gradient to be less dramatic.
Finally one may ask if the rotation inferred for the Source I masers is due to the ambient material as we have modelled or if it comes from rotation of the stellar wind. We note that the observed maser proper motions of $`15kms^1`$, which we take to be poloidal velocities along the shell, are an order of magnitude smaller than both typical young stellar object outflows and stellar winds. If this material were to originate from a disk wind it would have to emerge from a region quite far from the central source. Using $`v_{escape}v_{Keplerian}(R)`$ gives $`R40AU`$ which is quite far from the inner regions of the disk where launching is expected to be most effective (Ouyed et al., 2003). One might argue that the modest rotational gradient observed in the maser data indicates the presence of a magnetic field in the wind anchored to the disk such as would be the case for a disk wind. The field would then sustain the rotational motions of the wind material caught-up in the masers. However once any MHD launching mechanism takes a wind parcel out beyond the Alfvรฉn radius the field will no longer provide rotational support for the wind and angular momentum conservation will reduce the windโs rotational velocity as it expands. Since the Alfvรฉn radius tends to be a few times larger then the radius of the footpoint of the flow, and the flow is likely to form close to the star at the inner regions of the disk, it is unlikely that the winds rotational motion can be magnetically supported. Thus the rotation seen in the masers are most convincingly interpreted as ambient material that has been swept up or entrained along shocks with the fast-moving wind (Greenhill et al., 2003; Doeleman, Lonsdale, & Pelkey, 1999). In all models, the $`H_2`$ densities and temperatures necessary for maser action, and the small line-of-sight velocity shifts necessary for maser amplification occurs along the limbs of the outflow.
## 5 Conclusion
Using AMR simulations we have explored the evolution of an outflow driven by a spherical wind from a massive gravitating source interacting with a rotating infalling cavity. The dynamics of the flow in the shocked ambient material are essential aspects of our model that have been neglected in previous works. Specifically, the growth of Kelvin-Helmholtz unstable modes in the slip stream between outflowing wind material and infalling molecular material along the cavity walls provides mixing at the wind/ambient gas interface. This mixing results in the outward acceleration and eventual direction reversal of initially infalling ambient material.
We have also shown that the head of the outflow will be unstable to thin shell instabilities. While this conclusion is not relevant to the observations of Source I it is of general interest for studies of molecular outflows. The fragmentation of radiatively cooling outflow lobes has important consequences for their long term dynamics and observational properties. The fact that our simulations show a rapid transition to fragmentation implies that molecular outflows on larger scales can be expected to be โclumpyโ on a variety of scales. This issue should be addressed in future works.
We have compared our simulations with observations of maser spots in Source I in the BN/KL region in Orion. We find that when the wind evacuates a bipolar cavity, swept-up rotating ambient molecular material in the cavity walls is the likely source of the masers. This zone contains material which retains its rotation but which has been accelerated upward towards the head of the outflow lobe. The poloidal velocities seen in our simulations are of order those seen proper motions in the observations. The rotation retained by the shocked ambient material is consistent with the observed line of sight velocity of the maser spots observed about Source I. Thus we conclude that the line of sight motions in Source I inferred to be due to rotation (Greenhill et al., 2003, 2005) can be interpreted as originating in the rotation of the collapsing ambient material. Future work will need to explore the flow pattern in 3-D as well address issues related to formation of the masers in greater detail. We also conclude that the wide opening angle of the maser spot pattern is strong evidence that a significant wide angle wind component is at work in the Source I outflow.
We thank Lincoln Greenhill and Mark Reid for fruitful conversations which helped this work. Support for this work was provided by NFS grant AST 00-98442, NASA grant NAG5-8428, an HST grant, DOE grant DE-FG02-00ER54600, and the Laboratory for Laser Energetics.
## Appendix A $`H_2`$ Microphysics
### A.1 Cooling Functions
We employ a total cooling function $`\mathrm{\Lambda }=\mathrm{\Lambda }_{LS}+\mathrm{\Lambda }_{DM}+\mathrm{\Lambda }_{OI}+\mathrm{\Lambda }_{HeI_D}+\mathrm{\Lambda }_{HI_D}+\mathrm{\Lambda }_{H_{2D}}`$. Where $`\mathrm{\Lambda }_{DM}=n_{H_{ion}}n_e^{}\lambda _{DM}`$ and $`\lambda _{DM}`$ is the atomic line cooling function appropriate for interstellar gases by Dalgarno & McCray (1972), the dominant cooling process at $`T>10^4`$. $`\mathrm{\Lambda }_{LS}`$ is the $`H_2`$ cooling function of Lepp & Schull (1983) appropriate for a molecular gas. $`\mathrm{\Lambda }_{OI}=n_{HI}n_{OI}\lambda _{OI}(T)`$ dominates the cooling at $`T<300K`$ where $`\lambda _{OI}(T)`$ is the OI line cooling function tabulated by Launay & Roueff (1977). For temperatures greater than those listed in the table, $`\lambda _{OI}(T)`$ extrapolates with $`\sqrt{T}`$. Because the first ionization potential of $`O`$ is comparable to that of $`H`$, we use the abundance of molecular and atomic $`H`$ as a tracer for OI given as $`n_{OI}=(n_{H_2}+nHI)/n_{H_{ions}}f_{OI}`$ where $`f_{OI}=8.51\times 10^4`$ is the fractional coronal abundance of oxygen. The $`\mathrm{\Lambda }_{HeI_D}`$, $`\mathrm{\Lambda }_{HI_D}`$, and $`\mathrm{\Lambda }_{H_{2D}}`$ terms account for cooling due to dissociation and ionization.
$`\mathrm{\Lambda }_{HeI_D}`$ $`=`$ $`24.60eVD_{HeI}`$
$`\mathrm{\Lambda }_{HI_D}`$ $`=`$ $`13.59eVD_{HI}`$
$`\mathrm{\Lambda }_{H_{2D}}`$ $`=`$ $`4.48eVD_{H2}`$
Where D represents the dissociation rates given in the following sections. Figure 9 shows each of these cooling rates for typical ISM density.
### A.2 Dissociation
Lepp & Schull (1983) fit an analytic function for the $`H_2`$ dissociation rate which takes the form:
$$\mathrm{log}_{10}(k_D^{H_2,H})=\frac{\mathrm{log}_{10}k_H\mathrm{log}_{10}(k_H/k_L)}{1+n/n_{cr}}$$
where, for $`HH_2`$ collisions, $`\mathrm{log}_{10}(n_{cr})=4.000.416x0.327x^2`$ and, for $`H_2H_2`$ collisions, $`\mathrm{log}_{10}(n_{cr})=4.130.968x+0.119x^2`$ with $`x=\mathrm{log}_{10}(T/10^4K)`$. $`k_H`$ and $`k_L`$ refer to the high ($`n>>n_{cr}`$) and low ($`n<<n_{cr}`$) density limits for the reaction rate. Lepp & Schull (1983) give:
$$k_H(T)=\{\begin{array}{cc}3.52\times 10^9\mathrm{exp}(4.39\times 10^4/T)\text{for}HH_2\hfill & \\ 5.48\times 10^9\mathrm{exp}(5.30\times 10^4/T)\text{for}H_2H_2\hfill & \end{array}$$
Lim et. al. (2002) Improve the dissociation low density dissociation rate of Lepp & Schull (1983) by using the more recent calculations for $`k_L^{HH_2}`$ of Dove & Mandy (1986). Lim et. al. (2002) fit the results of Dove & Mandy (1986) to the form:
$$k_L(T)=4.69\times 10^{14}T^{0.746}\mathrm{exp}(5.55065\times 10^4/T)\text{for}HH_2$$
We modify the original rates further by using the more recent calculations for low density $`H_2H_2`$ dissociation of Martin et. al. (1998). We have also included corrections to the dissociation rate through the action of $`HeH_2`$ and $`e^{}H_2`$ collisions of Martin et. al. (1998) given by:
$$k_D^i(T)=\left(\frac{8dkT}{\pi \mu }\right)^{1/2}\frac{a(kT)^{b1}\mathrm{\Gamma }(b+1)\mathrm{exp}(Eo/kT)}{\left(1+CkT\right)^{b+1}}\text{for}H_2H_2$$
where $`d=1.894\times 10^{22}`$ $`k=3.167\times 10^6`$ and the constants a, b, c and Eo for each collision partner are given in table2 and $`\mu `$ is the reduced mass of the collision pair.
The total dissociation rate is given by:
$$D_{H_2}=\underset{i}{}n_in_{H_2}k_D^i$$
where i ranges over all collision partners: $`H_2`$,$`H`$,$`He`$,and $`e^{}`$.
### A.3 Molecular Recombination Rate
Hollenbach & McKee (1979) give an approximation to the $`H_2`$ recombination rate due to HI โstickingโ on dust surfaces as:
$`R_{H_2}`$ $`=`$ $`n_{H_{nuclei}}n_{HI}\mathrm{\hspace{0.17em}3}\times 10^{17}cm^3s^1{\displaystyle \frac{T_2^{1/2}f_a}{1+0.4(T_2+T_{dust\mathrm{\hspace{0.17em}2}})^{1/2}+0.2T_2+0.08T_2^2}}`$
$`T_2`$ $`=`$ $`T/100`$
$`f_a`$ $`=`$ $`0.5=\text{ fraction of molecules that do not evaporate on dust surfaces.}`$
### A.4 Shock Dissociation
As a test of these dissociation rates as implemented in our code we have computed the impulsively launched steady shock speed required for molecular dissociation. Figure 10 shows pre-shock density vs. shock speed for a steady shock resulting in 90% downstream $`H_2`$ dissociation. This result is consistent within a few $`kms^1`$ with the results of previous authors (Smith, 1994; Hollenbach & McKee, 1977).
### A.5 Ionization & Recombination
Mazzotta et. al. (1998) catalog ionization and recombination rates for many species. We have employed HeI and HI ionization rates originally from Arnaud & Rothenflug (1985). Verner & Ferland (1996) fit the radiative recombination rate coefficients for several species including the He and H recombinations rates used for this work. We have included the fit to dielectric contribution to the He recombination rate of Mazzotta et. al. (1998). |
warning/0506/hep-ph0506050.html | ar5iv | text | # Thermodynamics with density and temperature dependent particle masses and properties of bulk strange quark matter and strangelets
## I Introduction
One of the most exciting possibility from QCD is that hadronic matter undergoes a rich and varied phase landscape with increasing densities. At extremely high densities, the mass of strange quarks becomes unimportant and all the three flavors of $`u`$, $`d`$, and $`s`$ quarks can be treated on an equal footing. Consequently, quark matter is in the color-flavor-locked (CFL) phase Rajagopal2001PRL86 and/or a new gapless CFL phase (gCFL) Alford2004PRL92 . However, if densities are not that high, the strong interactions between quarks become important, and quark matter is in the unpaired phase. Presently, RHIC is teaching us about the properties of the hot but not asymptotically hot quark gluon plasma Gyulassy2004 ; Shuryak2004 . Actually the future FAIR project which will be built in the coming years at GSI in Germany is targeted towards the physics of ultradense matter like it is found in neutron stars FAIRweb .
The original idea of Witten is that strange quark matter (SQM), rather than the normal nuclear matter, might be the true ground state of the strong interaction Witten1984PRD30 . Immediately after Wittenโs conjecture, Farhi and Jaffe showed that SQM is absolutely stable near the normal nuclear saturation density for a wide range of parameters Farhi1984PRD30 . Now SQM has been investigated for more than two decades since the pioneer works of many authors Bodmer1971PRD4 ; Chin1979PRL43 ; Witten1984PRD30 ; Farhi1984PRD30 . Because the lattice gauge theory still has difficulty in the consistent implementation of chemical potential in numerical simulations presently Creutz2001NPB94c while the perturbative approach is unreliable at the strong coupling regime, phenomenological models reflecting the characteristics of QCD are widely used in the study of hadrons, and many of them have been successfully applied to investigating the stability and properties of SQM.
An important question in the study is how to incorporate pressure balance. Basically, one uses a model to give the thermodynamic potential, then add to it a constant to get mechanical equilibrium. This is the famous โbagโ mechanism. Many important investigations have been done in this direction Berger1987PRC35 ; Gilson1993PRL71 ; Madsen1993PRL70 ; Madsen2000PRL85 ; SchaffnerBielich1997PRD55 ; Parija1993PRC48 ; HeYB53PRC1903 ; Madsen87PRL172003 ; Schafer2002NPA ; Ratti2003 ; Wangp2003PRC67 ; Harris2004JPG30 . A common feature of these investigations is that quark masses are constant, so the normal thermodynamic formulas can be used without thermodynamic inconsistency problems. Actually, however, it is well known in nuclear physics that particle masses vary with environment, i.e., the density and temperature. Such masses are usually called effective masses Walecka1995OSNP16 ; Henley1990NPA617 ; Brown1991PRL66 ; Cohen1991PRL67 . Effective masses and effective bag constants for quark matter have been extensively discussed, e.g., within the NambuโJona-Lasinio model Buballa457PLB261 and within a quasiparticle model Schertler1997NPA616 . In principle, not only masses will change in the medium but also the coupling constant will run in the medium Fraga2004hepph0412298 . The question now is how to treat the thermodynamic formulas which do not violate the fundamental principles of thermodynamics, when introducing density and/or temperature dependent masses. In fact, a lot of problems have been caused in this field.
There exist in literature several kinds of thermodynamic treatments. The first one uses all the thermodynamic formulas exactly the same as in the constant-mass case Chakrabarty1989PLB229 ; Cleymans1994ZPC62 ; Letessier1994PLB323 . The second treatment adds a new term, originated from density dependence of quark masses, to both pressure and energy Benvenuto1989PRD51 . These two treatments were later proved to be inconsistent with the necessary thermodynamic requirement: their free energy minimum do not correspond to zero pressure Peng2000PRC62 . The third treatment adds the term from the density dependence of quark masses to the pressure but not to the energy, and the inconsistency disappear Peng2000PRC62 . Another treatment is the addition of a new term to the thermodynamic potential. This has been done when masses depend on either temperature Gorenstein1995PRD52 or chemical potential/density Schertler1997NPA616 ; Wangp2000PRC62 . However, if particle masses depend on both chemical/density and temperature, this way meets difficulties. We will discuss it further in the next section.
With the third thermodynamic treatment and the cubic root scaling Peng2000PRC61 , Zheng et al. have studied the viscosity of SQM and calculated the damping time scale due to the coupling of the viscosity and $`r`$ mode Zhengxp2004PRC70 . This model has also been applied to investigating the quark-diquark equation of state and compact star structure Lugonese2003IJMPD12 .
Another important progress has been made recently by Zhang et al. Zhangy2002PRC65 ; Zhangy2001EPL56 ; Zhangy2003PRC67 ; Zhangy2003MPLA18 . They extended the quark mass density dependent model to finite temperature to let the model be able to describe phase transition. In this case, the quark masses depend on both density and temperature and the permanent confinement is removed. They use the third thermodynamic treatment mentioned above Peng2000PRC62 and parametrize the common interacting part of quark masses as Zhangy2002PRC65 $`m_\mathrm{I}=B_0(3n_\mathrm{b})^1[1\left(T/T_c\right)^2]`$ with $`T_\mathrm{c}`$ being the critical temperature. Later, they found this parametrization causes an unreasonable result: the radius of strangelets decreases with increasing temperature Zhangy2001EPL56 . So they added a new linear term, and the parametrization became Zhangy2001EPL56 ; Zhangy2003PRC67
$$m_\mathrm{I}=\frac{B_0}{3n_\mathrm{B}}\left[1a\frac{T}{T_c}+b\left(\frac{T}{T_c}\right)^2\right].$$
(1)
This extension of the model has soon been applied to the study of strangelets Zhangy2001EPL56 ; Zhangy2003PRC67 ; Zhangy2003MPLA18 , dibaryons Zhangy2004JPG30 , and proto strange stars Gupta2004IJMPD12 , etc..
The purpose of the present paper is two-folded. First, we would like to point out that the thermodynamic derivation in Ref. Peng2000PRC62 is mainly concentrated on density dependent masses. When masses are also temperature dependent, there are special issues to be considered. We will prove that the temperature dependence of quark masses causes another term which contributes to the entropy and energy. With the new term, originated from the temperature dependence of masses, added to the entropy, the quark mass scaling Eq. (1) has a serious problem, i.e. $`lim_{T0}m_\mathrm{I}/T0`$ which leads to a non-zero entropy at zero temperature, violating the third law of thermodynamics. Therefore, our second purpose is to derive or suggest a new quark mass scaling, and then study the properties of bulk SQM and strangelets. An interesting new observation is that low mass strangelets near $`\beta `$ equilibrium are multiquark states with an anti-strange quark, such as the pentaquark $`(u^2d^2\overline{s})`$ for baryon number 1 and octaquark $`(u^4d^3\overline{s})`$ for dibaryon etc.
The paper is organized as follows. In Sec. II, we derive, in detail, the thermodynamic formulas at finite temperature with density and/or temperature dependent masses, and show why, in the present version of the quark mass density and temperature dependent model, one should require
$$\underset{T0}{lim}\frac{m_\mathrm{q}}{T}=0.$$
(2)
In the subsequent Sec. III, a new quark mass scaling at finite temperature is derived or suggested based on chiral and string model arguments, and accordingly in Sec. IV and Sec. V, respectively, the thermodynamic properties of bulk SQM and strangelets at both zero and finite temperature are calculated with the new thermodynamic formulas and the new quark mass scaling. Finally, a summary is given in Sec. VI.
## II Thermodynamics with density and temperature dependent masses
Quasi particle models have been explored in great detail over the past 10-15 years Peshier2002PRD66 . It is well understood how to construct thermodynamically consistent models when the masses depend on the chemical potential and temperature. In the present model, however, there are three differences. First, the particle masses depend on density and temperature, not directly on chemical potential and temperature. Secondly, the density and temperature dependence is independently determined from other arguments. And thirdly, the finite size effects have to be included. Therefore, we derive, in the following, the thermodynamic formulas suitable for the present case.
Suppose the thermodynamic potential is known as a function of the temperature $`T`$, volume $`V`$, chemical potentials $`\{\mu _i\}`$, and particle masses $`\{m_i\}`$, i.e.,
$$\overline{\mathrm{\Omega }}=\overline{\mathrm{\Omega }}(T,\{\mu _i\},\{m_i\},V).$$
(3)
We here explicitly write out the arguments to make the meaning of partial derivatives clear in the following. If the masses $`m_i`$ are constant, other thermodynamic quantities can be obtained from normal formulas available in textbooks. Here the quark masses are density and/or temperature dependent, i.e.,
$$m_i=m_i(n_\mathrm{b},T),$$
(4)
where the baryon number density $`n_\mathrm{b}`$ is connected to the particle numbers $`\{\overline{N}_i\}`$ and volume $`V`$ by
$$n_\mathrm{b}=\underset{i}{}n_i/3\text{with}n_i\underset{i}{}\overline{N}_i/V.$$
(5)
To study this question, we start from the fundamental derivation equality of thermodynamics, i.e.,
$$d\overline{E}=Td\overline{S}PdV+\underset{i}{}\mu _id\overline{N}_i.$$
(6)
This is nothing but the combination of the first and second laws of thermodynamics. It means that the energy $`\overline{E}`$ is the characteristic function, i.e., all other quantities can be obtained from it, if one takes the entropy $`\overline{S}`$, the volume $`V`$, and the particle numbers $`\{\overline{N}_i\}`$ as the full independent state variables. But it is sometimes convenient to take $`(T,V,\{\mu _i\})`$ as the full independent state variables. In this case, the characteristic function is the thermodynamic potential $`\overline{\mathrm{\Omega }}`$ which is defined to be
$$\overline{\mathrm{\Omega }}\overline{E}T\overline{S}\underset{i}{}\mu _i\overline{N}_i$$
(7)
because adding $`d(T\overline{S}+_i\mu _iN_i)`$ to both side of Eq. (6) will give
$$d\overline{\mathrm{\Omega }}=\overline{S}dTPdV\underset{i}{}\overline{N}_id\mu _i.$$
(8)
Another important characteristic function is the free energy $`\overline{F}`$. It is defined to be
$$\overline{F}\overline{E}T\overline{S}.$$
(9)
Then the corresponding basic derivation equation is
$$d\overline{F}=\overline{S}dTPdV+\underset{i}{}\mu _id\overline{N}_i$$
(10)
which can be obtained by adding $`d(T\overline{S})`$ to Eq. (6). Therefore, the independent state variables are $`(T,V,\{\overline{N}\})`$ in this case, i.e.,
$$\overline{F}=\overline{F}(T,V,\{N_i\}).$$
(11)
According to the second term on the right hand side of Eq. (10), one has a general expression for the pressure
$$P=\frac{dF}{dV}|_{T,\{N_i\}}.$$
(12)
Here $`\overline{F}`$ should be expressed as a function of $`(T,V,\{\overline{N}_i\})`$, and the derivative is taken with respect to the volume at fixed $`T`$ and $`\{\overline{N}_i\}`$. Comparing Eqs. (7) and (9) leads to the basic relation between thermodynamics and statistics, i.e.,
$$\overline{F}=\overline{\mathrm{\Omega }}+\underset{i}{}\mu _i\overline{N}_i$$
(13)
Substituting this into Eq. (12) gives
$$P=\frac{d\overline{\mathrm{\Omega }}}{dV}|_{T,\{N_i\}}\underset{i}{}\overline{N}_i\frac{\mu _i}{V}.$$
(14)
Because the independent state variables here are $`(T,V,\{N_i\})`$, the chemicals $`\{\mu _i\}`$ in $`\overline{\mathrm{\Omega }}`$ \[see Eq. (3)\] should be expressed as a function of $`(T,V,\{N_i\})`$, i.e.,
$$\mu _i=\mu _i(T,V,\{N_k\}).$$
(15)
So the total derivative of $`\overline{\mathrm{\Omega }}`$ with respect to $`V`$ at fixed $`T`$ and $`\{N_i\}`$ is
$$\frac{d\overline{\mathrm{\Omega }}}{dV}|_{T,\{N_k\}}=\frac{\overline{\mathrm{\Omega }}}{V}+\underset{i}{}\frac{\overline{\mathrm{\Omega }}}{\mu _i}\frac{\mu _i}{V}+\underset{i}{}\frac{\overline{\mathrm{\Omega }}}{m_i}\frac{m_i}{V},$$
(16)
where
$$\frac{m_i}{V}\frac{m_i}{V}|_{T,\{N_k\}}=\frac{m_i}{n_\mathrm{b}}\frac{n_\mathrm{b}}{V}=\frac{n_\mathrm{b}}{V}\frac{m_i}{n_\mathrm{b}}.$$
(17)
Consequently, substitution of Eq. (16) into Eq. (12) gives
$$P=\frac{\overline{\mathrm{\Omega }}}{V}+\frac{n_\mathrm{b}}{V}\underset{i}{}\frac{\overline{\mathrm{\Omega }}}{m_i}\frac{m_i}{n_\mathrm{b}}\underset{i}{}\left(\overline{N}_i+\frac{\overline{\mathrm{\Omega }}}{\mu _i}\right)\frac{\mu _i}{V}.$$
(18)
Similarly, for the entropy, we have
$$\overline{S}=\frac{d\overline{F}}{dT}|_{V,\{N_k\}}=\frac{d\overline{\mathrm{\Omega }}}{dT}|_{V,\{N_k\}}\underset{i}{}\overline{N}_i\frac{\mu _i}{T}.$$
(19)
Substitution of
$$\frac{d\overline{\mathrm{\Omega }}}{dT}|_{V,\{N_k\}}=\frac{\overline{\mathrm{\Omega }}}{T}+\underset{i}{}\frac{\overline{\mathrm{\Omega }}}{\mu _i}\frac{\mu _i}{T}+\underset{i}{}\frac{\overline{\mathrm{\Omega }}}{m_i}\frac{m_i}{T}$$
(20)
leads to
$$\overline{S}=\frac{\overline{\mathrm{\Omega }}}{T}\underset{i}{}\frac{\overline{\mathrm{\Omega }}}{m_i}\frac{m_i}{T}\underset{i}{}\left(\overline{N}_i+\frac{\overline{\mathrm{\Omega }}}{\mu _i}\right)\frac{\mu _i}{T}.$$
(21)
The energy can be obtained from Eq. (7)
$$\overline{E}=\overline{\mathrm{\Omega }}+\underset{i}{}\mu _i\overline{N}_i+T\overline{S}.$$
(22)
Replacing $`\overline{S}`$ here with the expression in Eq. (21), one has
$`\overline{E}`$ $`=`$ $`\overline{\mathrm{\Omega }}+{\displaystyle \underset{i}{}}\mu _i\overline{N}_iT{\displaystyle \frac{\overline{\mathrm{\Omega }}}{T}}T{\displaystyle \underset{i}{}}{\displaystyle \frac{\overline{\mathrm{\Omega }}}{m_i}}{\displaystyle \frac{m_i}{T}}`$ (23)
$`T{\displaystyle \underset{i}{}}\left(\overline{N}_i+{\displaystyle \frac{\overline{\mathrm{\Omega }}}{\mu _i}}\right){\displaystyle \frac{\mu _i}{T}}.`$
Now we need $`N_\mathrm{f}`$ (number of flavors) equations to connect $`T`$, $`V`$, $`\{\mu _i\}`$, and $`\overline{N}_i`$, so that the functions $`\mu _i(T,V,\{N_k\})`$ in Eq. (15) can be obtained. Presently, nearly all relevant models adopt Chakrabarty1989PLB229 ; Benvenuto1989PRD51 ; Peng2000PRC62 ; Gorenstein1995PRD52 ; Schertler1997NPA616 ; Wangp2000PRC62 ; Zhangy2002PRC65 ; Zhangy2003MPLA18 ; Zhangy2003PRC67 ; Zhangy2004JPG30 ; Gupta2004IJMPD12
$$\overline{N}_i=\frac{\overline{\mathrm{\Omega }}}{\mu _i}|_{T,V,\{m_k\}}.$$
(24)
Please note, $`\overline{N}_i`$ also appears on the right hand side of this equation through $`m_k=m_k(_i\overline{N}_i/[3V],T)`$. Then $`\mu _i`$ can be solved as a function of $`T`$, $`V`$, and $`\{N_k\}`$ from these equations.
If Eq. (24) applies, the last term in Eqs. (18), (21), and (23) vanishes. So all the formulas will take the simplest form. If define $`\mathrm{\Omega }\overline{\mathrm{\Omega }}/V`$, $`E\overline{E}/V`$, $`S\overline{S}/V`$, $`n_i\overline{N}_i/V`$, and use $`V=(4/3)\pi R^3`$, then Eqs. (18), (23), (21), (13), and Eq. (24) become, respectively,
$`P`$ $`=`$ $`\mathrm{\Omega }V{\displaystyle \frac{\mathrm{\Omega }}{V}}+n_\mathrm{b}{\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{\Omega }}{m_i}}{\displaystyle \frac{m_i}{n_\mathrm{b}}},`$ (25)
$`E`$ $`=`$ $`\mathrm{\Omega }{\displaystyle \underset{i}{}}\mu _i{\displaystyle \frac{\mathrm{\Omega }}{\mu _i}}T{\displaystyle \frac{\mathrm{\Omega }}{T}}T{\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{\Omega }}{m_i}}{\displaystyle \frac{m_i}{T}},`$ (26)
$`S`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{T}}{\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{\Omega }}{m_i}}{\displaystyle \frac{m_i}{T}},`$ (27)
$`F`$ $`=`$ $`\mathrm{\Omega }{\displaystyle \underset{i}{}}\mu _i{\displaystyle \frac{\mathrm{\Omega }}{\mu _i}},`$ (28)
$`n_i`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{\mu _i}}.`$ (29)
Compared with thermodynamic formulas in the normal case, both the entropy and energy have a new term from the temperature dependence of the masses, while the pressure has a new term due to the density dependence of the masses. The second term in Eq. (25) exists when the finite size effect can not be ignored, no matter the masses are constant or not. These new terms are understandable with a view to the equalities
$$S=\frac{d\overline{\mathrm{\Omega }}}{dT}|_{V,\mu },P=\frac{d\overline{\mathrm{\Omega }}}{dV}|_{T,\mu },$$
(30)
and the total derivative rules in mathematics. Here one should pay special attention to the difference between the total derivatives and the partial derivatives in Eq. (18) and Eq. (21).
ยฟFrom the viewpoint of quasiparticle models, the first and second terms in the pressure Eq. (25) are the normal quasiparticle contributions, while the last extra term is the contribution of mean-field interactions.
Normally, the first term on the right hand side of Eq. (27) goes to zero when the temperature approaches to zero, i.e., $`lim_{T0}\mathrm{\Omega }/T=0`$. However, the first factor of the second term does not, i.e., $`lim_{T0}\mathrm{\Omega }/m_i0`$. Therefore, one must require
$$\underset{T0}{lim}\frac{m_i}{T}=0$$
(31)
to be consistent with the third law of thermodynamics, i.e., $`lim_{T0}S=0`$. Because Eq. (27) depends on the model assumption Eq. (24), Eq. (31) is a model dependent requirement. Consequently, the thermodynamic formulas Eqs. (25)-(29) are merely valid for systems whose interactions meet this requirement. In the subsequent section, we will see that Eq. (31) is indeed fulfilled by a new quark mass scaling for the strong interactions of quarks.
In Ref. Peng2000PRC62 , we mentioned a necessary condition any consistent thermodynamic treatment must satisfy. At finite temperature, we also have a similar criterion:
$$P=\frac{d\overline{F}}{dn_\mathrm{b}}|_{T,\{\overline{N}_k\}}\frac{n_\mathrm{b}}{V}=n_\mathrm{b}^2\frac{d}{dn_\mathrm{b}}\left(\frac{F}{n_\mathrm{b}}\right)_{T,\{\overline{N}_k\}}.$$
(32)
In obtaining the second equality of Eq. (32), the equalities $`V=_i\overline{N}_i/(3n_\mathrm{b})`$, $`n_\mathrm{b}/V=n_\mathrm{b}/V`$, and $`\overline{F}=VF`$ have been used. Because $`F/n_\mathrm{b}=\overline{F}/(_i\overline{N}_i/3)`$ is the free energy per baryon, Eq. (32) shows explicitly that the free energy extreme occurs exactly at zero pressure. This is why people are interested in the free energy minimum, rather than the energy minimum, at finite temperature, to look for mechanically stable states.
The extra term or the last term in Eq. (25) is important even in the MIT bag model when one introduces a density dependent bag constant $`B(n_\mathrm{b})`$. In this context, the extra term is $`n_\mathrm{b}dB/dn_\mathrm{b}`$. If it is not included, the zero pressure will not be located at the free energy minimum. In Ref. Burgio2002PRC66 , this term has been considered in the calculation of hadron-quark phase transition in dense matter and neutron stars within the bag model.
The derivation process of the pressure is also instructive to the case when one wants to include the Coulomb contribution. In this case, the energy density, accordingly the thermodynamic potential density, gets a term $`E_{\mathrm{Coul}}(\{n_k\},R)`$. Correspondingly the pressure gets
$$P_{\mathrm{Coul}}=E_{\mathrm{Coul}}+\underset{i}{}n_i\frac{E_{\mathrm{Coul}}}{n_i}\frac{R}{3}\frac{E_{\mathrm{Coul}}}{R}.$$
(33)
In literature, there is another approach to have thermodynamic consistency by adding a term $`B^{}`$ to the original thermodynamic potential density. This makes the total thermodynamic potential density becomes $`\mathrm{\Omega }+B^{}`$ Gorenstein1995PRD52 ; Schertler1997NPA616 ; Wangp2000PRC62 . We comment that this can be unconditionally done merely in two special cases, i.e., at finite temperature with zero chemical potential and at finite chemical potential with zero temperature. The expression of the added term for the former case is Gorenstein1995PRD52 :
$$B^{}(T)=_{T_0}^T\frac{\mathrm{\Omega }}{m}|_{T,\mu =0}\frac{dm}{dT}dT,$$
(34)
while that for the later case is Schertler1997NPA616 :
$$B^{}(\mu )=_{\mu _0}^\mu \frac{\mathrm{\Omega }}{m}|_{T=0,\mu }\frac{dm}{d\mu }d\mu .$$
(35)
However, if one wants to extend this to both chemical potential/density and temperature dependent masses, difficulties arise. To perform the integration in Eqs. (34) and (35), one should use $`m=m(T)`$ to replace the $`m`$ on the right hand side of Eq. (34) or apply $`m=m(\mu )`$ to the right hand side of Eq. (35). If $`m_i=m_i(T,\{\mu _k\})`$, the directly extended integration is
$$B^{}(T,\mu )=\left[\underset{i}{}\frac{\mathrm{\Omega }}{m_i}\frac{m_i}{T}dT+\underset{i,j}{}\frac{\mathrm{\Omega }}{m_i}\frac{m_i}{\mu _j}d\mu _j\right].$$
(36)
In the above, ($`T_0,\mu _0`$) is some reference point. Eq. (36) is a multi-dimensional integration. To let it be path-independent, one must mathematically require Cauchy conditions. If all chemical potentials are equal, for example, the Cauchy condition is
$$\underset{i}{}\left[\frac{^2\mathrm{\Omega }}{m_i\mu }\frac{m_i}{T}\frac{^2\mathrm{\Omega }}{m_iT}\frac{m_i}{\mu }\right]=0.$$
(37)
Such an example has recently been given in Ref. Peshier2000PRC61 , where the Cauchy condition Eq. (37), or the equivalent Maxwell relation Eq. (7) in Ref. Peshier2000PRC61 , is fulfilled by solving the equation for the coupling in the masses, As the Eq. (8) of Ref. Peshier2000PRC61 indicated. However, if the masses are completely determined from other arguments, and so there are no adjustable parameters, this thermodynamic treatment may fail. When masses depend on density and temperature, rather than directly on chemical potential and temperature, the case becomes much more involved. Also, if finite size effects can not be ignored, or in other words, $`\mathrm{\Omega }`$ depends explicitly on the volume or radius, not merely the integration in Eq. (36) becomes much more difficult or impossible, the integration in Eq. (34) or in Eq. (35) has an unknown function of the volume or radius as well.
These difficulties are not surprising. In fact, we have proved, from Eq. (6) to Eq. (23), that Eqs. (25)-(28) are inevitable consequences of Eq. (3) with Eq. (24). In this paper, we will apply these formulas to the calculation of properties of both bulk SQM and strangelets.
## III Derivation of quark mass scaling
In the preceding section, we have derived the thermodynamic formulas suitable for systems with density and/or temperature dependent masses. In this section, we derive quark mass scaling by a similar method as in Ref. Peng2000PRC61 .
Letโs schematically write the QCD hamiltonian density for the three flavor case as
$$H_{\mathrm{QCD}}=H_\mathrm{k}+\underset{q=u,d,s}{}m_{q0}\overline{q}q+H_\mathrm{I},$$
(38)
where $`m_{q0}`$ (q=u,d,s) are the quarkโs current mass, $`H_\mathrm{k}`$ is the kinetic term, $`H_\mathrm{I}`$ is the interacting part.
Now we want to include interaction effects within an equivalent mass $`m_\mathrm{q}`$. For this purpose we define an hamiltonian density of the form
$$H_{\mathrm{eqv}}=H_\mathrm{k}+\underset{q=u,d,s}{}m_q\overline{q}q,$$
(39)
where $`m_q`$ is our equivalent mass to be determined. We firstly divide it into two parts, i.e.
$$m_q=m_{q0}+m_\mathrm{I}.$$
(40)
The first part $`m_{q0}`$ $`(q=u,d,s)`$ are the quark current masses while $`m_\mathrm{I}`$ is a common part for all the three flavors to mimic the strong interaction. Obviously we must require that the two hamiltonian densities $`H_{\mathrm{eqv}}`$ and $`H_{\mathrm{QCD}}`$ have the same eigenenergy for any eigenstate $`|\mathrm{\Psi }`$, i.e.
$$\mathrm{\Psi }|H_{\mathrm{eqv}}|\mathrm{\Psi }=\mathrm{\Psi }|H_{\mathrm{QCD}}|\mathrm{\Psi }.$$
(41)
Applying this equality, respectively, to the state $`|n_\mathrm{b},T`$ and the vacuum $`|0`$, and then taking the difference, we have
$$H_{\mathrm{eqv}}_{n_\mathrm{b},T}H_{\mathrm{eqv}}_0,=H_{\mathrm{QCD}}_{n_\mathrm{b},T}H_{\mathrm{QCD}}_0,$$
(42)
where the symbol definitions $`O_{n_\mathrm{b},T}n_\mathrm{b},T|O|n_\mathrm{b},T`$ and $`O_00|O|0`$ have been used for an arbitrary operator $`O`$. Then solving for $`m_\mathrm{I}`$ from this equation gives
$$m_\mathrm{I}=\frac{H_\mathrm{I}}{\underset{q=u,d,s}{}\left[\overline{q}q_{n_\mathrm{b},T}\overline{q}q_0\right]},$$
(43)
where $`H_\mathrm{I}H_\mathrm{I}_{n_\mathrm{b},T}H_\mathrm{I}_0`$ is the interacting part of the energy density from strong interactions between quarks. It can be linked to density $`n_\mathrm{b}`$ and temperature $`T`$ by
$$H_\mathrm{I}=3n_\mathrm{b}\mathrm{v}(\overline{r},T).$$
(44)
Here
$$\overline{r}=\left(\frac{2}{\pi n_\mathrm{b}}\right)^{1/3}$$
(45)
is the average distance of quarks at the density $`n_\mathrm{b}`$, $`\text{v}(\overline{r},T)`$ is the interaction between quarks at density $`n_\mathrm{b}`$ and temperature $`T`$. Because we are interested in the confinement effect, while the lattice simulation Belyaev1984PLB136 and string model investigation Isgur1983PLB124 show that the confinement is linear, we write
$$\text{v}(n_\mathrm{b},T)=\sigma (T)\overline{r}.$$
(46)
The temperature dependence of the string tension $`\sigma (T)`$ can be obtained by combining the Eqs. (94) and (91) in Ref. Ukawa :
$$\sigma (T)=\sigma _0\frac{4T}{a}\mathrm{exp}\left(\frac{2\sigma _0a}{T}\right),$$
(47)
where $`a`$ is the lattice spacing while $`\sigma _0`$ is the string tension at zero temperature. The value of $`\sigma _0`$ from potential models varies in the range of (0.18, 0.22) GeV<sup>2</sup> Veseli .
For convenience, letโs define a dimensionless constant $`\lambda 2\sigma _0a/T_\mathrm{c}`$ with $`T_c`$ being the critical temperature. Then substituting $`a=\lambda T_\mathrm{c}/(2\sigma _0)`$ into Eq. (47) gives
$$\sigma (T)=\sigma _0\left[1\frac{8T}{\lambda T_\mathrm{c}}\mathrm{exp}\left(\lambda \frac{T_\mathrm{c}}{T}\right)\right].$$
(48)
Because the string tension should become zero at the deconfinement temperature, the value of $`\lambda `$ is determined by the equation $`\sigma (T_\mathrm{c})=0`$ whose solution is
$$\lambda =\text{LambertW}(8)1.60581199632.$$
(49)
Accordingly, Eq. (46) becomes
$$\text{v}(n_\mathrm{b},T)=\frac{(2/\pi )^{1/3}\sigma _0}{n_\mathrm{b}^{1/3}}\left(1\frac{8T}{\lambda T_\mathrm{c}}e^{\lambda T_\mathrm{c}/T}\right).$$
(50)
The inter-quark potential has been also studied by comparison to lattice data in Ref. Wong65PRC034902 . Replacing the factor $`\mathrm{exp}(\mu r)`$ in the Eq. (2.8) there with $`1\mu r`$, one will find the linear confining part. The temperature factor used there is $`1(T/T_\mathrm{c})^2`$, and this temperature factor has also been used in Ref. Zhangy2002PRC65 . However, because of a problem with the radius-temperature relation mentioned in the introduction part, the factor has been modified to $`10.65T/T_\mathrm{c}0.35(T/T_\mathrm{c})^2`$ Zhangy2001EPL56 ; Zhangy2003PRC67 ; Zhangy2003MPLA18 ; Zhangy2004JPG30 . The temperature factor derived in the present paper is $`1\frac{8T}{\lambda T_\mathrm{c}}e^{\lambda T_\mathrm{c}/T}.`$ These temperature factors are compared in Fig. 1.
Chiral condensates have been extensively studied in literature qcrev ; Peng747NPA75 . In principle, they depend on both density and temperature. Presently for simplicity, we only consider their density dependence and use the model-independent result Drukarev353ZPA455 ; Cohen1991PRL67
$$\frac{\overline{q}q_{n_\mathrm{b}}}{\overline{q}q_0}=1\frac{n_\mathrm{b}}{\rho ^{}},$$
(51)
with
$$\rho ^{}=\frac{m_\pi ^2f_\pi ^2}{\sigma _\mathrm{N}}.$$
(52)
When taking 140 MeV for the pion mass $`m_\pi `$, 93.3 MeV for the pion decay constant $`f_\pi `$, and 45 MeV for pion-nucleon sigma term $`\sigma _\mathrm{N}`$, one has $`\rho ^{}`$ 0.49 fm<sup>-3</sup>.
Substituting Eqs. (51), (50), and (44) to Eq. (43), we obtain
$$m_\mathrm{I}(n_\mathrm{b},T)=\frac{D}{n_\mathrm{b}^z}\left[1\frac{8T}{\lambda T_c}\mathrm{exp}\left(\lambda \frac{T_\mathrm{c}}{T}\right)\right].$$
(53)
Here $`z=1/3`$. Please note, many density and temperature independent constants such as the vacuum condensates, the string tension at zero temperature, and $`\rho ^{}`$ et al. are grouped into a constant $`D`$, i.e.,
$$D=\frac{3(2/\pi )^{1/3}\sigma _0\rho ^{}}{_q\overline{q}q_0}.$$
(54)
Taking $`\sigma _0=0.18`$ GeV<sup>2</sup> and $`\rho ^{}=0.49`$ fm<sup>-3</sup>, $`\sqrt{D}`$ value is in the range of $`(147,270)`$ MeV when $`_q\overline{q}q_0`$ varies from 3$`\times `$(300 MeV)<sup>3</sup> to 3$`\times `$(200 MeV)<sup>3</sup>. Becasue of many uncertainties within the relevant quantities and the fact that Eq. (51) is not exactly valid, at least it ignores temperature dependence and higher orders in density, we do not try to use the relevant quantities to calculate the value for $`D`$ from Eq. (54). Instead, we treat $`D`$ as a free parameter to be determined by stability arguments, i.e., it makes the energy per baryon $`E/n_\mathrm{b}`$ at zero temperature is greater than 930 MeV for two flavor quark matter in order not to contradict standard nuclear physics, but less than 930 MeV for three flavor quark matter so that SQM can have a chance to be absolutely stable. Obviously, the range determined by this method depends on the thermodynamic formulas and the values of quark current masses. In the present calculation, we use $`m_{u0}=5`$ MeV, $`m_{d0}=10`$ MeV, and $`m_{s0}=120`$ MeV for the current masses involved. These conditions constrain $`\sqrt{D}`$ to a very narrow range of (154.8278, 156.1655) MeV, and we take $`D=(156\text{MeV})^2`$ in this paper.
The model described in the above is a combination of a chiral model and a string model. There should be nothing really wrong with it. In fact, its zero-temperature form has been successfully applied to studying the properties of SQM Peng2000PRC62 ; Peng2000PRC61 , calculating the damping time scale of strange stars due to the coupling of the viscosity and r mode Zhengxp2004PRC70 , and investigating the quark di-quark equation of state and compact star structure Lugonese2003IJMPD12 . Furthermore, various applications of the conventional quasi particle models with chemical and temperature dependent masses have turned out to be very successful Peshier2002PRD66 ; Peshier2000PRC61 . Therefore, we will apply the specific model presented here to the investigation of SQM in bulk and strangelets in the following.
## IV Properties of bulk strange quark matter at finite temperature
As usually done in this model, the quasiparticle contribution to the total thermodynamic potential density of SQM is written as
$$\mathrm{\Omega }=\underset{i}{}\mathrm{\Omega }_i(T,\mu _i,m_i).$$
(55)
where the summation index $`i`$ goes over $`u`$, $`d`$, $`s`$ quarks and electrons. Anti-particles are treated as a whole with particles, i.e., the contribution of the particle type $`i`$ to the thermodynamic potential density is
$`\mathrm{\Omega }_i`$ $`=`$ $`{\displaystyle \frac{g_iT}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}\{\mathrm{ln}[1+e^{(ฯต_{i,p}\mu _i)/T}]`$ (56)
$`+\mathrm{ln}[1+e^{(ฯต_{i,p}+\mu _i)/T}]\}p^2\text{d}p,`$
where $`m_i`$, $`\mu _i`$, and $`T`$ are, respectively, the particle masses, chemical potentials, temperature, and $`ฯต_{i,p}=(p^2+m_i^2)^{1/2}`$ is the dispersion relation. The particle number density corresponding to the particle type $`i`$ is obtained by $`n_i=\mathrm{\Omega }/\mu _i`$, giving
$$n_i=\frac{g_i}{2\pi ^2}_0^{\mathrm{}}\left[\frac{1}{1+e^{(ฯต_{i,p}\mu _i)/T}}\frac{1}{1+e^{(ฯต_{i,p}+\mu _i)/T}}\right]p^2\text{d}p.$$
(57)
The energy density is $`E=_iE_i(T,\mu _i,m_i)`$ with
$`E_i`$ $`=`$ $`{\displaystyle \frac{g_i}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}\left[{\displaystyle \frac{ฯต_{i,p}p^2}{1+e^{(ฯต_{i,p}\mu _i)/T}}}+{\displaystyle \frac{ฯต_{i,p}p^2}{1+e^{(ฯต_{i,p}+\mu _i)/T}}}\right]\text{d}p`$ (58)
$`T{\displaystyle \frac{\mathrm{\Omega }_i}{m_i}}{\displaystyle \frac{m_i}{T}}.`$
The free energy density $`F`$, the entropy density $`S`$, and the pressure $`P`$ are, respectively,
$$F=\underset{i}{}F_i=\underset{i}{}\left(\mathrm{\Omega }_i+\mu _in_i\right),$$
(59)
$$S=\underset{i}{}S_i=\underset{i}{}\left(\frac{\mathrm{\Omega }_i}{T}\frac{\mathrm{\Omega }_i}{m_i}\frac{m_i}{T}\right),$$
(60)
$$P=\underset{i}{}P_i=\underset{i}{}\left(\mathrm{\Omega }_i+n_\mathrm{b}\frac{m_i}{n_\mathrm{b}}\frac{\mathrm{\Omega }_i}{m_i}\right).$$
(61)
In the above, the partial derivatives relevant to $`\mathrm{\Omega }_i`$ are
$`{\displaystyle \frac{\mathrm{\Omega }_i}{m_i}}`$ $`=`$ $`{\displaystyle \frac{g_im_i}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}[{\displaystyle \frac{1}{1+e^{(ฯต_{i,p}\mu _i)/T}}}`$ (62)
$`+{\displaystyle \frac{1}{1+e^{(ฯต_{i,p}+\mu _i)/T}}}]{\displaystyle \frac{p^2\text{d}p}{ฯต_{i,p}}}`$
and
$`{\displaystyle \frac{\mathrm{\Omega }_i}{T}}`$ $`=`$ $`{\displaystyle \frac{g_i}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}\{\mathrm{ln}[1+e^{(ฯต_{i,p}\mu _i)/T}]`$ (63)
$`+{\displaystyle \frac{(ฯต_{i,p}\mu _i)/T}{1+e^{(ฯต_{i,p}\mu _i)/T}}}`$
$`+\mathrm{ln}[1+e^{(ฯต_{i,p}+\mu _i)/T}]`$
$`+{\displaystyle \frac{(ฯต_{i,p}+\mu _i)/T}{1+e^{(ฯต_{i,p}+\mu _i)/T}}}\}p^2\mathrm{d}p.`$
The masses of electrons/positrons (0.511 MeV) and neutrinos/anti-neutrinos (if any) are extremely small. So they can be treated to be zero. For massless particles, the relevant integrations in the above can be carried out to give
$`\mathrm{\Omega }_i`$ $`=`$ $`{\displaystyle \frac{g_i}{24}}\left({\displaystyle \frac{\mu _i^4}{\pi ^2}}+2\mu _i^2T^2+{\displaystyle \frac{7}{15}}\pi ^2T^4\right),`$ (64)
$`n_i`$ $`=`$ $`{\displaystyle \frac{g_i}{6}}\mu _i\left(T^2+{\displaystyle \frac{\mu _i^2}{\pi ^2}}\right),`$ (65)
$`S_i`$ $`=`$ $`{\displaystyle \frac{g_i}{6}}T\left(\mu _i^2+{\displaystyle \frac{7}{15}}\pi ^2T^2\right),`$ (66)
$`E_i`$ $`=`$ $`{\displaystyle \frac{g_i}{8}}\left({\displaystyle \frac{\mu _i^4}{\pi ^2}}+2\mu _i^2T^2+{\displaystyle \frac{7}{15}}\pi ^2T^4\right),`$ (67)
$`F_i`$ $`=`$ $`{\displaystyle \frac{g_i}{8}}\left({\displaystyle \frac{\mu _i^4}{\pi ^2}}+{\displaystyle \frac{2}{3}}\mu _i^2T^2{\displaystyle \frac{7}{45}}\pi ^2T^4\right).`$ (68)
At zero temperature, we have the familiar results
$`\mathrm{\Omega }_i`$ $`=`$ $`{\displaystyle \frac{g_i}{48\pi ^2}}[|\mu _i|\sqrt{\mu _i^2m_i^2}(2\mu _i^25m_i^2)`$ (69)
$`+3m_i^4\mathrm{ln}{\displaystyle \frac{|\mu _i|+\sqrt{\mu _i^2m_i^2}}{m_i}}],`$
$`n_i`$ $`=`$ $`{\displaystyle \frac{g_i\mu _i}{6\pi ^2|\mu _i|}}\left(\mu _i^2m_i^2\right)^{3/2},`$ (70)
$`E_i`$ $`=`$ $`{\displaystyle \frac{g_i}{16\pi ^2}}[|\mu _i|\sqrt{\mu _i^2m_i^2}(2\mu _i^2m_i^2)`$ (71)
$`m_i^4\mathrm{ln}{\displaystyle \frac{|\mu _i|+\sqrt{\mu _i^2m_i^2}}{m_i}}],`$
$`{\displaystyle \frac{\mathrm{\Omega }_i}{m_i}}`$ $`=`$ $`{\displaystyle \frac{g_im_i}{4\pi ^2}}[|\mu _i|\sqrt{\mu _i^2m_i^2}`$ (72)
$`m_i^2\mathrm{ln}{\displaystyle \frac{|\mu _i|+\sqrt{\mu _i^2m_i^2}}{m_i}}].`$
The above formulas show that the number density is an odd function of the corresponding chemical potential, i.e., particle and anti-particle numbers are opposite in sign. But other quantities, such as the thermodynamic potential, entropy, energy, and free energy, are all even functions.
Suppose weak equilibrium is always reached within SQM by the weak reactions such as
$$d,su+e+\overline{\nu _e},s+uu+d.$$
(73)
Correspondingly, relevant chemical potentials satisfy
$`\mu _d=\mu _s,`$ (74)
$`\mu _d+\mu _\nu =\mu _u+\mu _e.`$ (75)
We also have the baryon number density equality
$$\frac{1}{3}\left(n_u+n_d+n_s\right)=n_\mathrm{b}$$
(76)
and the charge neutrality condition
$$\frac{2}{3}n_u\frac{1}{3}n_d\frac{1}{3}n_sn_e=0.$$
(77)
Neutrinos are assumed to enter and leave the system freely. So their chemical potential $`\mu _\nu `$ is zero. From Eqs. (64)-(68), they have no contribution at zero temperature, but contribute at finite temperature.
For a given $`T`$ and $`n_\mathrm{b}`$, the quark masses are obtained from Eqs. (40) and (53):
$$m_q=m_{q0}+\frac{D}{n_\mathrm{b}^z}\left[1\frac{8T}{\lambda T_c}\mathrm{exp}\left(\lambda \frac{T_\mathrm{c}}{T}\right)\right].$$
(78)
The corresponding partial derivatives are easy to get
$`{\displaystyle \frac{m_q}{n_\mathrm{b}}}`$ $`=`$ $`{\displaystyle \frac{zD}{n_\mathrm{b}^{z+1}}}\left[1{\displaystyle \frac{8T}{\lambda T_c}}\mathrm{exp}\left(\lambda {\displaystyle \frac{T_\mathrm{c}}{T}}\right)\right],`$ (79)
$`{\displaystyle \frac{m_q}{T}}`$ $`=`$ $`{\displaystyle \frac{8D}{n_\mathrm{b}^z}}\left[{\displaystyle \frac{1}{\lambda T_\mathrm{c}}}+{\displaystyle \frac{1}{T}}\right]\mathrm{exp}\left(\lambda {\displaystyle \frac{T_\mathrm{c}}{T}}\right).`$ (80)
The chemical potentials $`\mu _i`$ $`(i=u,d,s,e)`$ are obtained by solving Eqs. (74)-(77), all other thermodynamic quantities can then be obtained.
Figure 2 shows the entropy per baryon as a function of temperature for different densities. It is an increasing function of the temperature and goes to zero at zero temperature. This is ensured by the fact that we have $`lim_{T0}m_q/T=0`$ from Eq. (80). It is interesting to note that we did not require this in deriving the scaling Eq. (78) in Sec. III. We got this automatically and naturally.
Figure 3 gives the temperature dependence of the energy and free energy. It is obvious that the energy is an increasing function of temperature while the free energy decreases with increasing temperature.
In Fig. 4, we plot the density dependence of the energy and free energy per baryon at different temperature. The free energy minimum corresponds exactly to the zero pressure, satisfying Eq. (32). However, these two points (zero pressure and the minimum) are generally not the same for the energy per baryon at finite temperature. But at zero temperature they coincide because the energy and free energy are equal to each other at zero temperature.
As pointed out in Ref. Zhangy2002PRC65 , a density and temperature dependent mass model should have the ability of investigating phase transition. This is demonstrated in Fig. 5. The figure is plotted by adjusting, for a given density, the temperature to such a value that it gives a definite pressure indicted in the legend. It is obviously shown that all lines go to the deconfinement temperature $`T_\mathrm{c}`$ as the density approaches to zero. But at zero temperature, the density is different for different pressure. Higher densities correspond to higher pressure.
## V Properties of strangelets
To study strangelets, the special problem is to include the finite size effect. We do this by applying the multi-reflection method, originally comprised by Balian and Bloch Balian1970AP60 , later developed by Madsen Madsen1993PRL70 , Farhi, Berger, and Jaffe Farhi1984PRD30 ; Berger1987PRC35 etc., and applied to the mass density and temperature dependent model by Zhang and Su Zhangy2003MPLA18 ; Zhangy2003PRC67 ; Zhangy2004JPG30 . We express the quasiparticle contribution to the thermodynamic potential density of strangelets as $`\mathrm{\Omega }=_i\mathrm{\Omega }_i(T,\mu _i,m_i,R)`$ with
$`\mathrm{\Omega }_i`$ $`=`$ $`T{\displaystyle _0^{\mathrm{}}}\{\mathrm{ln}[1+e^{(\sqrt{p^2+m_i^2}\mu _i)/T}]`$ (81)
$`+\mathrm{ln}[1+e^{(\sqrt{p^2+m_i^2}+\mu _i)/T}]\}n_i^{}\text{d}p`$
where the density of state $`n_i^{}(p,m_i,R)`$ is given in the multi-expansion approach Balian1970AP60 by
$$n_i^{}(p,m_i,R)=g_i\left[\frac{p^2}{2\pi ^2}+\frac{3p}{R}f_\mathrm{S}\left(x_i\right)+\frac{6}{R^2}f_\mathrm{C}\left(x_i\right)\right].$$
(82)
Here $`x_im_i/p`$, the functions $`f_\mathrm{S}(x_i)`$ Farhi1984PRD30 ; Berger1987PRC35 and $`f_\mathrm{C}(x_i)`$ Madsen1993PRL70 are
$$f_\mathrm{S}(x_i)=\frac{1}{4\pi ^2}\text{arctan}(x_i)$$
(83)
and
$$f_\mathrm{C}(x_i)=\frac{1}{12\pi ^2}\left[1\frac{3}{2x_i}\text{arctan}(x_i)\right].$$
(84)
Then all other thermodynamic quantities are straightforward from Eqs. (25)-(29). Here is the number density for flavor $`i`$:
$$n_i=_0^{\mathrm{}}\left[\eta _i^+\eta _i^{}\right]n_i^{}(p,m_i,R)dp,$$
(85)
where $`\eta _i^\pm `$ is the fermion distribution function, i.e.,
$$\eta _i^\pm \frac{1}{1+e^{(\sqrt{p^2+m_i^2}\mu _i)/T}}.$$
(86)
Because we treat the particles and anti-particle as a whole, the number densities $`n_i`$ can be both positive and negative theoretically. A negative particle number means anti-particles. Fig. 6 shows the chemical potential dependence of the number density for various temperature and mass. In general, the thermodynamic potential density is an even function while the number density is an odd function of the chemical potential, and the number density is zero at zero chemical potential. If finite size effects can be ignored, as in Eq. (57) of the preceding section, this function is monotonically increasing, and so positive chemical potentials correspond to positive number density. When the finite size effects are included (the surface term Farhi1984PRD30 ; Berger1987PRC35 and curvature term Madsen1993PRL70 ), the case becomes more complex. When the temperature is very high, the function is still monotonic. However, when the temperature becomes lower than some special value, the function becomes non-monotonic, and the number density is negative for some special positive chemical potentials. For smaller masses, the chemical potentials are smaller. So light quarks do not fall into this region of chemical potentials in actual cases. However, when quark masses become bigger, the corresponding chemical potentials becomes also bigger. It is therefore possible that quarks with comparatively bigger mass happen to have positive chemical potential and negative number density in some special cases. We will see such special examples a little later. One may concern that non-monotone of the number density leads to $`n/\mu <0`$ which violates the stability condition. However, our results are always located in the regime where the derivative is positive. The regime is marked with full lines in Fig. 6.
The free energy density is
$$F=\underset{i}{}(\mathrm{\Omega }_i+\mu _in_i,),$$
(87)
the energy density $`E`$ is
$`E`$ $`=`$ $`{\displaystyle \underset{i}{}}[{\displaystyle _0^{\mathrm{}}}(\eta _i^++\eta _i^{})\sqrt{p^2+m_i^2}n_i^{}(p,m_i,R)\mathrm{d}p`$ (88)
$`T{\displaystyle \frac{\mathrm{\Omega }_i}{m_i}}{\displaystyle \frac{m_i}{T}}],`$
the entropy density is
$$S=\underset{i}{}\left(\frac{\mathrm{\Omega }_i}{T}T\frac{\mathrm{\Omega }_i}{m_i}\frac{m_i}{T}\right),$$
(89)
and the pressure $`P`$ is
$$P=\underset{i}{}\left(\mathrm{\Omega }_i+n_\mathrm{b}\frac{m_i}{n_\mathrm{b}}\frac{\mathrm{\Omega }_i}{m_i}\frac{R}{3}\frac{\mathrm{\Omega }_i}{R}\right).$$
(90)
In the above, the relevant partial derivatives are
$$\frac{\mathrm{\Omega }_i}{T}=_0^{\mathrm{}}\mathrm{ln}\left[\frac{(1\eta _i^+)(1\eta _i^{})}{(1/\eta _i^+1)^{\eta _i^+}(1/\eta _i^{}1)^{\eta _i^{}}}\right]n_i^{}dp,$$
(91)
$`{\displaystyle \frac{R}{3}}{\displaystyle \frac{\mathrm{\Omega }_i}{R}}`$ $`=`$ $`g_i{\displaystyle _0^{\mathrm{}}}\left[2\sqrt{p^2+m_i^2}+T\mathrm{ln}(\eta _i^+\eta _i^{})\right]`$ (92)
$`\times \left[{\displaystyle \frac{p}{R}}f_\mathrm{S}(x_i)+{\displaystyle \frac{4}{R^2}}f_\mathrm{C}(x_i)\right]\mathrm{d}p,`$
and
$`{\displaystyle \frac{\mathrm{\Omega }_i}{m_i}}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\{{\displaystyle \frac{(\eta _i^++\eta _i^{})n_i^{}}{\sqrt{1+p^2/m_i^2}}}`$ (93)
$`+[2\sqrt{p^2+m_i^2}+T\mathrm{ln}(\eta _i^+\eta _i^{})]{\displaystyle \frac{n_i^{}}{m_i}}\}\mathrm{d}p`$
with
$$\frac{n_i^{}}{m_i}=\frac{3g_i}{4\pi ^2m_iR^2}\left[\frac{p}{m_i}\text{arctan}\left(\frac{m_i}{p}\right)\frac{1+Rm_i}{1+m_i^2/p^2}\right].$$
(94)
Quark masses and relevant derivatives are still given by Eqs. (78)-(80).
In Ref. Zhangy2001EPL56 , charge neutrality is also imposed for strangelets. This is convenient for checking whether the formulas are continuous from bulk SQM to finite baryon number. Fig. 7 gives the energy per baryon and chemical potentials versus baryon number at zero temperature and zero pressure. The horizontal lines are the corresponding values for bulk SQM. It is very clear that all quantities approach to the corresponding bulk values with increasing baryon number, and finite size effects destabilize low baryon number strangelets.
Now we treat strangelets in another different way. Instead of imposing the charge neutrality Eq. (77), we require that the electron and positron number densities are zero, i.e.,
$$n_e=0$$
(95)
because the radius of strangelets is much smaller than the Compton wavelength of electrons and positrons. Electrons and positrons are not involved in the strong interaction, or in other words, they are not confined, so the finite size terms in Eq. (82) vanish for them. From Eq. (65), we have
$$n_e=\frac{\mu _e}{3}\left(T^2+\frac{\mu _e^2}{\pi ^2}\right).$$
(96)
Therefore, zero $`n_e`$ means zero $`\mu _e`$. With a view to the chemical equilibrium Eqs. (74) and (75), we naturally get
$$\mu _u=\mu _d=\mu _s.$$
(97)
In fact, Eq. (97) is the condition to find out the configuration, which has the lowest energy per baryon, from the all possible strangelets with a fixed baryon number HeYB53PRC1903 . Due to Eq. (97), only one chemical potential is left independent. And it can then be determined by solving
$$\frac{1}{3}(n_u+n_d+n_s)=\frac{A}{(4/3)\pi R^3}$$
(98)
for a given baryon number $`A`$, temperature $`T`$, and radius $`R`$.
Figure 8 shows the energy and free energy per baryon as a function of the radius for $`A=20`$ at different temperature. The points marked with an open circle are the mechanically stable radius where the pressure is zero. The minimum of each line is marked with a triangle. Again we see that these two points are always the same on the free energy line. But they are different on the energy line at finite temperature. However, they coincide at zero temperature because the energy and free energy are equal at zero temperature.
For a given $`A`$ and $`T`$, the mechanically stable radius is obtained by adjusting it so that the free energy is minimized, or, simply by solving the equation
$$P=0.$$
(99)
The temperature dependence of the stable radius for $`A=20`$ is plotted in Fig. 9 with a solid line. It is obviously an increasing function of temperature. The corresponding energy per baryon is also plotted in the same figure, labeled on the right axis with a dotted line. It is also an increasing function of temperature. However, the free energy per baryon (dashed line) decreases with increasing temperature.
Because charge neutrality is not imposed, strangelets here are charged. The electric charge can be calculated by
$$Z=\frac{4}{3}\pi R^3\left(\frac{2}{3}n_u\frac{1}{3}n_d\frac{1}{3}n_s\right).$$
(100)
In Fig. 10, we plot the charge to baryon number ratio as a function of the baryon number at different temperature. This figure shows that the charge to baryon number ratio is a decreasing function of the baryon number and temperature.
A noticeable feature is that the charge to baryon number ratio at lower baryon numbers ($`A5`$ at zero temperature) becomes greater than 1/2. This is very different from that in the bag model where the charge to baryon number ratio is very small. For normal nuclei, this ratio reaches its biggest value 1/2. So it seems difficult to understand a heavy positive charge at first sight. In fact, it is caused by the fact that $`f_s`$ is negative for $`1A5`$ and $`T=0`$, i.e., anti-strange quarks, rather than strange quarks, appear. One can see this phenomenon clearly in Fig. 11 where the free energy per baryon and the chemical potential have also been shown. Because of finite size effects, stranelets with very low baryon numbers are metastable for the parameters chosen. If one choose a bigger value for the current quark mass, the value for $`D`$ should be smaller, and consequently, the mass for strangelets would be smaller.
Fig. 12 shows the quark configuration for low baryon numbers. To see the results clearly, we give the corresponding data in Tab. 1. The first column is the baryon number, the second to fourth column are, respectively, the quark numbers $`N_u`$, $`N_d`$, and $`N_s`$. Because shell effects are not taken into account and beta-equilibrium is imposed, fractional quark numbers appear in Tab. 1. Actual quark numbers should naturally be integers. So we approximate these real numbers to integers by $`\text{int}(N_i)+N_i/|N_i|`$ ($`i=u,d,s`$; int means the number before the decimal point). The results are shown in the fifth column. For $`A=1`$, we have the pentaquark state ($`u^2d^2s^1`$). For $`A=2`$, we have the dibaryon ($`u^4d^3s^1`$) or octaquark state. For $`A=3`$, 4, and 5, we respectively have the multi-quark states ($`u^5d^5s^1`$), ($`u^7d^6s^1`$), and ($`u^8d^8s^1`$). A common feature of these states is that they all include an anti-strange quark. So we use โ$`\overline{s}`$letโ as the title of the fifth column. The charge number of these $`\overline{s}`$lets are given in the sixth column, while the seventh column gives the energies calculated by the present parameters ($`D^{1/2}=156`$ MeV and $`m_{s0}=120`$ MeV) with perfect $`\beta `$ equilibrium. If one would like to produce 1540 MeV (the actual $`\mathrm{\Theta }^+`$ resonant mass) for $`A=1`$, then one has to take $`D^{1/2}=186`$ MeV and get 2856 MeV for $`A=2`$. So we expect that the mass of the octaquark ($`u^4d^3s^1`$), if truely exists, is near 2856 MeV. For $`D^{1/2}=186`$ MeV and $`A3`$, the strange quark number becomes positive. So in this case we have only the pentaquark and the octaquark. Because of uncertainties in parameters, and also many other factors e.g. the perturbative interaction has not been included (the quark mass scaling is derived by assuming that the linear confinement interaction dominates), the concrete values should not be taken seriously, and further studies are needed.
Recently, the pentaquark state $`\mathrm{\Theta }^+`$(1540) has aroused a lot of interest Eidelman2004PLB592 . The width of $`\mathrm{\Theta }^+`$(1540) is very narrow with upper limit as small as 9 MeV Nakano2004MPLA19 . Cahn and trilling have extracted $`\mathrm{\Gamma }(\mathrm{\Theta }^+)=0.9\pm 0.3`$ MeV from an analysis of Xenon bubble chamber Cahn2004PRD69 . Although other hadrons like $`\varphi (1020)`$ and $`\mathrm{\Lambda }(1520)`$ are also narrow \[$`\mathrm{\Gamma }(\varphi )=4.26\pm 0.05`$ MeV, $`\mathrm{\Gamma }(\mathrm{\Lambda })=15.6\pm 1.0`$ MeV\], we know what makes them narrow. However, we do not know, until now, why the $`\mathrm{\Theta }^+`$ should be so stable. We see from the above data that ($`u^2d^2s^1`$) is the multi-quark state which mostly satisfy the weak equilibrium for baryon number 1. This might serve as an explanation for the stability of $`\mathrm{\Theta }^+`$(1540). Other $`\overline{s}`$lets, e.g. the dibaryon ($`u^4d^3s^1`$), an octaquark state, should also exist, yet to be searched for by experiments though.
In Ref. SchaffnerBielich1997PRD55 , strangelets were also calculated to be heavily charged. But the electric charge is negative. There the investigation was concerned with how small metastable strangelets look like and might decay for different lifetime. With the similar ideas, Ref. Zhangy2003PRC67 studied strangelets within density-and-temperature dependent quark masses and extending the findings in Ref. SchaffnerBielich1997PRD55 to finite temperature. Present investigation concentrate on deriving thermodynamic formulas and quark mass scaling, and finding the lowest-energy configuration from the strangelets with a fixed baryon number. Naturally, the observation of heavily positively charged strangelets, or multi-quark states with an anti-strange quark, depends on the value of the parameter $`D`$. If we took a much larger $`D`$ value, the charge to baryon number ratio would also be small, or the anti-strange quark would not appear. However, bulk SQM has no chance to be absolutely stable in that case.
## VI Summary
When masses are density and/or temperature dependent, the thermodynamical formulas are different from that for constant masses. We have derived a new set of thermodynamical formulas which can be used to calculate the properties of quark matter within a density and/or temperature dependent quark mass model. The new formulas are also instructive when one introduces a density and/or temperature dependent bag constant in the bag model etc.
We have also argued for a new quark mass scaling at finite temperature. The basic feature is that quark masses and their partial derivative with respect to the temperature go to zero when the temperature approaches to zero. This ensures that all quantities restore to the density dependent model at zero temperature. It is especially important that the entropy goes naturally to zero when the temperature approaches to zero, satisfying the third law of thermodynamics.
With the new thermodynamical formulas and new quark mass scaling, we have studied the properties of bulk SQM and strangelets. It is shown that the free energy minimum corresponds exactly to the zero pressure, both at zero and finite temperature. The mechanically stable strangelet radius increases with temperature. An interesting new observation is that low mass strangelets are heavily positively charged, or appear as multi-quark states with an anti-strange quark, such as the pentaquark $`(u^2d^2\overline{s})`$ and the octaquark ($`u^4d^3\overline{s}`$) etc., if bulk SQM is absolutely stable.
###### Acknowledgements.
The authors would like to thank support from the DOE (DF-FC02-94ER40818) and NSFC (10375074, 90203004, 10475089, 10435080, 10275037). G.X.P also acknowledges hospitality at MIT-CTP. In particular, he is grateful to Prof. E. Farhi, R. Jackiw, R. L. Jaffe, J. W. Negele, K. Rajagopal, and F. Wilczek, for helpful conversations. |
warning/0506/math0506250.html | ar5iv | text | # A short proof of affability for certain Cantor minimal โคยฒ-systems
## 1 Introduction
In this paper, we would like to investigate the orbit structure of certain minimal dynamical systems on a Cantor set. Giordano, Putnam and Skau proved that equivalence relations arising from $``$-actions are orbit equivalent to AF equivalence relations in \[GPS1\]. Moreover, they gave the classification for AF equivalence relations. In a recent paper \[GPS3\], they continued their investigations and showed that equivalence relations arising from $`^2`$-actions are again orbit equivalent to AF equivalence relations under a hypothesis involving the existence of cocycles. An equivalence relation which is orbit equivalent to an AF equivalence relation is said to be affable. A crucial ingredient of their proof was the absorption theorem given in \[GPS2\]. They needed, however, sufficiently many cocycles in order to construct an AF subequivalence relation to which the absorption theorem can be applied. The aim of this paper is to show that the existence of cocycles is not necessary for certain $`^2`$-actions. We will give a short proof that the associated equivalence relations are orbit equivalent to AF equivalence relations, thus they are affable.
We recall some terminology that we shall use. Let $`X`$ be a Cantor set and let $``$ be an รฉtale equivalence relation on $`X`$. We define the $``$-equivalence class $`[x]_{}`$ of $`xX`$ by $`[x]_{}=\{yX:(x,y)\}`$. The equivalence relation $``$ is said to be minimal, if $`[x]_{}`$ is dense in $`X`$ for every $`xX`$. Let $`\phi :G\mathrm{Homeo}(X)`$ be a free action of a countable discrete group $`G`$, that is, $`\phi `$ is a group homomorphism and $`\phi ^g(x)x`$ for all $`xX`$ and $`gG\{e\}`$, where $`e`$ means the identity element. We put
$$_\phi =\{(x,\phi ^g(x))X\times X:xX,gG\}.$$
By transferring the product topology on $`X\times G`$ via the bijection $`(x,g)(x,\phi ^g(x))`$, we can topologize $`_\phi `$. It is easily verified that $`_\phi `$ becomes an รฉtale equivalence relation. We call $`(X,\phi )`$ a Cantor minimal $`G`$-system when $`_\phi `$ is minimal. In this paper, we deal with only Cantor minimal $``$-systems and Cantor minimal $`^2`$-systems.
Let $`(X,\phi )`$ and $`(Y,\psi )`$ be two Cantor minimal $`^2`$-systems. We say that $`\pi :(Y,\psi )(X,\phi )`$ is a factor map when $`\pi :YX`$ is a continuous map and $`\pi \psi ^a=\phi ^a\pi `$ for all $`a^2`$. The system $`(Y,\psi )`$ is called an extension of $`(X,\phi )`$. Our main theorem asserts that $`_\psi `$ is affable, if $`(X,\phi )`$ is conjugate to a product of two Cantor minimal $``$-systems. Suppose that $`(X,\phi )`$ is conjugate to the product of two Cantor minimal $``$-systems $`(X_1,\phi _1)`$ and $`(X_2,\phi _2)`$. From \[GPS1, Theorem 2.3\] we can see that $`_{\phi _1}`$ and $`_{\phi _2}`$ are affable. Since a product of AF equivalence relations is also AF, it is easily checked that $`_\phi `$ is affable. But, it looks impossible to mimic this simple argument in the case of the extension $`(Y,\psi )`$, because $`(Y,\psi )`$ itself is not a product. We will instead construct a โniceโ AF subequivalence relation of $`_\phi `$ and apply the absorption theorem to this relation.
## 2 Products of Cantor minimal $``$-systems
Throughout this section, let $`B_i=(V_i,E_i)`$ be simple properly ordered Bratteli diagrams for $`i=1,2`$. For each $`i=1,2`$, $`V_i`$ and $`E_i`$ can be written as a countable disjoint union of non-empty finite sets:
$$V_i=V_{i,0}V_{i,1}V_{i,2}\mathrm{}\text{ and }E_i=E_{i,1}E_{i,2}E_{i,3}\mathrm{}$$
with the source map $`s:E_{i,n}V_{i,n1}`$ and the range map $`r:E_{i,n}V_{i,n}`$. Without loss of generality, we may assume that all two vertices in consecutive levels are connected by more than three edges. We write the infinite path space associated with $`B_i`$ by $`X_i`$ for each $`i=1,2`$. Let $`p_i`$ be the unique maximal infinite path of $`X_i`$ and let $`\phi _i\mathrm{Homeo}(X_i)`$ be the Bratteli-Vershik transformation on $`X_i`$ (see \[HPS\]). It is well known that $`(X_i,\phi _i)`$ is a Cantor minimal $``$-system.
Set $`X=X_1\times X_2`$. Let $`\phi :^2\mathrm{Homeo}(X)`$ be the $`^2`$-action on $`X`$ induced by $`\phi _1\times \mathrm{id}`$ and $`\mathrm{id}\times \phi _2`$. Then, $`(X,\phi )`$ is a Cantor minimal $`^2`$-system.
For each $`n`$, we put
$$_n=\{((x_1,x_2),(y_1,y_2))X\times X:x_{i,m}=y_{i,m}\text{ for }i=1,2\text{ and }m>n\},$$
where $`x_{i,m},y_{i,m}E_{i,m}`$ mean the $`m`$-th coordinate of the infinite paths $`x_i,y_iX_i`$. It is not hard to see that $`_n`$ is a compact open subequivalence relation of $`_\phi `$ with the relative topology from $`_\phi `$. Therefore
$$=\underset{n}{}_n$$
is an AF subequivalence relation of $`_\phi `$. Note that $``$ is minimal because $`B_1`$ and $`B_2`$ are simple.
For $`i,j=0,1`$ and $`n`$, we define continuous functions $`\lambda _n^{ij}:X\{0,1\}`$ inductively as follows. Let $`(x_1,x_2)X`$. We denote the $`n`$-th coordinate of $`x_i`$ by $`x_{i,n}E_{i.n}`$. At first, put
$$\lambda _1^{00}(x_1,x_2)=\{\begin{array}{cc}1\hfill & x_{1,1}\text{ is maximal}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
$$\lambda _1^{01}(x_1,x_2)=\{\begin{array}{cc}1\hfill & x_{2,1}\text{ is minimal}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
$$\lambda _1^{11}(x_1,x_2)=\{\begin{array}{cc}1\hfill & x_{1,1}\text{ is minimal}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
and
$$\lambda _1^{10}(x_1,x_2)=\{\begin{array}{cc}1\hfill & x_{2,1}\text{ is maximal}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
Then, for $`n2`$, we define $`\lambda _n^{ij}`$ by
$$\lambda _n^{00}(x_1,x_2)=\{\begin{array}{cc}\lambda _{n1}^{00}(x_1,x_2)\hfill & x_{1,n}\text{ is maximal and }x_{2,n}\text{ is maximal}\hfill \\ 1\hfill & x_{1,n}\text{ is maximal and }x_{2,n}\text{ is not maximal}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
$$\lambda _n^{01}(x_1,x_2)=\{\begin{array}{cc}\lambda _{n1}^{01}(x_1,x_2)\hfill & x_{1,n}\text{ is maximal and }x_{2,n}\text{ is minimal}\hfill \\ 1\hfill & x_{1,n}\text{ is not maximal and }x_{2,n}\text{ is minimal}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
$$\lambda _n^{11}(x_1,x_2)=\{\begin{array}{cc}\lambda _{n1}^{11}(x_1,x_2)\hfill & x_{1,n}\text{ is minimal and }x_{2,n}\text{ is minimal}\hfill \\ 1\hfill & x_{1,n}\text{ is minimal and }x_{2,n}\text{ is not minimal}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
and
$$\lambda _n^{10}(x_1,x_2)=\{\begin{array}{cc}\lambda _{n1}^{10}(x_1,x_2)\hfill & x_{1,n}\text{ is minimal and }x_{2,n}\text{ is maximal}\hfill \\ 1\hfill & x_{1,n}\text{ is not minimal and }x_{2,n}\text{ is maximal}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
It is easily checked that $`\lambda _n^{ij}`$ is well-defined and continuous.
The following is an immediate consequence of the definition of $`\lambda _n^{ij}`$.
###### Lemma 2.1.
Let $`(i,j)\{0,1\}^2`$. For $`((x_1,x_2),(y_1,y_2))_n`$, if
$$\lambda _n^{ij}(x_1,x_2)=\lambda _n^{ij}(y_1,y_2),$$
then we have
$$\lambda _m^{ij}(x_1,x_2)=\lambda _m^{ij}(y_1,y_2),$$
for all $`m>n`$.
For every $`n`$, we define a subset $`_n^{}`$ of $`_n`$ by
$$_n^{}=\{((x_1,x_2),(y_1,y_2))_n:\lambda _n^{ij}(x_1,x_2)=\lambda _n^{ij}(y_1,y_2)\text{ for all }i,j=0,1\}.$$
###### Lemma 2.2.
For every $`n`$, $`_n^{}`$ is a compact open subequivalence relation of $`_n`$, and $`_n^{}`$ is contained in $`_{n+1}^{}`$.
###### Proof.
It is obvious that $`_n^{}`$ is a subequivalence relation of $`_n`$. Since $`\lambda _n^{ij}`$ is continuous, $`_n^{}`$ is compact and open. From the lemma above we can see $`_n^{}_{n+1}^{}`$. โ
Define
$$^{}=\underset{n}{}_n^{}.$$
By the lemma above, $`^{}`$ is an AF equivalence relation on $`X`$.
###### Lemma 2.3.
Let $`((x_1,x_2),(y_1,y_2))`$.
1. If $`x_1`$ is not in $`\{\phi _1^n(p_1):n\}`$, then $`((x_1,x_2),(y_1,y_2))^{}`$.
2. If $`x_2`$ is not in $`\{\phi _2^n(p_2):n\}`$, then $`((x_1,x_2),(y_1,y_2))^{}`$.
###### Proof.
It suffices to show (1). There exists $`n`$ such that $`((x_1,x_2),(y_1,y_2))_n`$. We can find a natural number $`m>n`$ such that $`x_{1,m}`$ is not maximal. Then, $`\lambda _m^{00}(x_1,x_2)`$ equals zero. From $`x_{1,m}=y_{1,m}`$, we get $`\lambda _m^{00}(x_1,x_2)=\lambda _m^{00}(y_1,y_2)=0`$. It is easy to see that $`\lambda _m^{01}(x_1,x_2)`$ depends only on $`x_{2,m}`$, and so we have $`\lambda _m^{01}(x_1,x_2)=\lambda _m^{01}(y_1,y_2)`$.
We can find a natural number $`l>n`$ such that $`x_{1,l}`$ is not minimal. It is clear that $`\lambda _l^{11}(x_1,x_2)=0`$ and $`\lambda _l^{10}(x_1,x_2)`$ depends only on $`x_{2,l}`$. In a similar fashion to the preceding paragraph, we get $`\lambda _l^{11}(x_1,x_2)=\lambda _l^{11}(y_1,y_2)`$ and $`\lambda _l^{10}(x_1,x_2)=\lambda _l^{10}(y_1,y_2)`$.
By virtue of Lemma 2.1, we can conclude that $`((x_1,x_2),(y_1,y_2))`$ is in $`_k^{}`$, where $`k`$ is the maximum of $`m`$ and $`l`$. โ
Put $`p=(p_1,p_2)X`$. The above lemma tells us that the four $``$-orbits $`[p]_{}`$, $`[\phi ^{(1,0)}(p)]_{}`$, $`[\phi ^{(0,1)}(p)]_{}`$ and $`[\phi ^{(1,1)}(p)]_{}`$ may split in $`^{}`$, but the other $``$-orbits do not split in $`^{}`$.
###### Lemma 2.4.
The equivalence relation $`^{}`$ is minimal.
###### Proof.
Let $`(x_1,x_2)X`$. It suffices to show that $`[(x_1,x_2)]_{^{}}`$ is dense in $`X`$. If $`x_1`$ does not belong to $`\{\phi _1^n(p_1):n\}`$ or $`x_2`$ does not belong to $`\{\phi _2^n(p_2):n\}`$, then we have nothing to do, because the $`^{}`$-orbit of $`(x_1,x_2)`$ is equal to the $``$-orbit of it. Suppose that $`(x_1,x_2)`$ is in $`\{\phi ^a(p):a^2\}`$. Without loss of generality, we may assume that $`(x_1,x_2)`$ belongs to $`[p]_{}`$. Take finite paths $`(e_{1,1},e_{1,2},\mathrm{},e_{1,n})`$ in $`B_1`$ and $`(e_{2,1},e_{2,2},\mathrm{},e_{2,n})`$ in $`B_2`$. Thus $`e_{i,k}E_{i,k}`$ and $`r(e_{i,k})=s(e_{i,k+1})`$. We can find $`m>n+2`$ such that both $`x_{1,m}`$ and $`x_{2,m}`$ are maximal. It follows that $`\lambda _m^{01}(x_1,x_2)=0`$, $`\lambda _m^{11}(x_1,x_2)=0`$ and $`\lambda _m^{10}(x_1,x_2)=1`$. We have two possibilities: $`\lambda _m^{00}(x_1,x_2)=0`$ or $`1`$.
Let us consider the case that $`\lambda _m^{00}(x_1,x_2)`$ is one. We can find edges $`e_{i,k}E_{i,k}`$ for $`i=1,2`$ and $`k=n+1,n+2,\mathrm{},m1`$ such that the following are satisfied.
* $`r(e_{i,k})=s(e_{i,k+1})`$ and $`r(e_{i,m1})=s(x_{i,m})`$ for all $`i=1,2`$ and $`k=n,n+1,\mathrm{},m2`$.
* $`e_{1,m1}`$ is maximal and $`e_{2,m1}`$ is not maximal.
Put
$$x_i^{}=(e_{i,1},e_{i,2},\mathrm{},e_{i,n},e_{i,n+1},\mathrm{},e_{i,m1},x_{i,m},x_{i,m+1}\mathrm{})X_i$$
for each $`i=1,2`$. Then it is clear that $`((x_1,x_2),(x_1^{},x_2^{}))_m`$. Moreover, it is not hard to see $`\lambda _m^{00}(x_1^{},x_2^{})=1`$, $`\lambda _m^{01}(x_1^{},x_2^{})=0`$, $`\lambda _m^{11}(x_1^{},x_2^{})=0`$ and $`\lambda _m^{10}(x_1^{},x_2^{})=1`$. Therefore we get $`((x_1,x_2),(x_1^{},x_2^{}))_m^{}`$.
Suppose that $`\lambda _m^{00}(x_1,x_2)`$ is zero. In this case we choose the edges $`e_{i,k}E_{i,k}`$ so that the following are satisfied.
* $`r(e_{i,k})=s(e_{i,k+1})`$ and $`r(e_{i,m1})=s(x_{i,m})`$ for all $`i=1,2`$ and $`k=n,n+1,\mathrm{},m2`$.
* $`e_{1,m1}`$ is not maximal.
Then, we can again obtain $`((x_1,x_2),(x_1^{},x_2^{}))_m^{}`$. Hence we can conclude that the $`^{}`$-orbit of $`(x_1,x_2)`$ is dense in $`X`$. โ
###### Lemma 2.5.
For every $`m\{1\}`$ we have the following.
1. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(p_1,\phi _2^{1m}(p_2))=(1,0,0,1)`$.
2. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(\phi _1(p_1),\phi _2^{1m}(p_2))=(0,0,1,0)`$.
3. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(p_1,\phi _2^m(p_2))=(1,0,0,0)`$.
4. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(\phi _1(p_1),\phi _2^m(p_2))=(0,1,1,0)`$.
5. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(\phi _1^{1m}(p_1),p_2)=(0,0,0,1)`$.
6. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(\phi _1^{1m}(p_1),\phi _2(p_2))=(1,1,0,0)`$.
7. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(\phi _1^m(p_1),p_2)=(0,0,1,1)`$.
8. $`lim_n\mathrm{}(\lambda _n^{00}\times \lambda _n^{01}\times \lambda _n^{11}\times \lambda _n^{10})(\phi _1^m(p_1),\phi _2(p_2))=(0,1,0,0)`$.
###### Proof.
Straightforward computation. โ
Take a clopen subset $`U_iX_i`$ which does not contain $`p_i`$ and $`\phi _i(p_i)`$ for each $`i=1,2`$. Put
$$B=(\{p_1\}\times U_2)(U_1\times \{p_2\})$$
and
$$B^{}=(\{\phi _1(p_1)\}\times U_2)(U_1\times \{\phi _2(p_2)\}).$$
###### Lemma 2.6.
Both $`B`$ and $`B^{}`$ are closed $`^{}`$-รฉtale thin subsets.
###### Proof.
It suffices to show the statement for $`B`$. Suppose that $`((x_1,x_2),(y_1,y_2))`$ is in $`^{}(B\times B)`$. Without loss of generality, we may assume $`x_1=p_1`$. Suppose $`y_2=p_2`$. Then $`x_2`$ must be $`\phi _2^{1m}(p_2)`$ for some $`m`$, and $`m`$ is not one because $`p_2`$ is not in $`U_2`$. Similarly $`y_1`$ must be $`\phi _1^{1l}(p_1)`$ for some $`l\{1\}`$. But $`((p_1,\phi _1^{1m}(p_2)),(\phi _1^{1l}(p_1),p_2))`$ never be in $`^{}`$ by the lemma above. Hence we have $`y_1=p_1`$. Thus $`((x_1,x_2),(y_1,y_2))`$ is equal to $`((p_1,x_2),(p_1,\phi _2^m(x_2)))`$ for some $`m`$. Define
$$V=\{((a,b),(c,d))^{}:a=c,d=\phi _2^m(b)\text{ and }b,dU_2\}.$$
Then $`V`$ is a clopen neighborhood of $`((p_1,x_2),(p_1,\phi _2^m(x_2)))`$ in $`^{}`$. For $`((a,b),(c,d))V`$, it is obvious that $`(a,b)B`$ if and only if $`(c,d)B`$, which implies that $`B`$ is รฉtale.
We would like to show that a probability measure on $`X=X_1\times X_2`$ is $``$-invariant if and only if it is $`^{}`$-invariant. If this is shown, thinness of $`B`$ easily follows. But, except for countably many $`(x_1,x_2)`$โs, the equivalence class $`[(x_1,x_2)]_{}`$ is equal to $`[(x_1,x_2)]_{^{}}`$. Since every invariant measure is nonatomic, we can finish the proof. โ
###### Lemma 2.7.
We have $`^{}(B\times B^{})=\mathrm{}`$.
###### Proof.
Suppose that $`((x_1,x_2),(y_1,y_2))`$ is contained in $`^{}(B\times B^{})`$. Without loss of generality, we may assume $`x_1=p_1`$. Then $`y_1`$ never be $`\phi _1(x_1)`$, because $`((p_1,x_2),(\phi _1(p_1),y_2))`$ does not belong to $``$. It follows that $`y_2=\phi _2(p_2)`$ and $`((x_1,x_2),(y_1,y_2))=((p_1,\phi _2^m(p_2)),(\phi _1^{1l}(p_1),\phi _2(p_2)))`$ for some $`m,l\{1\}`$. This pair, however, never belongs to $`^{}`$ by virtue of Lemma 2.5, which completes the proof. โ
We define a homeomorphism $`\beta :BB^{}`$ as follows. For $`(p_1,x_2)\{p_1\}\times U_2`$, we put $`\beta (p_1,x_2)=(\phi _1(p_1),x_2)`$. For $`(x_1,p_2)U_1\times \{p_2\}`$, we put $`\beta (x_1,p_2)=(x_1,\phi _2(p_2))`$.
###### Lemma 2.8.
The homeomorphism $`\beta :BB^{}`$ induces an isomorphism between $`^{}(B\times B)`$ and $`^{}(B^{}\times B^{})`$.
###### Proof.
Since the topology of $`^{}(B\times B)`$ and $`^{}(B^{}\times B^{})`$ is inherited from $``$, it suffices to show that $`\beta `$ is a well-defined bijection between $`^{}(B\times B)`$ and $`^{}(B^{}\times B^{})`$. Let $`((x_1,x_2),(y_1,y_2))^{}(B\times B)`$. Without loss of generality, we may assume that $`x_1=p_1`$ and $`x_2U_2`$. By the proof of Lemma 2.6, we get $`y_1=p_1`$ and $`y_2U_2`$. It follows that $`\beta (p_1,x_2)=(\phi _1(p_1),x_2)`$ and $`\beta (p_1,y_2)=(\phi _1(p_1),y_2)`$. If $`x_2`$ does not belong to $`\{\phi _2^n(p_2):n\}`$, then the $`^{}`$-orbit of $`(\phi _1(p_1),x_2)`$ is equal to the $``$-orbit of it. Hence we have $`((\phi _1(p_1),x_2),(\phi _1(p_1),y_2))^{}(B^{}\times B^{})`$.
Suppose that $`x_2`$ and $`y_2`$ belong to $`\{\phi _2^n(p_2):n\}`$. Since $`((p,x_2),(p,y_2))^{}`$, we have two possibilities: both $`x_2`$ and $`y_2`$ belong to $`\{\phi _2^{1n}(p_2):n\}`$, or both $`x_2`$ and $`y_2`$ belong to $`\{\phi _2^n(p_2):n\}`$. Without loss of generality, we may assume the latter. Thus, $`x_2=\phi _2^n(p_2)`$ and $`y_2=\phi _2^m(p_2)`$ for some $`n,m`$. Because $`x_2`$ and $`y_2`$ is in $`U_2`$, $`n`$ and $`m`$ are greater than one. It follows from Lemma 2.5 that $`((\phi _1(p_1),\phi _2^n(p_2)),(\phi _1(p_1),\phi _2^m(p_2)))`$ belongs to $`^{}`$. The proof is completed. โ
###### Lemma 2.9.
Let $`\stackrel{~}{}`$ be the equivalence relation generated by $`^{}`$ and the graph of $`\beta `$. Then $`_\phi `$ is generated by $`\stackrel{~}{}`$ and $`(p,\phi ^{(0,1)}(p))`$, $`(\phi ^{(0,1)}(p),\phi ^{(1,1)}(p))`$ and $`(\phi ^{(1,1)}(p),\phi ^{(1,0)}(p))`$.
###### Proof.
Evidently $`_\phi `$ is generated by $``$ and the graph of $`\beta `$. As mentioned before, if $`(x_1,x_2)`$ is not contained in the $`_\phi `$-orbit of $`p=(p_1,p_2)`$, then its $``$-orbit agrees with its $`^{}`$-orbit. It follows that the $`\stackrel{~}{}`$-orbit of $`(x_1,x_2)`$ agrees with the $`_\phi `$-orbit of it.
Let us consider $`[p]__\phi `$. Notice that it splits into four orbits in $``$, namely the $``$-orbits of $`p`$, $`\phi ^{(0,1)}(p)`$, $`\phi ^{(1,0)}(p)`$ and $`\phi ^{(1,1)}(p)`$. From Lemma 2.5 we can see that these orbits split into eight orbits in $`^{}`$, namely the $`^{}`$-orbits of $`p`$, $`\phi ^{(1,0)}(p)`$, $`\phi ^{(0,1)}(p)`$, $`\phi ^{(0,2)}(p)`$, $`\phi ^{(1,0)}(p)`$, $`\phi ^{(1,1)}(p)`$, $`\phi ^{(1,1)}(p)`$ and $`\phi ^{(2,1)}(p)`$. It can be easily seen that
$$[p]_\stackrel{~}{}=[p]_{^{}}[\phi ^{(1,1)}(p)]_{^{}},$$
$$[\phi ^{(0,1)}(p)]_\stackrel{~}{}=[\phi ^{(0,1)}(p)]_{^{}}[\phi ^{(1,0)}(p)]_{^{}},$$
$$[\phi ^{(1,0)}(p)]_\stackrel{~}{}=[\phi ^{(1,0)}(p)]_{^{}}[\phi ^{(2,1)}(p)]_{^{}}$$
and
$$[\phi ^{(1,1)}(p)]_\stackrel{~}{}=[\phi ^{(1,1)}(p)]_{^{}}[\phi ^{(0,2)}(p)]_{^{}}.$$
Therefore, by glueing the $`\stackrel{~}{}`$-orbits of $`p`$, $`\phi ^{(0,1)}(p)`$, $`\phi ^{(1,0)}(p)`$ and $`\phi ^{(1,1)}(p)`$, we can recover the equivalence relation $`_\phi `$. โ
By \[HPS, Theorem 4.6\], every minimal homeomorphism on the Cantor set is conjugate to a Bratteli-Vershik transformation on a simple properly ordered Bratteli diagram. Hence we can summarize the results obtained in this section as follows.
###### Theorem 2.10.
Let $`(X_1,\phi _1)`$ and $`(X_2,\phi _2)`$ be two Cantor minimal $``$-systems and let $`p_1X_1`$ and $`p_2X_2`$. Take clopen subsets $`U_1X_1`$ and $`U_2X_2`$ so that $`p_i`$ and $`\phi _i(p_i)`$ do not belong to $`U_i`$ for each $`i=1,2`$. Put $`B=(\{p_1\}\times U_2)(U_1\times \{p_2\})`$ and $`B^{}=(\{\phi _1(p_1)\}\times U_2)(U_1\times \{\phi _2(p_2)\})`$. Define $`\beta :BB^{}`$ by $`\beta (p_1,x_2)=(\phi _1(p_1),x_2)`$ and $`\beta (x_1,p_2)=(x_1,\phi _2(p_2))`$. Let $`\phi `$ be the $`^2`$-action on $`X=X_1\times X_2`$ induced by $`\phi _1\times \mathrm{id}`$ and $`\mathrm{id}\times \phi _2`$. Put $`p=(p_1,p_2)`$.
Then we can find a subequivalence relation $`^{}_\phi `$ such that the following are satisfied.
1. $`^{}`$ is a minimal AF equivalence relation, where the topology is given by $`_\phi `$.
2. Both $`B`$ and $`B^{}`$ are closed $`^{}`$-รฉtale thin subsets.
3. $`^{}(B\times B^{})`$ is empty.
4. $`\beta :BB^{}`$ induces an isomorphism between $`^{}(B\times B)`$ and $`^{}(B^{}\times B^{})`$.
5. The equivalence relation $`_\phi `$ is generated by $`^{}`$, the graph of $`\beta `$ and
$$\{(p,\phi ^{(0,1)}(p)),(\phi ^{(0,1)}(p),\phi ^{(1,1)}(p))(\phi ^{(1,1)}(p),\phi ^{(1,0)}(p))\}.$$
## 3 The main result
Let $`(X,\phi )`$ and $`(Y,\psi )`$ be two Cantor minimal $`^2`$-systems and let $`\pi :(Y,\psi )(X,\phi )`$ be a factor map.
###### Lemma 3.1.
Suppose that $``$ is an open subequivalence relation of $`_\phi `$. For
$$๐ฎ=\{(y,y^{})_\psi :(\pi (y),\pi (y^{}))\},$$
we have the following.
1. If $``$ is compact and open, then $`๐ฎ`$ is also compact and open.
2. If $``$ is AF, then $`๐ฎ`$ is also AF.
###### Proof.
(2) follows immediately from (1). Suppose that $``$ is compact and open. Since $`\pi \times \pi :_\psi _\phi `$ is proper and continuous, we can see that $`๐ฎ=_\psi (\pi \times \pi )^1()`$ is compact and open. โ
###### Lemma 3.2.
Suppose that $``$ is an open subequivalence relation of $`_\phi `$. Let $`๐ฎ=_\psi (\pi \times \pi )^1()`$. If $`BX`$ is a closed $``$-รฉtale thin subset, then $`\pi ^1(B)`$ is a closed $`๐ฎ`$-รฉtale thin subset.
###### Proof.
Let $`\mu `$ be an $`๐ฎ`$-invariant probability measure. Then we have
$$\mu (\pi ^1(B))=\pi _{}(\mu )(B)=0,$$
because $`\pi _{}(\mu )`$ is a $``$-invariant probability measure.
Take $`(y,y^{})๐ฎ`$. By the รฉtaleness of $`B`$, we can find a clopen neighborhood $`V`$ of $`(\pi (y),\pi (y^{}))`$ in $``$ such that, for $`(x,x^{})V`$, we have $`xB`$ if and only if $`x^{}B`$. It is clear that $`U=_\psi (\pi \times \pi )^1(V)`$ is a clopen neighborhood of $`(y,y^{})`$ in $`๐ฎ`$. Suppose $`(z,z^{})U`$. Because of $`(\pi (z),\pi (z^{}))V`$, we have
$$z\pi ^1(B)\pi (z)B\pi (z^{})Bz^{}\pi ^1(B).$$
It follows that $`\pi ^1(B)`$ is $`๐ฎ`$-รฉtale. โ
Now we are ready to prove the main theorem.
###### Theorem 3.3.
Let $`\pi :(Y,\psi )(X,\phi )`$ be a factor map between Cantor minimal $`^2`$-systems. If $`(X,\phi )`$ is conjugate to a product of two Cantor minimal $``$-systems, then $`_\psi `$ is affable.
###### Proof.
We may assume that $`(X,\phi )`$ is equal to the product of two Cantor minimal $``$-systems $`(X_1,\phi _1)`$ and $`(X_2,\phi _2)`$, that is, $`X=X_1\times X_2`$ and $`\phi ^{(n,m)}(x_1,x_2)=(\phi _1^n(x_1),\phi _2^m(x_2))`$ for all $`(n,m)^2`$. Let $`p=(p_1,p_2)`$, $`U_1,U_2`$, $`B`$, $`B^{}`$, $`\beta :BB^{}`$ and $`^{}`$ be as in Theorem 2.10.
Put $`๐ฎ=_\psi (\pi \times \pi )^1(^{})`$. Thanks to Theorem 2.10 (1) and Lemma 3.1, the equivalence relation $`๐ฎ`$ is AF, where the topology is given by $`_\psi `$. In order to show that $`๐ฎ`$ is minimal, let us choose $`x_iX_i\{p_i\}`$ and put $`x_0=(x_1,x_1)X`$. Take $`yY`$ arbitrarily. The closure of $`[\pi (y)]_{^{}}`$ is $`X`$, because $`^{}`$ is minimal. It follows that the closure of $`[y]_๐ฎ`$ contains a preimage of $`x_0`$, namely $`y_0Y`$. On account of $`[y_0]__\psi =[y_0]_๐ฎ`$, we can see that $`[y_0]_๐ฎ`$ is dense in $`Y`$. Therefore $`[y]_๐ฎ`$ is dense in $`Y`$.
Put $`C=\pi ^1(B)`$ and $`C^{}=\pi ^1(B^{})`$. By means of Theorem 2.10 (2) and Lemma 3.2, we have that both $`C`$ and $`C^{}`$ are closed $`๐ฎ`$-รฉtale thin subsets. Moreover, it is easily seen that $`๐ฎ(C\times C^{})`$ is empty.
We define a homeomorphism $`\gamma :CC^{}`$ as follows. Take $`yC`$. If $`\pi (y)=(p_1,x_2)`$ for some $`x_2U_2`$, then we set $`\gamma (y)=\psi ^{(1,0)}(y)`$. If $`\pi (y)=(x_1,p_2)`$ for some $`x_1U_1`$, then we set $`\gamma (y)=\psi ^{(0,1)}(y)`$. It is routine to check that $`\gamma `$ is a well-defined homeomorphism from $`C`$ to $`C^{}`$ and $`\gamma `$ induces an isomorphism between $`๐ฎ(C\times C)`$ and $`๐ฎ(C^{}\times C^{})`$.
Let $`\stackrel{~}{๐ฎ}`$ be the equivalence relation generated by $`๐ฎ`$ and the graph of $`\gamma `$. We can apply the absorption theorem \[GPS2, Theorem 4.18\] to $`๐ฎ`$ and $`\gamma :CC^{}`$ and get that $`\stackrel{~}{๐ฎ}`$ is affable.
The equivalence relation $`\stackrel{~}{๐ฎ}`$ is a little smaller than $`_\psi `$. We resolve this problem by using the absorption theorem three more times. Let $`D_1=\pi ^1(p)`$, $`D_2=\pi ^1(\phi ^{(0,1)}(p))`$, $`D_3=\pi ^1(\phi ^{(1,1)}(p))`$ and $`D_4=\pi ^1(\phi ^{(1,0)}(p))`$. At first, we apply the absorption theorem to $`\psi ^{(0,1)}:D_1D_2`$. Notice that
$$\stackrel{~}{๐ฎ}(D_i\times D_i)=\{(y,y):yD_i\}$$
for each $`i=1,2`$ and that $`\stackrel{~}{๐ฎ}(D_1\times D_2)`$ is empty. Therefore the hypothesis of the absorption theorem is trivially satisfied. It follows that the equivalence relation generated by $`\stackrel{~}{๐ฎ}`$ and
$$\{(y,\psi ^{(0,1)}(y)):yD_1\}$$
is affable. Theorem 2.10 (5) and two more applications of the absorption theorem imply that $`_\psi `$ is affable. โ
e-mail: matui@math.s.chiba-u.ac.jp
Graduate School of Science and Technology,
Chiba University,
1-33 Yayoi-cho, Inage-ku,
Chiba 263-8522,
Japan. |
warning/0506/math0506088.html | ar5iv | text | # Standard monomial bases & geometric consequences for certain rings of invariants
## Introduction
In , DeConcini-Procesi constructed a characteristic-free basis for the ring of invariants appearing in classical invariant theory (cf. ) for the action of the general linear, symplectic and orthogonal groups. In , the authors also considered the $`SL_n(K)`$-action on $`X=\underset{m\text{ copies}}{\underset{}{V\mathrm{}V}}\underset{q\text{ copies}}{\underset{}{V^{}\mathrm{}V^{}}}`$, $`V=K^n,K`$ an algebraically closed field of arbitrary characteristic and $`m,q>n`$, and described a set of algebra generators for $`K[X]^{SL_n(K)}`$.
The main goal of this paper is to prove the Cohen-Macaulayness of $`K[X]^{SL_n(K)}`$ (note that the Cohen-Macaulayness of $`K[X]^{GL_n(K)}`$ follows from the fact that
$`Spec(K[X]^{GL_n(K)})`$ is a certain determinantal variety inside $`M_{m,q}`$, the space of $`m\times q`$ matrices; note also that in characteristic $`0`$, the Cohen-Macaulayness of $`K[X]^{SL_n(K)}`$ follows from ). In recent times, among the several techniques of proving the Cohen-Macaulayness of algebraic varieties, two techniques have proven to be quite effective, namely, Frobenius-splitting technique and deformation technique. Frobenius-splitting technique is used in , for example, for proving the (arithmetic) Cohen-Macaulayness of Schubert varieties. Frobenius-splitting technique is also used in for proving the Cohen-Macaulayness of certain varieties. The deformation technique consists in constructing a flat family over $`๐ธ^1`$, with the given variety as the generic fiber (corresponding to $`tK`$ invertible). If the special fiber (corresponding to $`t=0`$) is Cohen-Macaulay, then one may conclude the Cohen-Macaulayness of the given variety. Hodge algebras (cf. ) are typical examples where the deformation technique affords itself very well. Deformation technique is also used in . The philosophy behind these works is that if there is a โstandard monomial basisโ for the co-ordinate ring of the given variety, then the deformation technique will work well in general (using the โstraightening relationsโ). It is this philosophy that we adopt in this paper in proving the Cohen-Macaulayness of $`K[X]^{SL_n(K)}`$. To be more precise, the proof of the Cohen-Macaulayness of $`K[X]^{SL_n(K)}`$ is accomplished in the following steps:
$``$ We first construct a $`K`$-subalgebra $`S`$ of $`K[X]^{SL_n(K)}`$ by prescribing a set of algebra generators $`\{f_\alpha ,\alpha H\},H`$ being a finite partially ordered set and $`f_\alpha K[X]^{SL_n(K)}`$.
$``$ We construct a โstandard monomialโ basis for $`S`$ by
(i) defining โstandard monomialsโ in the $`f_\alpha `$โs (cf. Definition 4.0.1)
(ii) writing down the straightening relation for a non-standard (degree $`2`$) monomial $`f_\alpha f_\beta `$ (cf. Theorem 4.1.1)
(iii) proving linear independence of standard monomials (by relating the generators of $`S`$ to certain determinantal varieties) (cf. ยง4.2)
(iv) proving the generation of $`S`$ (as a vector space) by standard monomials (using (ii)). In fact, to prove the generation for $`S`$, we first prove generation for a โgraded versionโ $`R(D)`$ of $`S`$, where $`D`$ is a distributive lattice obtained by adjoining $`\mathrm{๐},\mathrm{๐}`$ (the largest and the smallest elements of $`D`$) to $`H`$. We then deduce the generation for $`S`$. In fact, we construct a โstandard monomialโ basis for $`R(D)`$. While the generation by standard monomials for $`S`$ is deduced from the generation by standard monomials for $`R(D)`$, the linear independence of standard monomials in $`R(D)`$ is deduced from the linear independence of standard monomials in $`S`$ (cf. (iii) above).
$``$ We give a presentation for $`S`$ as a $`K`$-algebra (cf. Theorem 4.5.5)
$``$ We prove the normality and Cohen-Macaulayness of $`R(D)`$ by showing that Spec$`R(D)`$ flatly degenerates to the toric variety associated to the distributive lattice $`D`$ (cf. Theorem 5.4.3).
$``$ We deduce the normality and Cohen-Macaulayness of $`S`$ from the normality and Cohen-Macaulayness of $`R(D)`$ (cf. Theorem 5.4.4).
$``$ Using the normality of $`S`$ and a crucial Lemma concerning GIT (cf. Lemma 2.0.4 which gives a set of sufficient conditions for a $`\underset{ยฏ}{\mathrm{normal}}`$ sub algebra of $`K[X]^{SL_n(K)}`$
to equal $`K[X]^{SL_n(K)}`$), we show that $`S`$ is in fact $`K[X]^{SL_n(K)}`$, and hence conclude that $`K[X]^{SL_n(K)}`$ is Cohen-Macaulay.
As a consequence, we present (Theorem 6.0.2)
$``$ First fundamental Theorem for $`SL_n(K)`$-invariants, i.e., describing algebra generators for $`K[X]^{SL_n(K)}`$.
$``$ Second fundamental Theorem for $`SL_n(K)`$-invariants, i.e., describing generators for the ideal of relations among these algebra generators for $`K[X]^{SL_n(K)}`$.
$``$ A standard monomial basis for $`K[X]^{SL_n(K)}`$
As a by-product of our main results, we recover Theorem 3.3 of (which describes a set of algebra generators for $`K[X]^{SL_n(K)}`$). It should be pointed out that in , the authors remark (cf. , Remark (ii) following Theorem 3.3) โWe have in fact explicit bases for the rings $`K[X]^{SL_n(K)},K[X]^{GL_n(K)}`$โ. Of course, combining Theorems 1.2 & 3.1 of , one does obtain a basis for $`K[X]^{GL_n(K)}`$; nevertheless, there are no details given in regarding the basis for $`K[X]^{SL_n(K)}`$ (probably, the authors had in their minds the same basis for $`K[X]^{SL_n(K)}`$ as the one constructed in this paper). Our main goal in this paper is to prove the Cohen-Macaulayness of $`K[X]^{SL_n(K)}`$; as mentioned above, this is accomplished by first constructing a โstandard monomialโ basis for the subalgebra $`S`$ of $`K[X]^{SL_n(K)}`$, deducing Cohen-Macaulayness of $`S`$, and then proving that $`S`$ in fact equals $`K[X]^{SL_n(K)}`$. Thus we $`\underset{ยฏ}{\mathrm{do}\mathrm{not}}`$ use the results of (especially, Theorem 3.3 of ), we rather give a different proof of Theorem 3.3 of . Further, using Lemma 2.0.4, we get a GIT-theoretic proof (cf.) of the first and second fundamental theorems for the $`GL_n(K)`$-action in arbitrary characteristics which we have included in ยง2.2. (The GIT-theoretic proof as it appears in calls for a mild modification. Further, for the discussions in ยง3 we need the results on the ring of invariants for the $`GL_n(K)`$-action - specifically, first and second fundamental theorems for the $`GL_n(K)`$-action.)
The sections are organized as follows. In ยง1, after recalling some results (pertaining to standard monomial basis) for Schubert varieties (in the Grassmannian) and determinantal varieties, we derive the straightening relations for certain degree $`2`$ non-standard monomials. In ยง2, we first derive some lemmas concerning quotients leading to the main Lemma 2.0.4; we then give a GIT-theoretic proof of the first and second fundamental theorems for the $`GL_n(K)`$-action in arbitrary characteristics. In ยง3, we define the algebra $`S`$. In ยง4, we construct a standard monomial basis for $`S`$; we also introduce the algebra $`R(D)`$, and construct a standard monomial basis for $`R(D)`$. In ยง5, we first prove the normality and Cohen-Macaulayness of $`R(D)`$, and then deduce the normality and Cohen-Macaulayness of $`S`$. In ยง6, we show that $`S`$ is in fact $`K[X]^{SL_n(K)}`$ (using the crucial Lemma 2.0.4) and deduce the Cohen-Macaulayness of $`K[X]^{SL_n(K)}`$; we also present the first and second fundamental theorems for $`SL_n(K)`$-actions.
We thank C.S. Seshadri for many useful discussions (especially, pertaining to ยง2, ยง6).
## 1. Preliminaries
In this section, we recollect some basic results on determinantal varieties, mainly the standard monomial basis for the co-ordinate rings of determinantal varieties in terms of double standard tableaux. Since the results of ยง4 rely on an explicit description of the straightening relations (of a degree $`2`$ non-standard monomial) on a determinantal variety, in this section we derive such straightening relations (cf. Proposition 1.6.3) by relating determinantal varieties to Schubert varieties in the Grassmannian. We first recall some results on Schubert varieties in the Grassmannian, mainly the standard monomial basis for the homogeneous co-ordinate rings (for the Plรผcker embedding) of Schubert varieties. We then recall results for determinantal varieties (by identifying them as open subsets of suitable Schubert varieties in suitable Grassmannians). We then derive the desired straightening relations.
### 1.1. The Grassmannian Variety $`G_{d,n}`$
Let us fix the integers $`1d<n`$ and let $`V=K^n`$, $`K`$ being the base field which we suppose to be algebraically closed of arbitrary characteristic. Let $`G_{d,n}`$ be the Grassmannian variety consisting of $`d`$-dimensional subspaces of $`V`$.
Let $`\rho _d:G_{d,n}(^dV)`$ be the Plรผcker embedding.
Let $`I(d,n):=\{\underset{ยฏ}{i}=(i_1,\mathrm{},i_d)|1i_1<\mathrm{}<i_dn\}`$. We have a partial order $``$ on $`I(d,n)`$, namely, $`\underset{ยฏ}{i}\underset{ยฏ}{j}i_tj_t,t`$. Let $`N=\mathrm{\#}I(d,n)`$ (note that $`N=\left(\genfrac{}{}{0pt}{}{n}{d}\right)`$); we shall denote the projective coordinates of $`(^dV)`$ as $`p_{\underset{ยฏ}{i}},\underset{ยฏ}{i}I(d,n)`$, and refer to them as the Plรผcker coordinates.
For $`1tn`$, let $`V_t`$ be the subspace of $`V`$ spanned by $`\{e_1,\mathrm{},e_t\}`$. For $`\underset{ยฏ}{i}I(d,n)`$, let $`X(\underset{ยฏ}{i})`$ be the Schubert variety associated to $`\underset{ยฏ}{i}`$:
$$X(\underset{ยฏ}{i})=\{UG_{d,n}dim(UV_{i_t})t,1td\}.$$
###### Remark 1.1.1.
Note that under the set-theoretic bijection between the set of Schubert varieties and the set $`I(d,n)`$, the partial order on the set of Schubert varieties given by inclusion induces the partial order $``$ on $`I(d,n)`$.
Let $`R`$ be the homogeneous co-ordinate ring of $`G_{d,n}`$ for the Plรผcker embedding, and for $`wI(d,n)`$, let $`R(w)`$ be the homogeneous co-ordinate ring of the Schubert variety $`X(w)`$.
###### Definition 1.1.2.
A monomial $`f=p_{\tau _1}\mathrm{}p_{\tau _m}`$ is said to be standard if
(\*)
$$\tau _1\mathrm{}\tau _m.$$
Such a monomial is said to be standard on $`X(w)`$, if in addition to condition (\*), we have $`w\tau _1`$.
We recall the following fundamental result: ( cf. ; see also )
###### Theorem 1.1.3.
Standard monomials on $`X(w)`$ of degree $`m`$ give a basis for $`R(w)_m`$.
As a consequence, we have a qualitative description of a typical quadratic relation on a Schubert variety $`X(w)`$ as given by the following Proposition. First one definition:
###### Definition 1.1.4.
Given $`\tau ,\varphi I(d,n)`$, say, $`\tau =(a_1,\mathrm{},a_d),\varphi =(b_1,\mathrm{},b_d)`$, $`\tau \varphi :=(c_1,\mathrm{},c_d),\tau \varphi :=(e_1,\mathrm{},e_d)`$, where $`c_i=\mathrm{max}\{a_i,b_i\},e_i=\mathrm{min}\{a_i,b_i\},i`$ are called the *join* and *meet* of $`\tau `$ and $`\varphi `$ respectively. Note that $`\tau \varphi `$ (resp. $`\tau \varphi `$) is the smallest (resp. largest) element of $`I(d,n)`$ which is $`>`$ (resp. $`<`$) both $`\tau `$ and $`\varphi `$.
###### Proposition 1.1.5.
Let $`w,\tau ,\varphi I(d,n),w>\tau ,\varphi `$. Further let $`\tau ,\varphi `$ be non-comparable (so that $`p_\tau p_\varphi `$ is a non-standard degree $`2`$ monomial on $`X(w)`$). Let
(\*)
$$p_\tau p_\varphi =\underset{i}{}c_ip_{\alpha _i}p_{\beta _i},c_iK^{}$$
be the expression for $`p_\tau p_\varphi `$ as a sum of standard monomials on $`X(w)`$. Then
1. for every $`(\alpha ,\beta )`$ on the R.H.S. we have, $`\alpha >`$ both $`\tau \text{ and }\varphi `$, $`\beta <`$ both $`\tau \text{ and }\varphi `$.
2. for every $`(\alpha ,\beta )`$ on the right-hand side of (\*), we have $`\tau \dot{}\varphi =\alpha \dot{}\beta `$ (here $`\dot{}`$ denotes a disjoint union)
3. the term $`p_{\tau \varphi }p_{\tau \varphi }`$ occurs on the right-hand side of (\*) with coefficient $`1`$.
Such a relation as in (\*) is called a straightening relation.
###### Proof.
(1): Pick a minimal element in $`\{\alpha _i\}`$, call it $`\alpha _1`$. Restrict (\*) to $`X(\alpha _1)`$. Then R.H.S. is a non-zero sum of standard monomials on $`X(\alpha _1)`$. Hence linear independence of standard monomials on $`X(\alpha _1)`$ implies that the restriction of L.H.S. to $`X(\alpha _1)`$ is non-zero. Hence it follows that $`\alpha _1`$ both $`\tau \text{ and }\varphi `$ (note that restriction of $`p_\theta `$ to $`X(\alpha _1)`$ is non-zero if and only if $`\alpha _1\theta `$); we have in fact $`\alpha >\tau ,\varphi `$, for, if $`\alpha `$ equals one of $`\{\tau ,\varphi \}`$, say $`\alpha =\tau `$, then $`p_\tau p_\varphi =p_\alpha p_\varphi `$ would be standard, a contradiction. The assertion on $`\alpha `$ follows from this. The assertion on $`\beta `$ is proved similarly by working with $`w_0\tau ,w_0\varphi `$ (in the place of $`\tau ,\varphi `$), $`w_0`$ being the element of largest length in the Weyl group.
(2) follows from weight considerations (note that $`p_\tau ,\tau I(d,n)`$ \- say, $`\tau =(a_1,\mathrm{},a_d)`$ \- is a weight vector (for the $`T`$-action, $`T`$ being the maximal torus of diagonal matrices in $`GL_n(K)`$) of weight $`(ฯต_{a_1}+\mathrm{}+ฯต_{a_d})`$).
For a proof of (3), refer to , Proposition 7.33. โ
A presentation for $`R(w)`$. Let $`Z_w=\{\tau I(d,n)|w\tau \}`$. Consider the polynomial algebra $`K[x_\tau ,\tau Z_w]`$. For a pair $`\tau ,\varphi `$ in $`Z_w`$ such that $`\tau ,\varphi `$ are not comparable, denote $`F_{\tau ,\varphi }=x_\tau x_\varphi _ic_ix_{\alpha _i}x_{\beta _i},\alpha _i,\beta _i,c_i`$ being as in Proposition 1.1.5. Let $`I_w`$ be the ideal in $`K[x_\tau ,\tau Z_w]`$ generated by $`\{F_{\tau ,\varphi },\tau ,\varphi \mathrm{non}\mathrm{comparable}\}`$. Consider the surjective map $`f_w:K[x_\tau ,\tau Z_w]R(w),x_\tau p_\tau `$. We have
###### Proposition 1.1.6.
$`f_w`$ induces an isomorphism $`K[x_\tau ,\tau Z_w]/I_wR(w)`$.
See for a proof.
### 1.2. The opposite big cell in $`G_{d,n}`$
Let $`P_d`$ be the parabolic subgroup of $`G(=GL_n(K))`$ consisting of all matrices of the form
$$\left(\begin{array}{cc}& \\ 0& \end{array}\right),$$
where the $`0`$-matrix is of size $`nd\times d`$. Then we have an identification $`\phi _d:G/P_dG_{d,n}`$. Denote by $`O^{}`$ the sub group of $`G`$ consisting of matrices of the form
$$\left(\begin{array}{cc}I_d& 0_{d\times (nd)}\\ A_{(nd)\times d}& I_{nd}\end{array}\right)$$
where $`I_d`$ (resp. $`I_{nd}`$) is the $`d\times d`$ (resp. $`nd\times nd`$) identity matrix. We have that the restriction of the canonical morphism $`\theta _d:GG/P_d`$ to $`O^{}`$ is an open immersion, and $`\theta _d(O^{})=B^{}e_{id}`$, where $`e_{id}`$ is the coset $`P_d`$ of $`G/P_d`$, and $`B^{}`$ is the Borel sub group of lower triangular matrices in $`G`$; also, $`\phi _d(B^{}e_{id})`$ is the opposite big cell in $`G_{d,n}`$. Thus the opposite big cell in $`G_{d,n}`$ gets identified with $`O^{}`$, and in the sequel we shall denote the opposite big cell by just $`O^{}`$. Note that $`O^{}๐ธ^{d(nd)}`$.
### 1.3. The functions $`f_{\underset{ยฏ}{j}}`$ on $`O^{}`$:
Consider the morphism $`\xi _d:G(^dV)`$, where $`\xi _d=\rho _d\phi _d\theta _d`$, $`\rho _d,\phi _d,\theta _d`$ being as above. Then $`p_{\underset{ยฏ}{j}}(\xi _d(g))`$ is simply the $`d`$-minor of $`g`$ consisting of the first $`d`$ columns and rows given by $`j_1,\mathrm{},j_d`$. For $`\underset{ยฏ}{j}I(d,n)`$, we shall denote by $`f_{\underset{ยฏ}{j}}`$ the restriction of $`p_{\underset{ยฏ}{j}}`$ to $`O^{}`$. Under the identification
$$O^{}=\left\{\left(\begin{array}{c}I_d\\ A\end{array}\right),AM_{nd,d}\right\}$$
we have for $`z=\left(\begin{array}{c}I_d\\ A\end{array}\right)O^{}`$, $`f_{\underset{ยฏ}{j}}(z)`$ is simply a certain minor of $`A`$, which may be explicitly described as follows. Let $`\underset{ยฏ}{j}=(j_1,\mathrm{},j_d)`$, and let $`j_r`$ be the largest entry $`d`$. Let $`\{k_1,\mathrm{},k_{dr}\}`$ be the complement of $`\{j_1,\mathrm{},j_r\}`$ in $`\{1,\mathrm{},d\}`$. Then $`f_{\underset{ยฏ}{j}}(z)`$ is the $`(dr)`$-minor of $`A`$ with column indices $`k_1,\mathrm{}k_{dr}`$, and row indices $`j_{r+1},\mathrm{},j_d`$ (here the rows of $`A`$ are indexed as $`d+1,\mathrm{},n`$). Conversely, given a minor of $`A`$, say, with column indices $`b_1,\mathrm{},b_s`$, and row indices $`j_{ds+1},\mathrm{},j_d`$ (again, the rows of $`A`$ are indexed as $`d+1,\mathrm{},n`$), it is $`f_{\underset{ยฏ}{j}}(z)`$, where $`\underset{ยฏ}{j}=(j_1,\mathrm{},j_d)`$ is given as follows: $`\{j_1,\mathrm{},j_{ds}\}`$ is the complement of $`\{b_1,\mathrm{},b_s\}`$ in $`\{1,\mathrm{},d\}`$, and $`j_{ds+1},\mathrm{},j_d`$ are simply the row indices.
Convention. If $`\underset{ยฏ}{j}=(1,\mathrm{},d)`$, then $`f_{\underset{ยฏ}{j}}`$ evaluated at $`z`$ is $`1`$; we shall make it correspond to the minor of $`A`$ with row indices (and column indices) given by the empty set.
### 1.4. The opposite cell in $`X(w)`$
For a Schubert variety $`X(w)`$ in $`G_{d,n}`$, let us denote $`O^{}X(w)`$ by $`Y(w)`$; we refer to $`Y(w)`$ as the *opposite cell in* $`X(w)`$. We consider $`Y(w)`$ as a closed subvariety of $`O^{}`$. In view of Proposition 1.1.6, we obtain that the ideal defining $`Y(w)`$ in $`O^{}`$ is generated by
$$\{f_{\underset{ยฏ}{i}}\underset{ยฏ}{i}I(d,n),w\underset{ยฏ}{i}\}.$$
### 1.5. Determinantal Varieties
Let $`Z=M_{r,d}(K)`$, the space of all $`r\times d`$ matrices with entries in $`K`$. We shall identify $`Z`$ with $`๐ธ^{rd}`$. We have $`K[Z]=K[x_{i,j},1ir,1jd]`$.
The variety $`D_t`$. Let $`X=(x_{ij})`$, $`1ir`$, $`1jd`$ be a $`r\times d`$ matrix of indeterminates. Let $`A\{1,\mathrm{},r\},B\{1,\mathrm{},d\},\mathrm{\#}A=\mathrm{\#}B=s`$, where $`s\text{ min }\{r,d\}`$. We shall denote by $`p(A,B)`$ the $`s`$-minor of $`X`$ with row indices given by $`A`$, and column indices given by $`B`$. For $`t,1t\text{ min }\{r,d\}`$, let $`I_t(X)`$ be the ideal in $`K[x_{i,j}]`$ generated by $`\{p(A,B),A\{1,\mathrm{},r\},B\{1,\mathrm{},d\},\mathrm{\#}A=\mathrm{\#}B=t\}`$. Let $`D_t(M_{r,d})`$ (or just $`D_t`$) be the determinantal variety (a closed subvariety of $`Z`$), with $`I_t(X)`$ as the defining ideal. In the discussion below, we also allow $`t=d+1`$ in which case $`D_t=Z`$
Identification of $`D_t`$ with $`Y_\varphi `$. Let $`G=GL_n(K)`$. Let $`r,d`$ be such that $`r+d=n`$. Let $`X`$ be a $`r\times d`$ matrix of indeterminates. As in ยง1.2, let us identify the opposite cell $`O^{}`$ in $`G/P_d(G_{d,n})`$ as
$$O^{}=\left\{\left(\begin{array}{c}I_d\\ X\end{array}\right)\right\}.$$
As seen above (cf. ยง1.3), we have a bijection between $`\{f_{\underset{ยฏ}{i}},\underset{ยฏ}{i}I(d,n)\}`$ and $`\{\text{minors of }X\}`$ (note that as seen in ยง1.3, if $`\underset{ยฏ}{i}=(1,2,\mathrm{},d)`$, then $`f_{\underset{ยฏ}{i}}=`$ the constant function $`1`$ considered as the minor of $`X`$ with row indices (and column indices) given by the empty set).
For example, take $`r=3=d`$. We have,
$$O^{}=\left\{\left(\begin{array}{c}I_3\\ X_{3\times 3}\end{array}\right)\right\}.$$
We have, $`f_{(1,2,4)}=p(\{1\},\{3\}),f_{(2,4,6)}=p(\{1,3\},\{1,3\})`$.
Let $`\varphi `$ be the $`d`$-tuple, $`\varphi =(t,t+1,\mathrm{},d,n+2t,n+3t,\mathrm{},n)`$ (note that $`\varphi `$ consists of the two blocks $`[t,d]`$, $`[n+2t,n]`$ of consecutive integers - here, for $`i<j`$, $`[i,j]`$ denotes the set $`\{i,i+1,\mathrm{},j\}`$). If $`t=d+1`$, then we set $`\varphi =(n+1d,n+2d,\mathrm{},n)`$ (note then that $`Y_\varphi =O^{}(M_{r,d}(K)))`$.
###### Theorem 1.5.1.
(cf.) $`D_tY_\varphi `$.
###### Corollary 1.5.2.
$`K[D_t]R(\varphi )_{(p_{id})}`$, the homogeneous localization of $`R(\varphi )`$ at $`p_{id}`$.
### 1.6. The partially ordered set $`H_{r,d}`$
Let
$$H_{r,d}=\underset{0s\mathrm{min}\{r,d\}}{}I(s,r)\times I(s,d)$$
where our convention is that $`(\mathrm{},\mathrm{})`$ is the element of $`H_{r,d}`$ corresponding to $`s=0`$. We define a partial order $``$ on $`H_{r,d}`$ as follows:
$``$ We declare $`(\mathrm{},\mathrm{})`$ as the largest element of $`H_{r,d}`$.
$``$ For $`(A,B),(A^{},B^{})`$ in $`H_{r,d}`$, say, $`A=\{a_1,\mathrm{},a_s\},B=\{b_1,\mathrm{},b_s\},A^{}=\{a_1^{},\mathrm{},a_s^{}^{}\},B^{}=\{b_1^{},\mathrm{},b_s^{}^{}\}`$ for some $`s,s^{}1`$, we define $`(A,B)(A^{},B^{})`$ if $`ss^{},a_ja_j^{},b_jb_j^{},1js`$.
The bijection $`\theta `$: As above, let $`n=r+d`$. Then $``$ induces a partial order $``$ on the set of all minors of $`X`$, namely, $`p(A,B)p(A^{},B^{})`$ if $`(A,B)(A^{},B^{})`$. Given $`\underset{ยฏ}{i}I(d,n)`$, let $`m`$ be such that $`i_md,i_{m+1}>d`$. Set
$$A_{\underset{ยฏ}{i}}=\{n+1i_d,n+1i_{d1},\mathrm{},n+1i_{m+1}\},$$
$$B_{\underset{ยฏ}{i}}=\text{ the complement of }\{i_1,i_2\mathrm{},i_m\}\text{ in }\{1,2,\mathrm{},d\}.$$
Define $`\theta :I(d,n)\{\text{all minors of }X\}`$ by setting $`\theta (\underset{ยฏ}{i})=p(A_{\underset{ยฏ}{i}},B_{\underset{ยฏ}{i}})`$ (here, the constant function $`1`$ is considered as the minor of $`X`$ with row indices (and column indices) given by the empty set). Clearly $`\theta `$ is a bijection. Note that $`\theta `$ reverses the respective partial orders, i.e., given $`\underset{ยฏ}{i},\underset{ยฏ}{i}^{}I(d,n)`$, we have, $`\underset{ยฏ}{i}\underset{ยฏ}{i}^{}\theta (\underset{ยฏ}{i})\theta (\underset{ยฏ}{i}^{})`$. Using the partial order $``$, we define standard monomials in $`p(A,B)`$โs:
###### Definition 1.6.1.
A monomial $`p(A_1,B_1)\mathrm{}p(A_s,B_s),s`$ is standard if $`p(A_1,B_1)\mathrm{}p(A_s,B_s)`$.
In view of Theorem 1.1.3, Theorem 1.5.1, and ยง1.4, we obtain
###### Theorem 1.6.2.
Standard monomials in $`p(A,B)`$โs with # $`At1`$ form a basis for $`K[D_t]`$, the algebra of regular functions on $`D_t`$.
As a direct consequence of Proposition 1.1.5, we obtain
###### Proposition 1.6.3.
Let $`p(A_1,A_2),p(B_1,B_2)`$ (in $`K[D_t]`$) be not comparable. Let
$`()`$
$$p(A_1,A_2)p(B_1,B_2)=a_ip(C_{i1},C_{i2})p(D_{i1},D_{i2}),a_iK^{}$$
be the straightening relation in $`K[D_t]`$. Then for every $`i`$, $`C_{i1},C_{i2},D_{i1},D_{i2}`$ have cardinalities $`t1`$; further,
1. $`C_{i1}`$ both $`A_1`$ and $`B_1`$; $`D_{i1}`$ both $`A_1`$ and $`B_1`$.
2. $`C_{i2}`$ both $`A_2`$ and $`B_2`$; $`D_{i2}`$ both $`A_2`$ and $`B_2`$.
3. The term $`p((A_1,A_2)(B_1,B_2))p((A_1,A_2)(B_1,B_2))`$ occurs in (\*) with coefficient $`1`$.
Note that via the bijection $`\theta `$ (defined as above), join and meet (cf. Definition 1.1.4) of two non-comparable elements $`(A_1,A_2),(B_1,B_2)`$ of $`H_{r,d}`$ exist, and in fact are given by $`(A_1,A_2)(B_1,B_2)=(A_1B_1,A_2B_2),(A_1,A_2)(B_1,B_2)=(A_1B_1,A_2B_2)`$.
###### Remark 1.6.4.
On the R.H.S. of (\*), $`C_{i1},C_{i2}`$ could both be the empty set (in which case $`p(C_{i1},C_{i2})`$ is understood as $`1`$). For example, with $`X`$ being a $`2\times 2`$ matrix of indeterminates, we have
$$p_{1,2}p_{2,1}=p_{2,2}p_{1,1}p_\mathrm{},\mathrm{}p_{12,12}$$
###### Remark 1.6.5.
In the sequel, while writing a straightening relation as in Proposition 1.6.3, if for some $`i`$, $`C_{i1},C_{i2}`$ are both the empty set, we keep the corresponding $`p(C_{i1},C_{i2})`$ on the right hand side of the straightening relation (even though its value is $`1`$) in order to have homogeneity in the relation.
Taking $`t=d+1`$ (in which case $`D_t=Z=M_{r,d}(K)`$) in Theorem 1.6.2 and Proposition 1.6.3, we obtain
###### Theorem 1.6.6.
1. Standard monomials in $`p(A,B)`$โs form a basis for $`K[Z](K[x_{ij},1ir,1jd])`$.
2. Relations similar to those in Proposition 1.6.3 hold on $`Z`$.
###### Remark 1.6.7.
Note that Theorem 1.6.6,(1) recovers the result of Doubleit-Rota-Stein (cf. , Theorem 2):
###### Remark 1.6.8.
Theorem 1.6.2 is also proved in (Theorem 1.2 in ). But we had taken the above approach of deducing Theorem 1.6.2 from Theorems 1.1.3, 1.5.1 in order to derive the straightening relations as given by Proposition 1.6.3(which are crucial for the discussion in ยง4).
A presentation for $`K[(D_t)]`$. Let $`Z_t=\{(A,B)(\mathrm{},\mathrm{}),(A,B)H_{r,d},\mathrm{\#}A=\mathrm{\#}Bt1\}`$.
Consider the polynomial algebra $`K[x(A,B),(A,B)H_{r,d},\mathrm{\#}A=\mathrm{\#}Bt1]`$. For two non-comparable pairs (under $``$ (cf. ยง1.6)) $`(A_1,A_2),(B_1,B_2)`$ in $`Z_t`$, denote
$`F((A_1,A_2);(B_1,B_2))=x(A_1,A_2)(B_1,B_2)a_ix(C_{i1},C_{i2})x(D_{i1},D_{i2})`$, where
$`C_{i1},C_{i2},D_{i1},D_{i2},a_i`$ are as in Proposition 1.6.3. Let $`I_t`$ be the ideal generated by
$$\{F((A_1,A_2);(B_1,B_2)),(A_1,A_2),(B_1,B_2)\mathrm{non}\mathrm{comparable}\}$$
Consider the surjective map $`f_t:K[x(A,B),(A,B)Z_t]K[D_t],x(A,B)p(A,B)`$. Then in view of Proposition 1.1.6 and Theorem 1.5.1, we obtain
###### Proposition 1.6.9.
(A presentation for $`K[D_t]`$) $`f_t`$ induces an isomorphism
$`K[x(A,B),(A,B)Z_t]/I_tK[D_t]`$.
## 2. $`GL_n(K)`$-action
In this section, we first prove some Lemmas concerning quotients, to be applied to the following situation:
Suppose, we have an action of a reductive group $`G`$ on an affine variety $`X=SpecR`$. Suppose that $`S`$ is a subalgebra of $`R^G`$. We give below (cf. Lemma 2.0.4) a set of sufficient conditions for the equality $`S=R^G`$. We start with recalling
###### Theorem 2.0.1.
(Zariski Main Theorem, ,III.9) Let $`\phi :XY`$ be a morphism such that
1. $`\phi `$ is surjective
2. fibers of $`\phi `$ are finite
3. $`\phi `$ is birational
4. $`Y`$ is normal
Then $`\phi `$ is an isomorphism.
Let $`X=\mathrm{Spec}R`$ and a reductive group $`G`$ act linearly on $`X`$, i.e., we have a linear action of $`G`$ on an affine space $`๐ธ^r`$ and we have a $`G`$-equivariant closed immersion $`X๐ธ^r`$. Further, let $`R`$ be a graded $`K`$-algebra. Let $`X^{ss}`$ be the set of semi-stable points of $`X`$ (i.e., points $`x`$ such that $`0\overline{Gx}`$). Let $`X_1=ProjR,X_1^{ss}`$, the set of semi-stable points of $`X_1`$ (i.e., points $`yX_1`$ such that if $`\widehat{x}`$ is any point in $`K^{n+1}\mathrm{\hspace{0.17em}0}`$ lying over $`y`$, then $`\widehat{x}`$ is in $`X^{ss}`$). Let $`f_1,\mathrm{},f_N`$ be homogeneous $`G`$-invariant elements in $`R`$. Let $`S=K[f_1,\mathrm{},f_N]`$. Then for the morphism $`SpecR^GSpecS`$, the hypothesis (2) in Theorem 2.0.1 may be concluded if $`\{f_1,\mathrm{},f_N\}`$ is base-point free on $`X_1^{ss}`$ as given by the following
###### Lemma 2.0.2.
Suppose $`f_1,\mathrm{},f_N`$ are homogeneous $`G`$-invariant elements in $`R`$ such that for any $`xX^{ss}`$ , $`f_i(x)0`$, for at least one $`i`$. Then $`SpecR^GSpecS`$ has finite fibers.
###### Proof.
Case 1: Let $`f_1,\mathrm{},f_N`$ be of the same degree, say, $`d`$. Let $`Y=SpecR^G(=X//G`$, the categorical quotient) and $`\phi :XY`$ be the canonical quotient map. Let $`X_1=ProjR,X_1^{ss}`$ the set of semi-stable points of $`X_1`$. Let $`Y_1=ProjR^G(=X_1^{ss}//G)`$, and $`\phi _1:X_1^{ss}Y_1`$ be the canonical quotient map. Consider $`\psi :X๐ธ^N,x(f_1(x),\mathrm{},f_N(x))`$. This induces a map $`\rho :Y๐ธ^N`$ (since $`f_1,\mathrm{},f_N`$ are $`G`$-invariant). The commutative diagram
| $`X`$ | | |
| --- | --- | --- |
| $`\phi `$ ? | QQQQQQQQs $`\psi `$ |
| $`Y`$ | $`\underset{\rho }{\text{ }\text{-}}`$ | $`๐ธ^N`$ |
induces the commutative diagram
| $`X_1^{\text{ss}}`$ | | |
| --- | --- | --- |
| $`\phi _1`$ ? | QQQQQQQQs $`\psi _1`$ |
| $`Y_1`$ | $`\underset{\rho _1}{\text{ }\text{-}}`$ | $`^{N1}`$ |
Note that $`\psi _1:X_1^{ss}^{N1}`$ is defined in view of the hypothesis that for any $`xX^{ss}`$, $`f_i(x)0`$, for at least one $`i`$. Note also that $`f_1,\mathrm{},f_N`$ are sections of the ample line bundle $`๐ช_{X_1}(d)`$ as well as the basic fact from GIT that this line bundle descends to an ample line bundle on $`Y_1`$, which we denote by $`๐ช_{Y_1}(d)`$.
Claim 1: $`\rho _1`$ is a finite morphism.
Proof of Claim 1: Since $`f_1,\mathrm{},f_N`$ are $`G`$-invariant, we get that $`f_iH^0(Y_1,๐ช_{Y_1}(d))`$. Hence we obtain that
$$\rho _1^{}(๐ช_{^{N1}}(1))=๐ช_{Y_1}(d)$$
Thus, $`\rho _1^{}(๐ช_{^{N1}}(1))`$ is ample, and hence $`\rho _1`$ is finite (over any fiber $`(\rho _1)_z,z^{N1}`$,
$`\rho _1^{}(๐ช_{^{N1}}(1))|_{(\rho _1)_z}`$ is both ample and trivial, and hence dim$`(\rho _1)_z`$ is zero), and Claim 1 follows.
Claim 2: $`\rho `$ is a finite morphism.
Proof of Claim 2: Let $`S^{}=R^G`$. Let $`S^{(d)}=_nS_{nd}^{}`$. We have $`^{N1}=\mathrm{Proj}K[x_1,\mathrm{},x_N]`$. Since $`\rho _1`$ is finite, $`๐ช_{Y_1}`$ is a coherent $`๐ช_{^{N1}}`$-module.
We see that
$$H^0(^{N1},๐ช_{Y_1}๐ช_{^{N1}}(n))H^0(Y_1,\rho _1^{}(๐ช_{^{N1}}(n))$$
since the direct image of $`\rho _1^{}(๐ช_{^{N1}}(n))`$ by $`\rho _1`$ is $`๐ช_{Y_1}๐ช_{^{N1}}(n)`$ and $`\rho _1`$ is a finite morphism. On the other hand we have
$$\rho _1^{}(๐ช_{^{N1}}(n))๐ช_{Y_1}(nd).$$
Thus we have
$$H^0(^{N1},๐ช_{Y_1}๐ช_{^{N1}}(n))H^0(Y_1,๐ช_{Y_1}(nd))S_{nd}^{}.$$
Thus the graded $`K[x_1,\mathrm{},x_N]`$-module associated to the coherent sheaf $`๐ช_{Y_1}`$ on $`^{N1}`$ is $`S^{(d)}`$ and by the basic theorems of Serre, $`S^{(d)}`$ is of finite type over $`K[x_1,\mathrm{},x_N]`$. Now a $`d`$-th power of any homogeneous element of $`S^{}`$ is in $`S^{(d)}`$ and thus $`S^{}`$ is integral over $`K[x_1,\mathrm{},x_N]`$, which proves that $`\rho `$ is finite. Claim 2 and hence the required result follows from this.
Case 2: Let $`f_1,\mathrm{},f_N`$ be homogeneous possibly of different degrees, say, deg$`f_i=d_i`$. Let $`d=l.c.m.\{d_i\},e_i=\frac{d}{d_i}`$. Set $`g_i=f_i^{e_i},1iN`$. Then $`\{g_1,\mathrm{},g_N\}`$ is again base-point free on $`(ProjR)^{ss}`$. Hence by Case 1, we have that
$`R^G`$ is a finite $`K[g_1,\mathrm{},g_N]`$-module, and hence a finite $`K[f_1,\mathrm{},f_N]`$-module
(note that $`K[g_1,\mathrm{},g_N]K[f_1,\mathrm{},f_N]R^G`$).
In the Lemma below, we describe a set of sufficient conditions for (3) of Lemma 2.0.1, namely, birationality.
###### Lemma 2.0.3.
Suppose $`F:XY`$ is a surjective morphism of (irreducible) algebraic varieties, and $`U`$ is an open subset of $`X`$ such that
1. $`F|_U:UY`$ is an immersion
2. dim$`U`$ = dim$`Y`$.
Then $`F`$ is birational.
###### Proof.
Hypothesis (1) implies that $`F(U)`$ is locally closed in $`Y`$. This fact together with Hypothesis (2) implies that $`F(U)`$ is open in $`Y`$, and the result follows. โ
We now return to the situation of a linear action of a reductive group $`G`$ on an affine variety $`X=SpecR`$ with $`R`$ a graded $`K`$-algebra. Let $`f_1,\mathrm{},f_N`$ be homogeneous $`G`$-invariant elements in $`R`$. Let $`S=K[f_1,\mathrm{},f_N]`$. Combining Lemmas 2.0.1, 2.0.2, 2.0.3, we arrive at the following Lemma which gives a set of sufficient conditions for the equality $`S=R^G`$. Before stating the lemma, let us observe the following. Suppose that $`U`$ is a non-empty $`G`$-stable open subset in $`X`$. Since $`\phi :XSpecR^G`$ is surjective, $`\phi (U)`$ contains a non-empty open subset. Hence by shrinking $`U`$, if necessary, we can suppose that $`\phi (U)`$ is open. We suppose that this is the case and denote it by $`U//G`$.
###### Lemma 2.0.4.
Let notation be as above. Let $`\psi :X๐ธ^N`$ be the map, $`x(f_1(x),\mathrm{},f_N(x))`$. Denote $`D=SpecS`$. Then $`D`$ is the categorical quotient of $`X`$ by $`G`$ and $`\psi :XD`$ is the canonical quotient map, provided the following conditions are satisfied:
(i) For $`xX^{\text{ss}}`$, $`\psi (x)(0)`$.
(ii) There is a $`G`$-stable open subset $`U`$ of $`X`$ such that $`\psi |_{U//G}:U//GD`$ is an immersion.
(iii) $`dimD=dimU//G`$.
(iv) $`D`$ is normal.
###### Remark 2.0.5.
Suppose that $`U`$ is a (non-empty) $`G`$-stable open subset of $`X`$, $`G`$ operates freely with $`U/G`$ as quotient, and $`\psi `$ induces an immersion of $`U/G`$ in $`A^N`$. Then (ii) is satisfied:
This assertion is immediately seen, for we have
$$U/GU//G๐ธ^N$$
and the fact that $`U/GA^N`$ is an immersion implies that $`U//G๐ธ^N`$ immersion.
In the following subsection, using Lemma 2.0.4, we give a GIT-theoretic proof of the first and second fundamental theorems for the $`GL_n(K)`$-action in arbitrary characteristics.
### 2.1. Classical Invariant Theory:
Let $`V=K^n`$, $`X=\underset{m\text{copies}}{\underset{}{V\mathrm{}V}}\underset{q\text{copies}}{\underset{}{V^{}\mathrm{}V^{}}}`$, where $`m,q>n`$.
The $`GL(V)`$-action on $`X`$: Writing $`\underset{ยฏ}{u}=(u_1,u_2,\mathrm{},u_m)`$ with $`u_iV`$ and $`\underset{ยฏ}{\xi }=(\xi _1,\xi _2,\mathrm{},\xi _q)`$ with $`\xi _iV^{}`$, we shall denote the elements of $`X`$ by $`(\underset{ยฏ}{u},\underset{ยฏ}{\xi })`$. The (natural) action of $`GL(V)`$ on $`V`$ induces an action of $`GL(V)`$ on $`V^{}`$, namely, for $`\xi V^{},gGL(V)`$, denoting $`g\xi `$ by $`\xi ^g`$, we have
$$\xi ^g(v)=\xi (g^1v),vV$$
The diagonal action of $`GL(V)`$ on $`X`$ is given by
$$g(\underset{ยฏ}{u},\underset{ยฏ}{\xi })=(g\underset{ยฏ}{u},\underset{ยฏ}{\xi }^g)=(gu_1,gu_2,\mathrm{},gu_m,\xi _1^g,\xi _2^g,\mathrm{},\xi _q^g),gG,(\underset{ยฏ}{u},\underset{ยฏ}{\xi })X$$
The induced action on $`K[X]`$ is given by
$$(gf)(\underset{ยฏ}{u},\underset{ยฏ}{\xi })=f(g^1(\underset{ยฏ}{u},\underset{ยฏ}{\xi })),fK[X],gGL(V)$$
Consider the functions $`\phi _{ij}:XK`$ defined by $`\phi _{ij}((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))=\xi _j(u_i)`$, $`1im,1jq`$. Each $`\phi _{ij}`$ is a $`GL(V)`$-invariant: For $`gGL(V)`$, we have,
$`(g\phi _{ij})((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))`$ $`=`$ $`\phi _{ij}(g^1(\underset{ยฏ}{u},\underset{ยฏ}{\xi }))`$
$`=`$ $`\phi _{ij}((g^1u,\xi ^{g^1})`$
$`=`$ $`\xi _j^{g^1}(g^1u_i)`$
$`=`$ $`\xi _j(u_i)`$
$`=`$ $`\phi _{ij}((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))`$
It is convenient to have a description of the above action in terms of coordinates. So with respect to a fixed basis, we write the elements of $`V`$ as row vectors and those of $`V^{}`$ as column vectors. Thus denoting by $`M_{a,b}`$ the space of $`a\times b`$ matrices with entries in $`K,`$ $`X`$ can be identified with the affine space $`M_{m,n}\times M_{n,q}`$. The action of $`GL_n(K)(=GL(V))`$ on $`X`$ is then given by
$$A(U,W)=(UA,A^1W),AGL_n(K),UM_{m,n},WM_{n,q}$$
And the action of $`GL_n(K)`$ on $`K[X]`$ is given by
$$(Af)(U,W)=f\left(A^1(U,W)\right)=f(UA^1,AW),fK[X]$$
Writing $`U=\left(u_{ij}\right)`$ and $`W=\left(\xi _{kl}\right)`$ we denote the coordinate functions on $`X`$, by $`u_{ij}`$ and $`\xi _{kl}.`$ Further, if $`u_i`$ denotes the $`i`$-th row of $`U`$ and $`\xi _j`$ the $`j`$-th column of $`W,`$ the invariants $`\phi _{ij}`$ described above are nothing but the entries $`u_i,\xi _j(=\xi _j(u_i))`$ of the product $`UW.`$
In the sequel, we shall denote $`\phi _{ij}(\underset{ยฏ}{u},\underset{ยฏ}{\xi })`$ also by $`u_i,\xi _j`$.
The function $`p(A,B)`$: For $`AI(r,m),BI(r,q),1rn`$, let $`p(A,B)`$ be the regular function on $`X`$: $`p(A,B)((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))`$ is the determinant of the $`r\times r`$-matrix $`(u_i,\xi _j)_{iA,jB}`$. Let $`S`$ be the subalgebra of $`R^G`$ generated by $`\{p(A,B)\}`$. We shall now show (using Lemma 2.0.4) that $`S`$ is in fact equal to $`R^G`$.
### 2.2. The first and second fundamental Theorems of classical invariant theory (cf. ) for the action of $`GL_n(K)`$:
###### Theorem 2.2.1.
Let $`G=GL_n(K)`$. Let $`X`$ be as above. The morphism $`\psi :XM_{m,q}`$,
$`(\underset{ยฏ}{u},\underset{ยฏ}{\xi })\left(\phi _{ij}(\underset{ยฏ}{u},\underset{ยฏ}{\xi })\right)(=\left(u_i,\xi _j\right))`$ maps $`X`$ into the determinantal variety $`D_{n+1}(M_{m,q})`$, and the induced homomorphism $`\psi ^{}:K[D_{n+1}(M_{m,q})]K[X]`$ between the coordiante rings induces an isomorphism $`\psi ^{}:K[D_{n+1}(M_{m,q})]K[X]^G`$, i.e. the determinantal variety $`D_{n+1}(M_{m,q})`$ is the categorical quotient of $`X`$ by $`G`$.
###### Proof.
Clearly, $`\psi (X)D_{n+1}(M_{m,q})`$ (since, $`\psi (X)=SpecS`$, and clearly $`SpecSD_{n+1}(M_{m,q})`$ (since any $`n+1`$ vectors in $`V`$ are linearly independent)). We shall prove the result using Lemma 2.0.4. To be very precise, we shall first check the conditions (i)-(iii) of Lemma 2.0.4 for $`\psi :XM_{m,q}`$, deduce that the inclusion $`SpecSD_{n+1}(M_{m,q})`$ is in fact an equality, and hence conclude the normality of $`SpecS`$ (condition (iv) of Lemma 2.0.4).
(i) Let $`x=(\underset{ยฏ}{u},\underset{ยฏ}{\xi })=(u_1,\mathrm{},u_m,\xi _1,\mathrm{},\xi _q)X^{\text{ss}}`$. Let $`W_{\underset{ยฏ}{u}}`$ be the subspace of $`V`$ spanned by $`x_i`$โs and $`W_{\underset{ยฏ}{\xi }}`$ the subspace of $`V^{}`$ spanned by $`\xi _j`$โs. Assume if possible that $`\psi ((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))=0`$, i.e. $`u_i,\xi _j=0`$ for all $`i,j`$.
Case (a): $`W_{\underset{ยฏ}{\xi }}=0`$, i.e., $`\xi _j=0`$ for all $`j`$.
Consider the one parameter subgroup $`\mathrm{\Gamma }=\{g_t,t0\}`$ of $`GL(V)`$, where $`g_t=tI_n`$, $`I_n`$ being the $`n\times n`$ identity matrix. Then $`g_tx=g_t(\underset{ยฏ}{u},\underset{ยฏ}{0})=(t\underset{ยฏ}{u},\underset{ยฏ}{0})`$, so that $`g_tx(0)`$ as $`t0`$. Thus the origin 0 is in the closure of $`Gx`$, and consequently $`x`$ is not semi-stable, which is a contradiction.
Case (b): $`W_{\underset{ยฏ}{\xi }}0`$.
Since the case $`W_{\underset{ยฏ}{u}}=0`$ is similar to Case (a), we may assume that $`W_{\underset{ยฏ}{u}}0`$. Also the fact that $`W_{\underset{ยฏ}{\xi }}0`$ together with the assumption that $`x_i,\xi _j=0`$ for all $`i,j`$ implies that dim$`W_{\underset{ยฏ}{u}}<n`$. Let $`r=dimW_{\underset{ยฏ}{u}}`$ so that we have $`0<r<n`$. Hence, we can choose a basis $`\{e_1,\mathrm{},e_n\}`$ of $`V`$ such that $`W_{\underset{ยฏ}{u}}=e_1,\mathrm{},e_r`$, $`r<n`$, and $`W_{\underset{ยฏ}{\xi }}e_{r+1}^{},\mathrm{},e_n^{}`$, where $`\{e_1^{},\mathrm{},e_n^{}\}`$ is the dual basis in $`V^{}`$. Consider the one parameter subgroup $`\mathrm{\Gamma }=\{g_t,t0\}`$ of $`GL(V)`$, where
$$g_t=\left(\begin{array}{cc}tI_r& 0\\ 0& t^1I_{nr}\end{array}\right).$$
We have $`g_t(\underset{ยฏ}{u},\underset{ยฏ}{\xi })=(t\underset{ยฏ}{u},t\underset{ยฏ}{\xi })0`$ as $`t0`$. Thus, by the same reasoning as in Case (a), the point $`(\underset{ยฏ}{u},\underset{ยฏ}{\xi })`$ is not semi-stable, which leads to a contradiction. Hence we obtain $`\psi ((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))0`$.
(ii) Let
$$U=\{(\underset{ยฏ}{u},\underset{ยฏ}{\xi })X\{u_1,\mathrm{},u_n\},\{\xi _1,\mathrm{},\xi _n\}\text{are linearly independent}\}$$
Clearly, $`U`$ is a $`G`$-stable open subset of $`X`$.
Claim : $`G`$ operates freely on $`U`$, $`UUmodG`$ is a $`G`$-principal fiber space, and $`\psi `$ induces an immersion $`U/GM_{m,q}`$.
Proof of Claim: We have a $`G`$-equivariant identification
$`()`$
$$UG\times G\times \underset{(mn)\text{copies}}{\underset{}{V\times \mathrm{}\times V}}\times \underset{(qn)\text{copies}}{\underset{}{V^{}\times \mathrm{}\times V^{}}}$$
from which it is clear that and $`G`$ operates freely on $`U`$. Further, we see that
$`UmodG`$ may be identified with the fiber space with base $`(G\times G)modG`$ ($`G`$ acting on $`G\times G`$ as $`g(g_1,g_2)=(g_1g,g^1g_2),g,g_1,g_2G`$), and fiber $`\underset{(mn)\text{copies}}{\underset{}{V\times \mathrm{}\times V}}\times \underset{(qn)\text{copies}}{\underset{}{V^{}\times \mathrm{}\times V^{}}}`$ associated to the principal fiber space $`G\times G(G\times G)/G`$. It remains to show that $`\psi `$ induces an immersion $`U/G๐ธ^N`$, i.e., to show that the map $`\psi :U/GM_{m,q}`$ and its differential $`d\psi `$ are both injective. We first prove that $`\psi :U/GM_{m,q}`$ is injective. Let $`x,x^{}`$ in $`U/G`$ be such that $`\psi (x)=\psi (x^{})`$. Let $`\eta ,\eta ^{}U`$ be lifts for $`x,x^{}`$ respectively. Using the identification $`()`$ above, we may write
$$\begin{array}{c}\eta =(A,u_{n+1},\mathrm{},u_m,B,\xi _{n+1},\mathrm{},\xi _q),A,BG\\ \eta ^{}=(A^{},u_{n+1}^{},\mathrm{},u_m^{},B^{},\xi _{n+1}^{},\mathrm{},\xi _q^{}),A^{},B^{}G\end{array}$$
(here, $`u_i,1in`$ are given by the rows of $`A`$, while $`\xi _i,1in`$ are given by the columns of $`B`$; similar remarks on $`u_i^{},\xi _i^{}`$). The hypothesis that $`\psi (x)=\psi (x^{})`$ implies in particular that
$$u_i,\xi _j=u_i^{},\xi _j^{},1i,jn$$
which may be written as
$$AB=A^{}B^{}$$
This implies that $`A^{}=Ag`$, where $`g=BB^1`$. Hence on $`U/G`$, we may suppose that
$$\begin{array}{c}x=(u_1,\mathrm{},u_n,u_{n+1},\mathrm{},u_m,\xi _1,\mathrm{},\xi _q)\\ x^{}=(u_1,\mathrm{},u_n,u_{n+1}^{},\mathrm{},u_m^{},\xi _1^{},\mathrm{},\xi _q^{})\end{array}$$
where $`\{u_1,\mathrm{},u_n\}`$ is linearly independent.
For a given $`j`$, we have,
$$u_i,\xi _j=u_i,\xi _j^{},1in,\text{implies }\xi _j=\xi _j^{}$$
(since $`\{u_1,\mathrm{},u_n\}`$ is linearly independent). Thus we obtain
$`()`$
$$\xi _j=\xi _j^{},\text{for all }j$$
On the other hand, we have (by definition of $`U`$) that $`\{\xi _1,\mathrm{},\xi _n\}`$ is linearly independent. Hence fixing an $`i,n+1im`$, we get
$$u_i,\xi _j=u_i^{},\xi _j(=u_i^{},\xi _j^{}),1jn,\text{implies }u_i=u_i^{}$$
Thus we obtain
$`()`$
$$u_i=u_i^{},\text{for all }i$$
The injectivity of $`\psi :U/GM_{m,q}`$ follows from (\*) and (\**).
To prove that the differential d$`\psi `$ is injective, we merely note that the above argument remains valid for the points over $`K[ฯต]`$, the algebra of dual numbers ($`=KKฯต`$, the $`K`$-algebra with one generator $`ฯต`$, and one relation $`ฯต^2=0`$), i.e., it remains valid if we replace $`K`$ by $`K[ฯต]`$, or in fact by any $`K`$-algebra.
(iii) We have
$$dimU/G=dimUdimG=(m+q)nn^2=dimD_{n+1}(M_{m,q}).$$
The immersion $`U/GSpecS(D_{n+1}(M_{m,q}))`$ together with the fact above that $`dimU/G=dimD_{n+1}(M_{m,q})`$ implies that $`SpecS`$ in fact equals $`dimD_{n+1}(M_{m,q})`$.
(iv) The normality of $`SpecS(=D_{n+1}(M_{m,q}))`$ follows from Theorem 1.5.1 (and the normality of Schubert varieties). โ
Combining the above Theorem with Theorem 1.6.2, we obtain the following
###### Corollary 2.2.2.
Let $`X`$ and $`G`$ be as above. Let $`\phi _{ij}`$ denote the regular function $`(\underset{ยฏ}{u},\underset{ยฏ}{\xi })u_i,\xi _j`$ on $`X`$, $`1im`$, $`1jq`$, and let $`f`$ denote the $`m\times q`$ matrix $`\left(\phi _{ij}\right)`$. The ring of invariants $`K[X]^G`$ has a basis consisting of standard monomials in the regular functions $`p_{\lambda ,\mu }(f)`$ with $`\mathrm{\#}\lambda n`$, where $`\mathrm{\#}\lambda =t`$ is the number of elements in the sequence $`\lambda =(\lambda _1,\mathrm{},\lambda _t)`$ and $`p_{\lambda ,,\mu }(f)`$ is the $`t`$-minor with row indices $`\lambda _1,\mathrm{},\lambda _t`$ and column indices $`\mu _1,\mathrm{},\mu _t`$.
As a consequence of the above Theorem, we obtain the first and second fundamental Theorems of classical invariant theory (cf. ). Let notation be as above.
###### Theorem 2.2.3.
1. First fundamental theorem
The ring of invariants $`K[X]^{GL(V)}`$ is generated by $`\phi _{ij},1im,1jq`$.
2. Second fundamental theorem
The ideal of relations among the generators in (1) is generated by the $`(n+1)`$-minors of the $`m\times q`$-matrix $`(\phi _{ij})`$.
Further, we have (cf. Corollary 2.2.2) :
###### Theorem 2.2.4.
A standard monomial basis for the ring of invariants: The ring of invariants $`K[X]^{GL(V)}`$ has a basis consisting of standard monomials in the regular functions $`p_(A,B),AI(r,m),BI(r,q),rn`$.
## 3. The $`K`$-algebra $`S`$
Let $`X`$ be as above. We shall denote $`K[X]`$ by $`R`$ so that $`R=K[u_{ij},\xi _{kl}\mathrm{\hspace{0.17em}1}im,\mathrm{\hspace{0.17em}1}j,kn,1lq]`$.
The functions $`u(I),\xi (J)`$: As above, let
$`U=\left(u_{ij}\right)_{1im,\mathrm{\hspace{0.17em}1}jn}`$ and $`W=\left(\xi _{kl}\right)_{1kn,\mathrm{\hspace{0.17em}1}lq}`$. For $`II(n,m),JI(n,q)`$, let $`u(I),\xi (J))`$ denote the following regular functions on $`X`$:
$`u(I)((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))=`$ the $`n`$-minor of $`U`$ with row indices given by $`I`$.
$`\xi (J)((\underset{ยฏ}{u},\underset{ยฏ}{\xi }))=`$ the $`n`$-minor of $`W`$ with column indices given by $`J`$.
Note that for the diagonal action of $`SL_n(K)(=SL(V))`$ on $`X`$, we have, $`u(I),\xi (J)`$ are in $`R^{SL_n(K)}`$.
The $`K`$-algebra $`S`$: Let $`S`$ be the $`K`$-subalgebra of $`R`$ generated by
$`\{u(I),\xi (J),p(A,B),II(n,m),JI(n,q),AI(r,m),BI(r,q),1rn\}`$. We shall denote the set $`I(n,m)`$ indexing the $`u(I)`$โs by $`H_u`$ and the set $`I(n,q)`$ indexing the $`\xi (J)`$โs by $`H_\xi .`$ Also, we shall denote $`H_p:=\underset{1rn}{}(I(r,m)\times I(r,q))`$, and set
$`H`$ $`=`$ $`H_u\dot{}H_\xi H_p`$
$`=`$ $`I(n,m)\dot{}I(n,q)\underset{1rn}{}(I(r,m)\times I(r,q)),`$
where $`\dot{}`$ denotes a disjoint union. (If $`m=q`$, then $`H_u,H_\xi `$ are to be considered as two disjoint copies of $`I(n,m)`$.) Then the algebra generators
$`\{u(I),\xi (J),p(A,B),II(n,m),JI(n,q),AI(r,m),BI(r,q),1rn\}`$ of $`S`$ are indexed by the set $`H`$. Clearly $`SR^{SL(V)}`$.
###### Remark 3.0.1.
The $`K`$-algebra $`S`$ could have been simply defined as the $`K`$-subalgebra of $`R^G`$ generated by $`\{u_i,\xi _j\}`$ (i.e., by $`\{p(A,B),\mathrm{\#}A=\mathrm{\#}B=1\}`$) and $`\{u(I),\xi (J)\}`$. But we have a purpose in defining it as above, namely, the standard monomials (in $`S`$) will be built out of the $`p(A,B)`$โs with $`\mathrm{\#}An`$, the $`u(I)`$โs and $`\xi (J)`$โs (cf. Definition 4.0.1).
Our goal is to show that $`S`$ equals $`R^{SL(V)}`$.
A partial order on $`H`$: Define a partial order on $`H`$ as follows:
1. The partial order on $`H_p`$ is as in ยง1.6 (note that $`H_pH_{m,q}`$)
2. The partial order on $`H_u`$ and $`H_\xi `$ are as in ยง1.1.
3. Any element of $`H_u`$ and any element of $`H_\xi `$ are not comparable.
4. No element of $`H_u,H_\xi `$ is greater than any element of $`H_p`$.
5. For $`IH_u`$ and $`(A,B)H_p`$, we define $`I(A,B)`$ if $`IA`$ (the partial order being as in ยง1.6). Similarly, for $`JH_\xi `$ and $`(A,B)H_p`$, we define $`J(A,B)`$ if $`JB`$.
###### Lemma 3.0.2.
$`H`$ is a ranked poset of rank $`d:=(m+q)nn^2`$, i.e., all maximal chains in $`H`$ have the same cardinality $`=(m+q)nn^2+1`$.
###### Proof.
Clearly, $`H`$ is a ranked poset (since it is composed of ranked posets). To compute the rank of $`H`$, we consider the maximal chain consisting of $`\tau _1,\mathrm{},\tau _N`$, where
the first $`q`$ of them are given by $`(m,q),(m,q1),\mathrm{},(m,1)`$ (of $`H_p`$),
the next $`(m1)`$ of them are given by $`(m1,1),(m2,1),\mathrm{},(1,1)`$ (of $`H_p`$).
(thus contributing $`m+q1`$ to the cardinality of the chain).
This is now followed by the $`q1`$ elements of $`H_p`$ :
$`((1,m),(1,q)),((1,m),(1,q1)),\mathrm{},((1,m),(1,2))`$ ,
followed by the $`m2`$ elements of $`H_p`$ :
$`((1,m1),(1,2)),((1,m2),(1,2)),\mathrm{},((1,2),(1,2))`$
(thus contributing $`m+q3`$ to the cardinality of the chain).
Thus proceeding, finally, we end up with $`((1,2,\mathrm{},n),(1,2,\mathrm{},n))`$ (in $`H_p`$). This is now followed by either $`(1,2,\mathrm{},n)`$ of $`H_u`$ or $`(1,2,\mathrm{},n)`$ of $`H_\xi `$.
The number of elements in the above chain equals
$`[m+q1+(m+q3)+\mathrm{}+m+q(2n1)]+1=(m+q)nn^2+1`$
## 4. Standard monomials in the $`K`$-algebra $`S`$
###### Definition 4.0.1.
A monomial $`F`$ in the $`p(A,B)`$โs, $`u(I)`$โs, and $`\xi (J)`$โs, is said to be *standard* if $`F`$ satisfies the following conditions:
1. If $`F`$ involves $`u(I)`$, for some $`I`$ (resp. $`\xi (J)`$ for some $`J`$), then $`F`$ does not involve $`\xi (J^{})`$ for any $`J^{}`$ (resp. $`u(I^{})`$, for any $`I^{}`$).
2. If $`F=p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$
(resp. $`p(A_1,B_1)\mathrm{}p(A_r,B_r)\xi (J_1)\mathrm{}\xi (J_t)`$), where $`r,s,t`$ are integers $`0`$, then
$$A_1\mathrm{}A_rI_1\mathrm{}I_s(\mathrm{resp}.B_1\mathrm{}B_rJ_1\mathrm{}J_t)$$
### 4.1. Quadratic relations
In this subsection, we describe certain straightening relations to be used while proving the linear independence of standard monomials and generation (of $`S`$ as a $`K`$-vector space) by standard monomials.
###### Theorem 4.1.1.
1. Let $`IH_u,JH_\xi `$. We have
$$u(I)\xi (J)=p(I,J)$$
2. Let $`I,I^{}H_u`$ be not comparable. We have,
$$u(I)u(I^{})=\underset{๐}{}b_ru(I_r)u(I_r^{}),b_rK^{}$$
where for all $`r`$, $`I_r`$ both $`I,I^{}`$, and $`I_r^{}`$ both $`I,I^{}`$.
3. Let $`J,J^{}H_\xi `$ be not comparable. We have,
$$\xi (J)\xi (J^{})=\underset{๐ }{}c_s\xi (J_s)\xi (J_s^{}),c_sK^{}$$
where for all $`s`$, $`J_s`$ both $`J,J^{}`$, and $`J_s^{}`$ both $`J,J^{}`$.
4. Let $`(A_1,A_2),(B_1,B_2)H_p`$ be not comparable. Then we have
$$p(A_1,A_2)p(B_1,B_2)=\underset{๐}{}a_ip(C_{i1},C_{i2})p(D_{i1},D_{i2}),a_iK^{},$$
where $`(C_{i1},C_{i2}),(D_{i1},D_{i2})`$ belong to $`H_p`$; further, for every $`i`$, we have
1. $`C_{i1}`$ both $`A_1`$ and $`B_1`$; $`D_{i1}`$ both $`A_1`$ and $`B_1`$.
2. $`C_{i2}`$ both $`A_2`$ and $`B_2`$; $`D_{i2}`$ both $`A_2`$ and $`B_2`$.
5. Let $`IH_u,(A,B)H_p`$ be such that $`AI`$. We have,
$$p(A,B)u(I)=\underset{๐ก}{}d_tp(A_t,B_t)u(I_t),d_tK^{}$$
where for every $`t`$, we have, $`A_t`$ (resp. $`I_t`$) both $`A`$ and $`I`$, and $`B_tB`$.
6. Let $`JH_\xi ,(A,B)H_p`$ be such that $`BJ`$. We have,
$$p(A,B)\xi (J)=\underset{๐}{}e_lp(A_l,B_l)\xi (J_l),e_lK^{}$$
where for every $`l`$, we have, $`A_lA`$ , and $`B_l`$ (resp. $`J_l`$) both $`B`$ and $`J`$.
###### Proof.
In the course of the proof, we will be repeatedly using the fact that the subalgebra generated by $`\{p(A,B),AI(r,m),BI(r,q)\mathrm{\hspace{0.17em}1}rn\}`$ being $`R^{GL(V)}`$ (cf. Theorem 2.2.3,(1)), the results given in Theorem 2.2.3,(1), Theorem 2.2.4 apply to this subalgebra.
(1) is clear from the definitions of $`u(I),\xi (J)`$ and $`p(I,J)`$.
(2). We shall denote a minor of $`U(=(u_{ij})_{1im,\mathrm{\hspace{0.17em}1}jn})`$ with rows and columns given by $`I,J`$ (where $`I,JI(r,m)`$ for some $`rn`$) by $`\mathrm{\Delta }(I,J)`$. Observe that if #$`I=n`$, then $`J=(1,2,\mathrm{},n)`$ necessarily (since $`U`$ has size $`m\times n`$). Thus for $`IH_u`$, we have that $`u(I)=\mathrm{\Delta }(I,I_n),u(I^{})=\mathrm{\Delta }(I^{},I_n)`$ (as minors of $`U`$), where $`I_n=(1,2,\mathrm{},n)`$, we have, in view of Theorem 1.6.6, (2),
$$u(I)u(I^{})=\mathrm{\Delta }(I,I_n)\mathrm{\Delta }(I^{},I_n)=\underset{๐}{}b_i\mathrm{\Delta }(C_{i1},C_{i2})\mathrm{\Delta }(D_{i1},D_{i2}),a_iK^{},$$
where we have for every $`i`$, $`C_{i1}`$ both $`I`$ and $`I^{}`$; $`D_{i1}`$ both $`I`$ and $`I^{}`$; $`C_{i2}`$ $`I_n`$; $`D_{i2}I_n`$ which forces $`\mathrm{\#}D_{i2}=n`$ (in view of the partial order (cf. ยง1.6); note that $`D_{i2}`$ being the column indices of a minor of the $`m\times n`$ matrix $`U`$, we have that $`\mathrm{\#}D_{i2}n`$). Hence we obtain that $`D_{i2}=I_n`$, for all $`i`$. In particular, we obtain that $`\mathrm{\#}D_{i1}(=\mathrm{\#}D_{i2})=n`$. This in turn implies (by consideration of the degrees in $`u_{ij}`$โs of the terms in the above sum) that $`\mathrm{\#}C_{i1}=\mathrm{\#}C_{i2}=n`$. Hence $`C_{i2}=I_n`$ (again note that $`C_{i2}`$ gives the column indices of the $`n`$-minor $`\mathrm{\Delta }(C_{i1},C_{i2})`$ of the $`m\times n`$ matrix $`U`$). Thus the above relation becomes
$$u(I)u(I^{})=\underset{๐}{}b_iu(C_{i1})u(D_{i1}),$$
with $`C_{i1}`$ both $`I`$ and $`I^{}`$; $`D_{i1}`$ both $`I`$ and $`I^{}`$. This proves (2).
Proof of (3) is similar to that of (2).
(4) is a direct consequence of Theorem 2.2.3,(2) and Proposition 1.6.3.
(5). If $`\mathrm{\#}A=n=\mathrm{\#}B`$, then $`p(A,B)u(I)=u(A)u(I)\xi (B)`$. By (2),
$$u(A)u(I)=\underset{๐}{}d_iu(C_i)u(D_i),d_iK^{}$$
where $`C_i`$ both $`A,I`$, and $`D_i`$ both $`A,I`$. Hence
$$p(A,B)u(I)=\underset{๐}{}d_iu(C_i)u(D_i)\xi (B)=\underset{๐}{}d_ip(C_i,B)u(D_i)$$
where $`C_i`$ both $`A,I`$, and $`D_i`$ both $`A,I`$, and the result follows.
Let then $`\mathrm{\#}A<n`$. By (1), we have $`u(I)\xi (I_n)=p(I,I_n)`$. Hence, $`p(A,B)u(I)\xi (I_n)=p(A,B)p(I,I_n)`$. The hypothesis that $`AI`$ implies that $`p(A,B)p(I,I_n)`$ is not standard (note that the facts that $`\mathrm{\#}A<n,\mathrm{\#}I=n`$ implies that $`IA`$). Hence (4) implies that
$$p(A,B)p(I,I_n)=a_ip(C_{i1},C_{i2})p(D_{i1},D_{i2}),a_iK^{},$$
where $`(C_{i1},C_{i2}),(D_{i1},D_{i2})`$ belong to $`H_p`$; further, for every $`i`$, $`C_{i1}`$ both $`A`$ and $`I`$; $`D_{i1}`$ both $`A`$ and $`I`$; $`C_{i2}`$ both $`B`$ and $`I_n`$; $`D_{i2}`$ both $`B`$ and $`I_n`$ which forces $`D_{i2}=I_n`$ (note that in view of Theorem 2.2.4, all minors in the above relation have size $`n`$); and hence $`\mathrm{\#}D_{i1}=n`$, for all $`i`$. Hence $`p(D_{i1},D_{i2})=u(D_{i1})\xi (I_n)`$, for all $`i`$. Hence cancelling $`\xi (I_n)`$, we obtain
$$p(A,B)u(I)=a_ip(C_{i1},C_{i2})u(D_{i1}),$$
where $`C_{i1}`$ both $`A`$ and $`I`$, $`D_{i1}`$ both $`A`$ and $`I`$, and $`C_{i2}B`$. This proves (5).
Proof of (6) is similar to that of (5).
### 4.2. Linear independence of standard monomials:
In this subsection, we prove the linear independence of standard monomials.
###### Lemma 4.2.1.
Let $`(A,B)H_p,IH_u,JH_\xi `$.
1. The set of standard monomials in the $`p(A,B)`$โs is linearly independent.
2. The set of standard monomials in the $`u(I)`$โs is linearly independent.
3. The set of standard monomials in the $`\xi (J)`$โs is linearly independent.
###### Proof.
(1) follows from Theorem 2.2.4.
(2), (3) follow from Theorem 1.6.6,(1) applied to $`K[u_{ij},1im,1jn]`$,
$`K[\xi _{kl},1kn,1lq]`$ respectively. โ
###### Proposition 4.2.2.
Standard monomials are linearly independent.
###### Proof.
For a monomial $`M`$, by $`u`$-*degree* (resp. $`\xi `$-*degree*) of $`M`$, we shall mean the degree of $`M`$ in the variables $`u_{ij}`$โs (resp. $`\xi _{kl}`$โs ). We have
$$\begin{array}{c}u\mathrm{degree}\mathrm{of}p(A_1,B_1)\mathrm{}p(A_r,B_r)=\xi \mathrm{degree}\mathrm{of}p(A_1,B_1)\mathrm{}p(A_r,B_r)=\underset{๐}{}\mathrm{\#}A_i\\ u\mathrm{degree}\mathrm{of}u(I_1)\mathrm{}u(I_s)=ns,\xi \mathrm{degree}\mathrm{of}u(I_1)\mathrm{}u(I_s)=0\\ \xi \mathrm{degree}\mathrm{of}\xi (J_1)\mathrm{}\xi (J_t)=nt,u\mathrm{degree}\mathrm{of}\xi (J_1)\mathrm{}\xi (J_t)=0\end{array}$$
By considering the $`u`$-degree and the $`\xi `$-degree, and using Lemma 4.2.1 we see that
$`\{p(A_1,B_1)\mathrm{}p(A_r,B_r),u(I_1)\mathrm{}u(I_s),\xi (J_1)\mathrm{}\xi (J_t),r,s,t0\}`$ is linearly independent.
Let
$`()`$
$$F:=R+S=0$$
be a relation among standard monomials, where $`R=a_iM_i,S=b_iN_i`$ such that each $`M_i`$ (resp. $`N_i`$) is a standard monomial of the form $`p(A_1,B_1)\mathrm{}p(A_{r_i},B_{r_i})`$ (resp. $`p(A_1,B_1)\mathrm{}p(A_{q_i},B_{q_i})u(I_1)\mathrm{}u(I_{s_i})\xi (J_1)\mathrm{}\xi (J_{t_i}),q_i0`$, and at least one of $`\{s_i,t_i\}>0`$ ). If $`g`$ is in $`GL_n(K)`$, with det$`g`$ a root of unity, then using the facts that $`gp(A,B)=p(A,B),gu(I)=(detg)u(I),g\xi (J)=(detg)\xi (J)`$, we have, $`FgF=b_i(1(detg)^{s_i+t_i})N_i=0`$. Hence if we show that the $`N_i`$โs are linearly independent, then (in view of Lemma 4.2.1,(1)), we would obtain that (\*) is the trivial relation. Thus we may suppose that
$`()`$
$$F=b_iN_i=0,$$
where each $`N_i`$ is a standard monomial of the form
$$p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)\xi (J_1)\mathrm{}\xi (J_t)$$
where $`r0`$ and at least one of $`\{s,t\}>0`$; in fact, $`N_i`$โs being standard, in any $`N_i`$, precisely one of $`\{s_i,t_i\}`$ is non-zero.
We first multiply (\**) by $`u(I_n)^N`$ ($`I_n`$ being $`(1,2,\mathrm{},n)`$), for a sufficiently large $`N`$ ($`N`$ could be taken to be any integer greater than all of the $`t`$โs, appearing in the $`\xi (J_1)\mathrm{}\xi (J_t)`$โs); we then replace a $`\xi (J)u(I_n)`$ by $`p(I_n,J)`$ (cf. Theorem 4.1.1, (1)). Then in the resulting sum, any monomial will involve only the $`p(A,B)`$โs and the $`u(I)`$โs. Thus we may suppose that (\**) is of the form
$`()`$
$$G:=c_iG_i=0$$
where each $`G_i`$ is of the form $`p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$. Note that for each standard monomial *M*=$`p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$
(resp. $`p(A_1,B_1)\mathrm{}p(A_r,B_r)\xi (J_1)\mathrm{}\xi (J_t)`$) appearing in (\**), *M*$`u(I_n)^N`$ is again standard. Again, considering $`GgG,gGL_n(K)`$, with det$`g`$ a root of unity, as above, we may suppose that in each monomial $`p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$ appearing in (\***), $`s>0`$. Further, in view of Lemma 4.2.1,(2), we may suppose that for at least one monomial $`r>0`$. Now considering the $`\xi `$-degree of the monomials, we may suppose (in view of Lemma 4.2.1,(2)) that in each monomial $`p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$ appearing in (\***), $`r>0`$.
Thus, for each monomial $`p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$ appearing in (\***), we have, $`r,s>0`$. Now the $`\xi `$-degree (as well as the $`u`$-degree) being the same for all of the monomials in (\***), for any two monomials $`G_i,G_i^{}`$, say
$$G_i=p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s),G_i^{}=p(A_1^{},B_1^{})\mathrm{}p(A_r^{}^{},B_r^{}^{})u(I_1^{})\mathrm{}u(I_s^{}^{})$$
we have $`\underset{1ir}{}\mathrm{\#}A_i=\underset{1ir^{}}{}\mathrm{\#}A_i^{}`$. This together with the fact that the $`u`$-degree is the same for all of the terms $`G_k`$โs in (\***) implies that $`s=s^{}`$. Thus we obtain that in all of the monomials $`p(A_1,B_1)\mathrm{}p(A_r,B_r)u(I_1)\mathrm{}u(I_s)`$ in (\***), the integer $`s`$ is the same (and $`s>0`$). Now we multiply (\***) through out by $`\xi (I_n)^s`$ (where $`I_n=(1,2,\mathrm{},n)`$) to arrive at a linear sum
$$d_iH_i=0$$
where each $`H_i`$ is a standard monomial in the $`p(A,B)`$โs (note that
$`H_i=p(A_1,B_1)\mathrm{}p(A_r,B_r)p(I_1,I_n)\mathrm{}p(I_s,I_n)`$ is standard). Now the required result follows from the linear independence of $`p(A,B)`$โs (cf. Lemma 4.2.1,(1)).
### 4.3. The algebra $`S(D)`$
To prove the generation of $`S`$ (as a $`K`$-vector space) by standard monomials, we define a $`K`$-algebra $`S(D)`$, construct a standard monomial basis for $`S(D)`$ and deduce the results for $`S`$ (in fact, it will turn out that $`S(D)S`$). We first define the $`K`$-algebra $`R(D)`$ as follows:
Let
$$D=H\{\mathrm{๐}\}\{\mathrm{๐}\}$$
$`H`$ being as in the beginning of ยง3. Extend the partial order on $`H`$ to $`D`$ by declaring $`\{\mathrm{๐}\}`$ (resp. $`\{\mathrm{๐}\}`$) as the largest (resp. smallest) element. Let $`P(D)`$ be the polynomial algebra
$$P(D):=K[X(A,B),Y(I),Z(J),X(\mathrm{๐}),X(\mathrm{๐}),(A,B)H_p,IH_u,JH_\xi ]$$
Let $`๐(D)`$ be the homogeneous ideal in the polynomial algebra $`P(D)`$ generated by the six relations of Theorem 4.1.1 ($`X(A,B),Y(I),Z(J)`$ replacing $`p(A,B),u(I),\xi (J)`$ respectively), with relations (1) and (4) homogenized as follows: (1) is homogenized as
$`()`$
$$X(I)Y(J)=X(I,J)X(\mathrm{๐})$$
while (4) is homogenized as
$$X(A_1,A_2)X(B_1,B_2)=a_iX(C_{i1},C_{i2})X(D_{i1},D_{i2})$$
where $`X(C_{i1},C_{i2})`$ is to be understood as $`X(\mathrm{๐})`$ if both $`C_{i1},C_{i2}`$ equal the empty set (cf. Remark 1.6.5). Let
$$R(D)=P(D)/๐(D)$$
We shall denote the classes of $`X(A,B),Y(I),Z(J),X(\mathrm{๐}),X(\mathrm{๐})`$ in $`R(D)`$ by
$`x(A,B),y(I),z(J),x(\mathrm{๐}),x(\mathrm{๐})`$ respectively.
The algebra $`M(D)`$: Set $`M(D)=R(D)_{(x(\mathrm{๐}))}`$, the homogeneous localization of $`R(D)`$ at $`x(\mathrm{๐})`$. We shall denote $`\frac{x(\mathrm{๐})}{x(\mathrm{๐})},\frac{x(A,B)}{x(\mathrm{๐})},\frac{y(I)}{x(\mathrm{๐})},\frac{z(J)}{x(\mathrm{๐})}`$ (in $`M(D)`$) by $`q(\mathrm{๐}),r(A,B),s(I),t(J)`$ respectively.
A grading for $`M(D)`$: We give a grading for $`M(D)`$ by assigning degree one to $`s(I),t(J)`$, and degree 2 to $`q(\mathrm{๐}),r(A,B)`$, where as above $`IH_u,JH_\xi ,(A,B)H_p`$.
The algebra $`S(D)`$: Set $`S(D)=M(D)_{(q(\mathrm{๐}))}`$, the homogeneous localization of $`M(D)`$ at $`q(\mathrm{๐})`$. We shall denote $`\frac{r(A,B)}{q(\mathrm{๐})},\frac{s(I)}{q(\mathrm{๐})},\frac{t(J)}{q(\mathrm{๐})}`$ (in $`S(D)`$) by $`c(A,B),d(I),e(J)`$ respectively.
Let $`\phi _D:S(D)S`$ be the map, $`\phi _D(c(A,B))=p(A,B),\phi _D(d(I))=u(I),\phi _D(e(J))=\xi (J)`$. Consider the canonical maps
$$\theta _D:R(D)M(D),\delta _D:M(D)S(D)$$
Denote $`\gamma _D:R(D)S`$ as the composite $`\gamma _D=\phi _D\delta _D\theta _D`$.
### 4.4. A standard monomial basis for $`R(D)`$:
We define a monomial in
$`x(A,B),y(I),z(J),x(\mathrm{๐}),x(\mathrm{๐})`$ (in $`R(D)`$) to be standard in exactly the same way as in Definition 4.0.1 (we declare $`x(\mathrm{๐})`$ (resp. $`x(\mathrm{๐})`$) as the largest (resp. smallest)).
###### Proposition 4.4.1.
The standard monomials in the $`x(A,B)`$โs, $`y(I)`$โs, $`z(J)`$โs, $`x(\mathrm{๐})`$โs, $`x(\mathrm{๐})`$โs are linearly independent.
###### Proof.
The result follows by considering $`\gamma _D:R(D)S`$, and using the linear independence of standard monomials in $`S`$ (cf. Proposition 4.2.2). โ
Generation of $`R(D)`$ by standard monomials: We shall now show that any non-standard monomial $`F`$ in $`R(D)`$ is a linear sum of standard monomials. Observe that if $`M`$ is a standard monomial, then $`x(\mathrm{๐})^lM`$ (resp. $`Mx(\mathrm{๐})^l`$) is again standard; hence we may suppose $`F`$ to be:
$$F=x(A_1,B_1)\mathrm{}x(A_r,B_r)y(I_1)\mathrm{}y(I_s)z(J_1)\mathrm{}z(J_t)$$
Using the relations $`y(I)z(J)=x(I,J)x(\mathrm{๐})`$, we may suppose that
$`F=x(A_1,B_1)\mathrm{}x(A_r,B_r)y(I_1)\mathrm{}y(I_s)`$ or $`F=x(A_1,B_1)\mathrm{}x(A_r,B_r)z(J_1)\mathrm{}z(J_t)`$, say, $`F=x(A_1,B_1)\mathrm{}x(A_r,B_r)y(I_1)\mathrm{}y(I_s)`$.
Fix an integer $`N`$ sufficiently large. To each element $`A_{r=1}^nI(r,m)`$, we associate an $`(n+1)`$-tuple as follows: Let $`AI(r,m)`$, for some $`r`$, say, $`A=(a_1,\mathrm{},a_r)`$. To $`A`$, we associate the $`n+1`$-tuple
$$\overline{A}:=(a_1,\mathrm{},a_r,m,m,\mathrm{},m,1)$$
Similarly, for $`B_{r=1}^nI(r,q)`$, say, $`B=(b_1,\mathrm{},b_r)`$, we associate the $`n+1`$-tuple
$$\overline{B}:=(b_1,\mathrm{},b_r,q,q,\mathrm{},q,1)$$
To $`F`$, we associate the integer $`n_F`$ (and call it the weight of $`F`$) which has the entries of $`\overline{A_1},\overline{B_1},\overline{A_2},\overline{B_2},\mathrm{},\overline{A_r},\overline{B_r},\overline{I_1},\mathrm{},\overline{I_s}`$ as digits (in the $`N`$-ary presentation). The hypothesis that $`F`$ is non-standard implies that
either $`x(A_i,B_i)x(A_{i+1},B_{i+1})`$ is non-standard for some $`ir1`$, or, $`x(A_r,B_r)y(I_1)`$ is non-standard or $`y(I_j)y(I_{j+1})`$ is non-standard for some $`js1`$. Straightening these using Theorem 4.1.1, we obtain that $`F=a_iF_i`$ where $`n_{F_i}>n_F,i`$, and the result follows by decreasing induction on $`n_F`$ (note that while straightening a degree $`2`$ relation using Theorem 4.1.1, (4), if $`x(\mathrm{๐})`$ occurs in a monomial $`G`$, then the digits in $`n_G`$ corresponding to $`x(\mathrm{๐})`$ are taken to be $`(\underset{n+1\text{ times}}{\underset{}{m,m\mathrm{},m}},\underset{n+1\text{ times}}{\underset{}{q,q\mathrm{},q}})`$. Also note that the largest $`F`$ of degree $`r`$ in $`x(A,B)`$โs and degree $`s`$ in the $`y(I)`$โs is $`x(\{m\},\{q\})^ru(I_0)^s`$ (where $`I_0`$ is the $`n`$-tuple $`(m+1n,m+2n,\mathrm{},m)`$) which is clearly standard (the starting point of the decreasing induction).
Hence we obtain
###### Proposition 4.4.2.
Standard monomials in $`x(A,B),y(I),z(J),x(\mathrm{๐}),x(\mathrm{๐})`$ generate $`R(D)`$ as a $`K`$-vector space.
Combining Propositions 4.4.1, 4.4.2, we obtain
###### Theorem 4.4.3.
Standard monomials in $`x(A,B),y(I),z(J),x(\mathrm{๐}),x(\mathrm{๐})`$ give a basis for the $`K`$-vector space $`R(D)`$.
### 4.5. Standard monomial bases for $`M(D),S(D)`$
Standard monomials in
$`r(A,B),s(I),t(J),q(\mathrm{๐}))`$ in $`M(D)`$ (resp. $`c(A,B),d(I),e(J))`$ in $`S(D)`$) are defined in exactly the same way as in Definition 4.0.1.
###### Proposition 4.5.1.
Standard monomials in $`r(A,B),s(I),t(J),q(\mathrm{๐})`$ give a basis for the $`K`$-vector space $`M(D)`$.
###### Proof.
The linear independence of standard monomials follows as in the proof of Prop 4.4.1 by considering $`\phi _D\delta _D:M(D)S`$, and using the linear independence of standard monomials in $`S`$ (cf. Proposition 4.2.2).
To see the generation of $`M(D)`$ by standard monomials, consider a non-standard monomial $`F`$ in $`M(D)`$. Since $`q(\mathrm{๐})^l`$ is the largest monomial of a given degree $`l`$, we may suppose $`F`$ to be:
$$F=r(A_1,B_1)\mathrm{}r(A_i,B_i)s(I_1)\mathrm{}s(I_k)t(J_1)\mathrm{}t(J_l)$$
In view of Theorem 4.1.1, (1), we may suppose that
$`F=r(A_1,B_1)\mathrm{}r(A_i,B_i)s(I_1)\mathrm{}s(I_k)`$ or $`r(A_1,B_1)\mathrm{}r(A_i,B_i)t(J_1)\mathrm{}t(J_l)`$, say, $`F=r(A_1,B_1)\mathrm{}r(A_i,B_i)s(I_1)\mathrm{}s(I_k)`$. Then $`F=\theta _D(H)`$, where
$`H=x(A_1,B_1)\mathrm{}x(A_i,B_i)y(I_1)\mathrm{}y(I_k)`$. The required result follows from Proposition 4.4.2. โ
###### Proposition 4.5.2.
Standard monomials in $`c(A,B),d(I),e(J)`$ give a basis for the $`K`$-vector space $`S(D)`$.
The proof is completely analogous to that of Proposition 4.5.1 (in view of the fact that $`S(D)=M(D)_{(q(\mathrm{๐}))}`$).
###### Theorem 4.5.3.
Standard monomials in $`p(A,B),u(I),\xi (J)`$ form a basis for the $`K`$-vector space $`S`$.
###### Proof.
We already have established the linear independence of standard monomials (cf. Proposition 4.2.2). The generation by standard monomials follows by considering the surjective map $`\phi _D:S(D)S`$ and using the generation of $`S(D)`$ by standard monomials (cf. Theorem 4.5.2). โ
###### Theorem 4.5.4.
The map $`\phi _D:S(D)S`$ is an isomorphism of $`K`$-algebras.
###### Proof.
Under $`\phi _D`$, the standard monomials in $`S(D)`$ are mapped bijectively onto the standard monomials in $`S`$. The result follows from Proposition 4.5.2 and Theorem 4.5.3. โ
###### Theorem 4.5.5.
A presentation for $`\mathrm{S}`$:
1. The $`K`$-algebra $`S`$ is generated by $`\{p(A,B),u(I),\xi (J),(A,B)H_p,IH_u,JH_\xi \}`$.
2. The ideal of relations among the generators $`\{p(A,B),u(I),\xi (J),(A,B)H_p`$,
$`IH_u,JH_\xi \}`$ is generated by the six type of relations as given by Theorem 4.1.1.
###### Proof.
The result follows from Theorem 4.5.4, Proposition 4.5.2 (and the definition of $`S(D)`$) โ
## 5. Normality and Cohen-Macaulayness of the $`K`$-algebra $`S`$
In this section, we prove the normality and Cohen-Macaulayness of $`SpecS`$ by relating it to a toric variety. From ยง3, ยง4, we have
$``$ $`\{u(I),\xi (J),p(A,B),IH_u,JH_\xi ,(A,B)H_p\}`$ generates $`S`$ as a $`K`$-algebra.
$``$ Standard monomials in $`\{u(I),\xi (J),p(A,B),IH_u,JH_\xi ,(A,B)H_p\}`$ form a $`K`$-basis for $`S`$.
$``$ Considering $`S`$ as a quotient of the polynomial algebra
$$K[X(A,B),Y(I),Z(J),(A,B)H_p,IH_u,JH_\xi ]$$
the ideal $`๐`$ of relations is generated by the six kinds of quadratic relations as given in Theorem 4.1.1.
### 5.1. The algebra associated to a distributive lattice
###### Definition 5.1.1.
A lattice is a partially ordered set $`(,)`$ such that, for every pair of elements $`x,y`$, there exist elements $`xy`$, $`xy`$, called the join, respectively the meet of $`x`$ and $`y`$, satisfying:
$$xyx,xyy,\text{ and if }zx\text{ and }zy,\text{ then }zxy,$$
$$xyx,xyy,\text{ and if }zx\text{ and }zy,\text{ then }zxy.$$
###### Definition 5.1.2.
A lattice is called distributive if the following identities hold:
$`x(yz)`$ $`=(xy)(xz)`$
$`x(yz)`$ $`=(xy)(xz)`$
###### Definition 5.1.3.
Given a finite lattice $``$, the ideal associated to $``$, denoted by $`I()`$, is the ideal of the polynomial algebra $`K[](=K[x_\alpha ,\alpha ])`$ generated by the set of binomials
$$๐ข_{}=\{xy(xy)(xy)x,y\text{ non-comparable}\}.$$
Set $`A()=K[]/I()`$, *the algebra associated to* $``$.
The chain lattice $`๐(n_1,\mathrm{},n_d)`$: Given an integer $`n1`$, let $`๐(n)`$ denote the chain $`\{1<\mathrm{}<n\}`$, and for $`n_1,\mathrm{},n_d>1`$, let $`๐(n_1,\mathrm{},n_d)`$ denote the chain product lattice $`๐(n_1)\times \mathrm{}\times ๐(n_d)`$ consisting of all $`d`$-tuples $`(i_1,\mathrm{},i_d)`$, with $`1i_1n_1,\mathrm{},1i_dn_d`$. For $`(i_1,\mathrm{},i_d)`$, $`(j_1,\mathrm{},j_d)`$ in $`๐(n_1,\mathrm{},n_d)`$, we define
$$(i_1,\mathrm{},i_d)(j_1,\mathrm{},j_d)i_1j_1,\mathrm{},i_dj_d.$$
We have
$$\begin{array}{cc}\hfill (i_1,\mathrm{},i_d)(j_1,\mathrm{},j_d)& =(\mathrm{max}\{i_1,j_1\},\mathrm{},\mathrm{max}\{i_d,j_d\})\hfill \\ \hfill (i_1,\mathrm{},i_d)(j_1,\mathrm{},j_d)& =(\mathrm{min}\{i_1,j_1\},\mathrm{},\mathrm{min}\{i_d,j_d\}).\hfill \end{array}$$
Clearly, $`๐(n_1,\mathrm{},n_d)`$ is a finite distributive lattice.
### 5.2. Flat degenerations of certain $`K`$-algebras:
Let $``$ be a finite lattice, and $`R`$ a $`K`$-algebra with generators $`\{p_\alpha \alpha \}`$.
###### Definition 5.2.1.
A monomial $`p_{\alpha _1}\mathrm{}p_{\alpha _r}`$ is said to be standard if $`\alpha _1\mathrm{}\alpha _r`$.
Suppose that the standard monomials form a $`K`$-basis for $`R`$. Given any nonstandard monomial $`F`$, the expression
$$F=c_iF_i,c_iK^{}$$
for $`F`$ as a sum of standard monomials will be referred to as a straightening relation. Consider the surjective map
$$\pi :K[]R,x_\alpha p_\alpha .$$
Let us denote $`\mathrm{ker}\pi `$ by $`I`$.
For $`\alpha ,\beta H`$ with $`\alpha >\beta `$, we set
$$]\beta ,\alpha [=\{\gamma \alpha >\gamma >\beta \}.$$
Recall the following theorem (cf., Theorem 5.2)
###### Theorem 5.2.2.
Let $`,R,I`$ be as above. Suppose that there exists a lattice embedding $`๐`$, where $`๐=๐(n_1,\mathrm{},n_d)`$ for some $`n_1,\mathrm{},n_d1`$, such that the entries of the $`d`$-tuple $`(\theta _1,\mathrm{},\theta _d)`$ representing an element $`\theta `$ of $``$ form a non-decreasing sequence, i.e. $`\theta _1\mathrm{}\theta _d`$. Suppose that $`I`$ is generated as an ideal by elements of the form $`x_\tau x_\phi c_{\alpha \beta }x_\alpha x_\beta `$(where $`\tau ,\phi `$ are non-comparable, and $`\alpha \beta `$). Further suppose that in the straightening relation
(\*)
$$p_\tau p_\varphi =c_{\alpha \beta }p_\alpha p_\beta ,$$
the following hold:
(a) $`p_{\tau \varphi }p_{\tau \varphi }`$ occurs on the right-hand side of (\*) with coefficient $`1`$.
(b) $`\tau ,\varphi ]\beta ,\alpha [,`$ for every pair $`(\alpha ,\beta )`$ appearing on the right-hand side of $`(`$*$`)`$.
(c) Under the embedding $`๐`$, we have $`\tau \dot{}\varphi =\alpha \dot{}\beta `$, for every $`(\alpha ,\beta )`$ on the right-hand side of (\*).
Then there exists a flat family over Spec$`K[t]`$ whose special fiber ($`t=0`$) is $`SpecA()`$ and general fiber ($`t`$ invertible) is $`SpecR`$.
###### Corollary 5.2.3.
$`SpecR`$ flatly degenerates to a (normal) toric variety. In particular, $`SpecR`$ is normal and Cohen-Macaulay.
###### Proof.
We have (cf. ) that $`A()`$ is a normal domain. Hence we obtain that $`I()`$ is a binomial prime ideal. On the other hand, we have (cf. ) that a binomial prime ideal is a toric ideal (in the sense of ). It follows that $`SpecA()`$ is a (normal) toric variety and we obtain the first assertion. The first assertion together with Theorem 5.2.2 and the fact that a toric variety is Cohen-Macaulay implies that $`SpecR`$ is normal and Cohen-Macaulay. โ
### 5.3. The distributive lattice $`D`$:
Consider the partially ordered set
$$D=H\{\mathrm{๐}\}\{\mathrm{๐}\}$$
defined in ยง4.3. We equip $`D`$ with the structure of a distributive lattice by embedding it inside the chain lattice $`๐(\underset{ยฏ}{m},\underset{ยฏ}{q}):=๐(\underset{n+1\text{ times}}{\underset{}{m,m\mathrm{},m}},\underset{n+1\text{ times}}{\underset{}{q,q\mathrm{},q}})`$, as follows:
To each element of $`D`$, we associate a $`2n+2`$-tuple:
For $`A=(a_1,\mathrm{},a_r)I(r,m),B=(b_1,\mathrm{},b_r)I(r,q)`$, let $`\overline{A},\overline{B}`$ denote the $`n+1`$-tuples:
$$\overline{A}:=(a_1,\mathrm{},a_r,m,m,\mathrm{},m,1),\overline{B}:=(b_1,\mathrm{},b_r,q,q,\mathrm{},q,1)$$
(i) Let $`(A,B)H_p`$, say, $`AI(r,m),BI(r,q)`$, for some $`r,1rn`$. We let $`\overline{(A,B)}`$ be the $`(2n+2)\text{-tuple }`$: $`\overline{(A,B)}=(\overline{A},\overline{B})`$.
(ii) Let $`IH_u`$, say, $`I=(i_1,\mathrm{},i_n)(I(n,m))`$ . We let $`\stackrel{~}{I}`$ be the $`(2n+2)\text{-tuple }`$: $`\stackrel{~}{I}=(i_1,\mathrm{},i_n,1,\underset{n+1\text{ times}}{\underset{}{1,\mathrm{},1}})`$
(iii) Let $`\xi H_\xi `$, say, $`J=(j_1,\mathrm{},j_n)(I(n,m))`$), we let $`\stackrel{~}{J}`$ be the $`(2n+2)\text{-tuple }`$: $`\stackrel{~}{J}=(\underset{n+1\text{ times}}{\underset{}{1,\mathrm{},1}},j_1,\mathrm{},j_n,1)`$.
(iv) Corresponding to $`\mathrm{๐},\mathrm{๐}`$, we let $`\stackrel{~}{\mathrm{๐}},\stackrel{~}{\mathrm{๐}}`$ be the $`(2n+2)\text{-tuples }`$:
$$\stackrel{~}{\mathrm{๐}}=(\underset{n+1\text{ times}}{\underset{}{m,m\mathrm{},m}},\underset{n+1\text{ times}}{\underset{}{q,q\mathrm{},q}}),\stackrel{~}{\mathrm{๐}}=(\underset{2n+2\text{ times}}{\underset{}{1,\mathrm{},1}})$$
This induces a canonical embedding of $`D`$ inside the chain lattice $`๐(\underset{n+1\text{ times}}{\underset{}{m,m\mathrm{},m}},\underset{n+1\text{ times}}{\underset{}{q,q\mathrm{},q}})`$.
###### Lemma 5.3.1.
Let $`\tau _1,\tau _2๐(\underset{ยฏ}{m},\underset{ยฏ}{q})`$. Suppose $`\tau _1,\tau _2D`$. Then $`\tau _1\tau _2,\tau _1\tau _2`$ are also in $`D`$. Thus $`D`$ acquires the structure of a distributive lattice.
###### Proof.
Clearly the Lemma requires a proof only when $`\tau _1,\tau _2`$ are non-comparable. We consider the following cases. For two $`s`$-tuples $`E=\{e_1,\mathrm{},e_s\},F=\{f_1,\mathrm{},f_s\}`$, we shall denote
$$\begin{array}{cc}\hfill EF:& =(\mathrm{max}\{e_1,f_1\},\mathrm{},\mathrm{max}\{e_s,f_s\})\hfill \\ \hfill EF:& =(\mathrm{min}\{e_1,f_1\},\mathrm{},\mathrm{min}\{e_s,f_s\}).\hfill \end{array}$$
Case 1: $`\tau _1,\tau _2H_p`$, say $`\tau _1=(\overline{A_1},\overline{B_1}),\tau _2=(\overline{A_2},\overline{B_2})`$. We have
$$\tau _1\tau _2=(\overline{A_1}\overline{A_2},\overline{B_1}\overline{B_2}),\tau _1\tau _2=(\overline{A_1}\overline{A_2},\overline{B_1}\overline{B_2})$$
Clearly $`\tau _1\tau _2,\tau _1\tau _2`$ are in $`H_p`$, and hence in $`D`$.
Case 2: $`\tau _1H_p,\tau _2H_u`$, say $`\tau _1=(\overline{A},\overline{B}),\tau _2=\stackrel{~}{I}`$ (for some $`IH_u`$). Let $`\overline{I}`$ be the $`n+1`$-tuple $`(I,1)`$ (entries of $`I`$ followed by $`1`$). We have
$$\tau _1\tau _2=(\overline{A}\overline{I},\overline{B}),\tau _1\tau _2=(\overline{A}\overline{I},(\underset{n+1\text{ times}}{\underset{}{1,\mathrm{},1}}))$$
Clearly $`\tau _1\tau _2H_p,\tau _1\tau _2H_u`$.
Case 3: $`\tau _1H_p,\tau _2H_\xi `$, say $`\tau _1=(\overline{A},\overline{B}),\tau _2=\stackrel{~}{J}`$ (for some $`JH_\xi `$). Let $`\overline{J}`$ be the $`n+1`$-tuple $`(J,1)`$ (entries of $`I`$ followed by $`1`$). We have
$$\tau _1\tau _2=(\overline{A},\overline{B}\overline{J}),\tau _1\tau _2=(\underset{n+1\text{ times}}{\underset{}{1,\mathrm{},1}},\overline{B}\overline{J})$$
Clearly $`\tau _1\tau _2H_p,\tau _1\tau _2H_\xi `$.
Case 4: $`\tau _1,\tau _2H_u`$, say $`\tau _1=\stackrel{~}{I_1},\tau _2=\stackrel{~}{I_2}`$ (for some $`I_1,I_2H_u`$). We have
$$\tau _1\tau _2=\stackrel{~}{I_1I_2},\tau _1\tau _2=\stackrel{~}{I_1I_2}$$
Clearly $`\tau _1\tau _2,\tau _1\tau _2`$ are in $`H_u`$.
Case 5: $`\tau _1,\tau _2H_\xi `$.
This case is similar to Case 4.
Case 6: $`\tau _1H_u,\tau _2H_\xi `$, say $`\tau _1=\stackrel{~}{I},\tau _2=\stackrel{~}{J}`$ (for some $`I,J`$ in $`H_u,H_\xi `$ respectively). We have
$$\tau _1\tau _2=(\overline{I},\overline{J}),\tau _1\tau _2=\stackrel{~}{0}$$
Clearly $`\tau _1\tau _2H_p`$, $`\tau _1\tau _2D`$.
###### Lemma 5.3.2.
We have rank$`(D)=(m+q)nn^2+2(=d+2`$, where $`d=(m+q)nn^2`$). In particular, dim$`A(D)=d+3`$
This is immediate from Lemma 3.0.2.
### 5.4. Flat degeneration of Spec$`R(D)`$ to the toric variety Spec$`A(D)`$
In this subsection, we show that Spec$`R(D)`$ flatly degenerates to the toric variety Spec$`A(D)`$ by showing that $`R(D)`$ satisfies the hypotheses of Lemma 5.2.2. We first prove some preparatory Lemmas.
###### Lemma 5.4.1.
Let $`\tau ,\varphi `$ be two non-comparable elements of $`H`$. Then in the straightening relation for $`p_\tau p_\varphi `$ as given by Theorem 4.1.1, $`p_{\tau \varphi }p_{\tau \varphi }`$ occurs with coefficient $`1`$ (here for an element $`\phi `$ of $`H,p_\phi `$ stands for $`p(A,B),u(I)`$ or $`\xi (J)`$ according as $`\phi =(A,B)H_p,IH_u`$ or $`JH_\xi `$).
###### Proof.
The assertion is clear if the relation is of the type (1) of Theorem 4.1.1.
If the relation is of the type (4) of Theorem 4.1.1, then the result follows from Proposition 1.1.5,(3) (one uses the identification - as described in ยง1.5, ยง1.6 \- of $`\{p(A,B),(A,B)H_p\}`$ with the Plรผcker co-ordinates $`\{p_\tau ,\tau I(q,m+q)\}`$ restricted to the opposite cell in $`G_{q,m+q}`$).
Similarly, if the relation is of the type (2) (resp. (3)) of Theorem 4.1.1, by identifying $`M_{m,n}`$ (resp. $`M_{n,q}`$ with the opposite cell in $`G_{n,m+n}`$ (resp. $`G_{q,n+q}`$) (and using the identifications as described in ยง1.5, ยง1.6), the result follows as above (in view of Proposition 1.1.5,(3))
Let then the relation be of the type (5) or (6) of Theorem 4.1.1, say of type (5) (the proof is similar if it is of type (6)):
$`()`$
$$p(A,B)u(I)=\underset{๐ก}{}c_tp(A_t,B_t)u(I_t)$$
where $`II(n,m),(A,B)H_p`$, and $`AI`$. As in the proof of Theorem 4.1.1, (5), we multiply through out by $`\xi (I_n)`$ to arrive at
$`()`$
$$p(A,B)p(I,I_n)=a_ip(C_{i1},C_{i2})p(D_{i1},D_{i2}),a_iK^{}$$
where $`(C_{i1},C_{i2}),(D_{i1},D_{i2})`$ belong to $`H_p`$. As above, using Proposition 1.1.5,(3), we obtain that $`p((A,B)(I,I_n))p((A,B)(I,I_n))`$ occurs in (\**) with coefficient $`1`$. We have (in view of Lemma 5.3.1, rather its proof),
$$p((A,B)(I,I_n))p((A,B)(I,I_n))=p(\overline{A}\overline{I},\overline{B})p(\overline{A}\overline{I},\overline{I_n})=p(\overline{A}\overline{I},\overline{B})u(\overline{A}\overline{I})\xi (\overline{I_n})$$
Also from the proof of Theorem 4.1.1, (5), we have, for every $`i,D_{i2}=I_n`$ (in (\**)). Hence writing $`p(D_{i1},D_{i2})=u(D_{i1})\xi (I_n)`$, cancelling out $`\xi (I_n)`$ (note that L.H.S. of (\**)= $`p(A,B)u(I)\xi (I_n))`$, we obtain that $`p(\overline{A}\overline{I},\overline{B})u(\overline{A}\overline{I})`$ occurs in (\*) with coefficient $`1`$ (note that by Case 2 in the proof of Lemma 5.3.1, we have $`(A,B)I=(\overline{A}\overline{I},\overline{B}),(A,B)I=(\overline{A}\overline{I},(\underset{n+1\text{ times}}{\underset{}{1,\mathrm{},1}}))`$).
Thus the result follows if the relation is of the type (5) (or (6)) of Theorem 4.1.1.
###### Lemma 5.4.2.
Let $`\tau ,\varphi `$ be two non-comparable elements of $`D`$. Then for every $`(\alpha ,\beta )`$ on the right-hand side of the straightening relation (in $`R(D)`$, as given by Theorem 4.1.1), we have
1. $`\tau ,\varphi ]\beta ,\alpha [,`$
2. $`\tau \dot{}\varphi =\alpha \dot{}\beta `$
(here, $`\dot{}`$ denotes a disjoint union).
###### Proof.
The assertions follow from Theorem 4.1.1 (and the identification of $`D`$ as a sublattice of $`๐(\underset{ยฏ}{m},\underset{ยฏ}{q})`$).
###### Theorem 5.4.3.
There exists a flat family over $`๐ธ^1`$, with $`SpecR(D)`$ as the generic fiber and $`SpecA(D)`$ as the special fiber. In particular, $`R(D)`$ is a normal Cohen-Macaulay ring of dimension $`d+3`$ (where $`d=(m+q)nn^2`$).
###### Proof.
In view of Theorem 5.2.2, and Corollary 5.2.3, it suffices to show that (a)- (c) of Theorem 5.2.2 hold for $`R_D`$.
(a) follows from Lemma 5.4.1; (b) and (c) follow from Lemma 5.4.2.
Clearly $`R(D)`$ has dim $`d+3`$ (since dim$`A(D)=d+3`$ (cf. Lemma 5.3.2)). โ
###### Theorem 5.4.4.
The $`K`$-algebra $`S`$ is normal, Cohen-Macaulay of dimension
$`(m+q)nn^2+1`$.
###### Proof.
The algebra $`M(D)(=R(D)_{(x(\mathrm{๐}))})`$ being a homogeneous localization of the normal, Cohen-Macaulay ring $`R(D)`$, is a normal, Cohen-Macaulay ring of dim$`d+2`$.
Considering $`M(D)`$ as a graded ring (cf. ยง4.3), we have $`S(D)=M(D)_{(x(\mathrm{๐}))}`$. Hence $`S(D)`$ being a homogeneous localization of the normal, Cohen-Macaulay ring $`M(D)`$, is a normal, Cohen-Macaulay ring of dimension $`d+1`$. This together with Theorem 4.5.4 implies that $`S`$ is a normal, Cohen-Macaulay ring of dimension $`d+1`$ (note that $`d=(m+q)nn^2`$). โ
## 6. The ring of invariants $`K[X]^{SL_n(K)}`$
We preserve the notation of ยง3, ยง4. In this section, we shall show that the inclusion $`SR^{SL_n(K)}`$ is in fact an equality, i.e., $`S=R^{SL_n(K)}`$.
We now apply Lemma 2.0.4 to our situation. Let $`G=SL_n(K)`$. Consider
$$X=\underset{m\text{ copies}}{\underset{}{V\mathrm{}V}}\underset{q\text{ copies}}{\underset{}{V^{}\mathrm{}V^{}}}=SpecR,๐ธ^N=M_{m,q}(K)\times K^{\left(\genfrac{}{}{0pt}{}{m}{n}\right)}\times K^{\left(\genfrac{}{}{0pt}{}{q}{n}\right)}$$
Let $`\{u_i,\xi _j),1im,1jq,u(I),\xi (J),IH_u,JH_\xi \}`$ be denoted by $`\{f_1,\mathrm{},f_N\}`$ (note that $`f_1,\mathrm{},f_N`$ are $`G`$-invariant elements in $`R`$). Let $`x=(\underset{ยฏ}{u},\underset{ยฏ}{\xi })X`$. Let $`\psi :X๐ธ^N`$ be the map, $`\psi (x)=(f_1(x),\mathrm{},f_N(x))`$. Clearly $`\psi (X)=SpecS`$. Let us denote $`Y=SpecS`$.
###### Proposition 6.0.1.
With $`X,๐ธ^N,\psi ,Y`$ as above, the hypotheses of Lemma 2.0.4 are satisfied.
###### Proof.
(i) Let $`xX^{ss}`$. We need to show that $`\psi (x)0`$. If possible, let us assume that $`\psi (x)=0`$. Let $`x=(\underset{ยฏ}{u},\underset{ยฏ}{\xi })`$. Let $`W_u`$ (resp. $`W_\xi `$) be the span of $`\{u_1,\mathrm{},u_m\}`$ (resp.$`\{\xi _1,\mathrm{},\xi _q\}`$). Further, let dim $`W_u=r`$, dim $`W_\xi =s`$. The assumption that $`\psi (x)=0`$ implies in particular that $`u(I)(x)=0,II(n,m),\xi (J)(x)=0,JI(n,q)`$. Hence, $`W_u`$ (resp. $`W_\xi `$) is not equal to $`V`$ (resp.$`V^{}`$). Therefore, we get $`r<n,s<n`$. Also at least one of $`\{r,s\}`$ is non-zero; otherwise, $`r=0=s`$ would imply $`u_i=0,i,\xi _j=0,j`$, i.e., $`x=0`$ which is not possible, since $`xX^{ss}`$. Let us suppose that $`r0`$. (The proof is similar if $`s0`$.) The assumption that $`\psi (x)=0`$ implies in particular that $`u_i,\xi _j=0`$, for all $`i,j`$; hence, $`W_\xi W_u^{}`$. Therefore, $`snr`$. Hence we can choose a basis $`\{e_1,\mathrm{},e_n\}`$ of $`V`$ such that $`W_u=`$ the $`K`$-span of $`\{e_1,\mathrm{},e_r\}`$, and $`W_\xi `$ the $`K`$-span of $`\{e_{r+1}^{},\mathrm{},e_n^{}\}`$. Writing each vector $`u_i`$ as a row vector (with respect to this basis), we may represent the $`u`$โs by the $`m\times n`$ matrix $`๐ฐ`$ given by
$$๐ฐ:=\left(\begin{array}{ccccccc}u_{11}& u_{12}& \mathrm{}& u_{1r}& 0& \mathrm{}& 0\\ u_{21}& u_{22}& \mathrm{}& u_{2r}& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ u_{m1}& u_{m2}& \mathrm{}& u_{mr}& 0& \mathrm{}& 0\end{array}\right)$$
Similarly, writing each vector $`\xi _j`$ as a column vector (with respect to the above basis), we may represent $`\xi `$โs by the $`n\times q`$ matrix $`\mathrm{\Lambda }`$ given by
$$\mathrm{\Lambda }:=\left(\begin{array}{cccc}0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& 0& \mathrm{}& 0\\ \xi _{r+\mathrm{1\hspace{0.17em}1}}& \xi _{r+\mathrm{1\hspace{0.17em}2}}& \mathrm{}& \xi _{r+1q}\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ \xi _{n\mathrm{\hspace{0.17em}1}}& \xi _{n\mathrm{\hspace{0.17em}2}}& \mathrm{}& \xi _{nq}\end{array}\right)$$
Choose integers $`a_1,\mathrm{},a_r,b_{r+1},\mathrm{},b_n`$, all of them $`>0`$ so that $`a_i=b_j`$.
Let $`g_t`$ be the diagonal matrix in $`G(=SL_n(K)),g_t=diag(t^{a_1},\mathrm{}t^{a_r},t^{b_{r+1}},\mathrm{},t^{b_n})`$. We have, $`g_tx=g(๐ฐ,\mathrm{\Lambda })=(๐ฐg_t,g_t^1\mathrm{\Lambda })`$ (cf. ยง2.1) $`=(๐ฐ_t,\mathrm{\Lambda }_t)`$ , where
$$๐ฐ_t=\left(\begin{array}{ccccccc}t^{a_1}u_{11}& t^{a_2}u_{12}& \mathrm{}& t^{a_r}u_{1r}& 0& \mathrm{}& 0\\ t^{a_1}u_{21}& t^{a_2}u_{22}& \mathrm{}& t^{a_r}u_{2r}& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ t^{a_1}u_{m1}& t^{a_2}u_{m2}& \mathrm{}& t^{a_r}u_{mr}& 0& \mathrm{}& 0\end{array}\right)$$
and
$$\mathrm{\Lambda }_t=\left(\begin{array}{cccc}0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& 0& \mathrm{}& 0\\ t^{b_{r+1}}\xi _{r+\mathrm{1\hspace{0.17em}1}}& t^{b_{r+1}}\xi _{r+\mathrm{1\hspace{0.17em}2}}& \mathrm{}& t^{b_{r+1}}\xi _{r+1q}\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ t^{b_n}\xi _{n\mathrm{\hspace{0.17em}1}}& t^{b_n}\xi _{n\mathrm{\hspace{0.17em}2}}& \mathrm{}& t^{b_n}\xi _{nq}\end{array}\right)$$
Hence $`g_tx0`$ as $`t0`$, and this implies that $`0\overline{Gx}`$ which is a contradiction to the hypothesis that $`x`$ is semi-stable. Therefore our assumption that $`\psi (x)=0`$ is wrong and (i) of Lemma 2.0.4 is satisfied.
(ii) Let
$$U=\{(\underset{ยฏ}{u},\underset{ยฏ}{\xi })X\{u_1,\mathrm{},u_n\},\{\xi _1,\mathrm{},\xi _n\}\text{are linearly independent}\}$$
Clearly, $`U`$ is a $`G`$-stable open subset of $`X`$.
Claim : $`G`$ operates freely on $`U`$, $`UUmodG`$ is a $`G`$-principal fiber space, and $`F`$ induces an immersion $`U/G๐ธ^N`$.
Proof of Claim: Let $`H=GL_n(K)`$. We have a $`G`$-equivariant identification
$`()`$
$$UH\times H\times \underset{(mn)\text{copies}}{\underset{}{V\times \mathrm{}\times V}}\times \underset{(qn)\text{copies}}{\underset{}{V^{}\times \mathrm{}\times V^{}}}=E\times F,\mathrm{say}$$
where $`E=H\times H,F=\underset{(mn)\text{copies}}{\underset{}{V\times \mathrm{}\times V}}\times \underset{(qn)\text{copies}}{\underset{}{V^{}\times \mathrm{}\times V^{}}}`$. From this it is clear that $`G`$ operates freely on $`U`$. Further, we see that $`UmodG`$ may be identified with the fiber space with base $`(H\times H)modG`$ ($`G`$ acting on $`H\times H`$ as $`g(h_1,h_2)=(h_1g,g^1h_2),gG,h_1,h_2H`$), and fiber $`\underset{(mn)\text{copies}}{\underset{}{V\times \mathrm{}\times V}}\times \underset{(qn)\text{copies}}{\underset{}{V^{}\times \mathrm{}\times V^{}}}`$ associated to the principal fiber space $`H\times H(H\times H)/G`$. It remains to show that $`\psi `$ induces an immersion $`U/G๐ธ^N`$, i.e., to show that the map $`\psi :U/G๐ธ^N`$ and its differential $`d\psi `$ are both injective. We first prove the injectivity of $`\psi :U/G๐ธ^N`$. Let $`x,x^{}`$ in $`U/G`$ be such that $`\psi (x)=\psi (x^{})`$. Let $`\eta ,\eta ^{}U`$ be lifts for $`x,x^{}`$ respectively. Using the identification (\*) above, we may write
$$\begin{array}{c}\eta =(A,u_{n+1},\mathrm{},u_m,B,\xi _{n+1},\mathrm{},\xi _q),A,BH\\ \eta ^{}=(A^{},u_{n+1}^{},\mathrm{},u_m^{},B^{},\xi _{n+1}^{},\mathrm{},\xi _q^{}^{}),A^{},B^{}H\end{array}$$
(here, $`u_i,1in`$ are given by the rows of $`A`$, while $`\xi _i,1in`$ are given by the columns of $`B`$; similar remarks on $`u_i^{},\xi _i^{}`$). The hypothesis that $`\psi (x)=\psi (x^{})`$ implies in particular that
$$u_i,\xi _j=u_i^{},\xi _j^{},1i,jn$$
which may be written as $`AB=A^{}B^{}`$. This implies that
$`()`$
$$A^{}=Ag,$$
where $`g=BB^1(H)`$. Further, the hypothesis that $`u(I)(x)=u(I)(x^{}),I`$, implies in particular that $`u(I_n)(x)=u(I_n)(x^{})`$ (where $`I_n=(1,2,\mathrm{},n)`$). Hence we obtain
$`()`$
$$detA=detA^{}$$
Now (\**) and (\***) imply that $`g`$ in fact belongs to $`G(=SL_n(K))`$. Hence on $`U/G`$, we may suppose that
$$\begin{array}{c}x=(u_1,\mathrm{},u_n,u_{n+1},\mathrm{},u_m,\xi _1,\mathrm{},\xi _q)\\ x^{}=(u_1,\mathrm{},u_n,u_{n+1}^{},\mathrm{},u_m^{},\xi _1^{},\mathrm{},\xi _q^{})\end{array}$$
where $`\{u_1,\mathrm{},u_n\}`$ is linearly independent.
For a given $`j`$, we have,
$$u_i,\xi _j=u_i,\xi _j^{},1in,\text{implies},\xi _j=\xi _j^{}$$
(since, $`\{u_1,\mathrm{},u_n\}`$ is linearly independent). Thus we obtain
$`()`$
$$\xi _j=\xi _j^{},\text{for all }j$$
On the other hand, we have (by definition of $`U`$) that $`\{\xi _1,\mathrm{},\xi _n\}`$ is linearly independent. Hence fixing an $`i,n+1im`$, we get
$$u_i,\xi _j=u_i^{},\xi _j(=u_i^{},\xi _j^{}),1jn\text{implies},u_i=u_i^{}.$$
Thus we obtain
$`()`$
$$u_i=u_i^{},\text{for all }i$$
The injectivity of $`\psi :U/G๐ธ^N`$ follows from ($``$),($``$).
To prove that the differential d$`\psi `$ is injective, we merely note that the above argument remains valid for the points over $`K[ฯต]`$, the algebra of dual numbers ($`=KKฯต`$, the $`K`$-algebra with one generator $`ฯต`$, and one relation $`ฯต^2=0`$), i.e., it remains valid if we replace $`K`$ by $`K[ฯต]`$, or in fact by any $`K`$-algebra.
(iii) The above Claim implies in particular that dim $`U/G=`$ dim$`U`$ \- dim$`G`$ =
$`(m+q)n(n^21)`$ = dim$`SpecS`$ (cf. Theorem 5.4.4).
The condition (iv) of Lemma 2.0.4 follows from Theorem 5.4.4.
###### Theorem 6.0.2.
Let $`V=K^n`$, $`X=\underset{m\text{copies}}{\underset{}{V\mathrm{}V}}\times \underset{q\text{copies}}{\underset{}{V^{}\mathrm{}V^{}}}`$, where $`m,q>n`$. Then for the diagonal action of $`G:=SL_n(K)`$, we have
1. First Fundamental Theorem for $`\mathrm{S}\mathrm{L}_\mathrm{n}(\mathrm{K})`$-invariants: $`K[X]^G`$ is generated by $`\{p(A,B),u(I),\xi (J),(A,B)H_p,IH_u,JH_\xi \}`$.
2. Second Fundamental Theorem for $`\mathrm{S}\mathrm{L}_\mathrm{n}(\mathrm{K})`$-invariants: The ideal of relations among the generators $`\{p(A,B),u(I),\xi (J),(A,B)H_p,IH_u,JH_\xi \}`$ is generated by the six type of relations as given by Theorem 4.1.1.
3. A standard monomial basis for $`\mathrm{S}\mathrm{L}_\mathrm{n}(\mathrm{K})`$-invariants: Standard monomials in $`\{p(A,B),u(I),\xi (J),(A,B)H_p,IH_u,JH_\xi \}`$ form a $`K`$-basis for $`K[X]^G`$.
4. $`K[X]^G`$ is Cohen-Macaulay.
###### Proof.
Lemma 2.0.4 implies that $`SpecS`$ is the categorical quotient of $`X`$ by $`G`$ and $`\psi :XSpecS`$ is the canonical quotient map. Assertion (1) follows from this. Assertion (2) follows from Theorem 4.5.5. Assertion (3) follows from Theorem 4.5.3. Assertion (4) follows from Theorem 5.4.4 |
warning/0506/astro-ph0506141.html | ar5iv | text | # Resonant absorption in dissipative flux tubes
## 1 Introduction
Ionson (1978) was first to suggest that the resonant absorption of MHD waves in coronal plasmas could be a primary mechanism in coronal heating. Since then, much analytical and numerical work has been done on the subject. Rae and Roberts (1982) investigated both eikonal and differential equation approaches for the propagation of MHD waves in inhomogeneous plasmas. Hollweg (1987a,b) considered a dissipative layer of planar geometry to study the resonant absorption of coronal loops. Poedts et al. (1989, 1990) developed a finite element code to elaborate on the resonant absorption of Alfvรฉn waves in circular cylinders.
Davila (1987) and Steinolfson & Davila (1993) did much work on resonant absorption through resistivity. Ofman et al. (1994) included viscous dissipations in their analysis and concluded that the heating rate due to shear viscosity is comparable in magnitude to the resistive resonant heating. Also, they concluded that the heating caused by compressive viscosity is negligible. Goossens et al. (2002) used the TRACE data of Ofman & Aschwanden (2002) to infer the width of the inhomogeneous layer for 11 coronal loops. Ruderman & Roberts (2002) did similar analysis with the data of Nakariakov et al. (1999). Van Doorsselaere et al. (2004a) used the LEDA code to study the resistive absorption of the kink modes of cylindrical models. They concluded that, when the width of the nonuniform layer was increased, their numerical results differed by as much as 25$`\%`$ from those obtained with the analytical approximation. Van Doorsselaere et al. (2004b) investigated the effect of longitudinal curvature on quasi modes of a typical coronal loop. They found that the frequencies and damping rates of ideal quasi modes were not influenced much by the curvature. Andries et al. (2005) studied the effect of density stratification on coronal loop oscillations, and conclude that longitudinal mode numbers are coupled due to the density stratification.
In the absence of resonance, Edwin & Roberts (1983) and Roberts et al. (1984) gave a comprehensive account of the theoretical and physical concepts related to coronal oscillations. Karami et al. (2002; hereafter paper I) studied the full spectrum of MHD modes of oscillations in zero-$`\beta `$ magnetic flux tubes with discontinuous Alfvรฉn speeds at the tubeโs surface. In the vicinity of singularity, field gradients are large. Recognizing this, Sakurai et al. (1991a, b) and Goossens et al. (1992, 1995) developed a method to analyze dissipative processes in such regimes and to neglect them elsewhere.
Here we combine the two techniques of paper I and of Sakurai et al. (1991a) to obtain the resonant damping rates for the full spectrum of the normal modes of magnetic flux tubes.
## 2 Equations of motions
The linearized MHD equations for a zero-$`\beta `$, but resistive and viscous, plasma are
$$\frac{\delta ๐ฏ}{t}=\frac{1}{4\pi \rho }\left\{\left(\times \delta ๐\right)\times ๐+\left(\times ๐\right)\times \delta ๐\right\}+\frac{\eta }{\rho }^2\delta ๐ฏ,$$
(1)
$$\frac{\delta ๐}{t}=\times (\delta ๐ฏ\times ๐)+\frac{c^2}{4\pi \sigma }^2\delta ๐,$$
(2)
where $`\delta ๐ฏ`$ and $`\delta ๐`$ are the Eulerian perturbations in the velocity and the magnetic fields; $`\rho `$, $`\sigma `$, $`\eta `$, and $`c`$ are the mass density, the electrical conductivity, the viscosity and the speed of light, respectively. The simplifying assumptions are:
* Under coronal conditions gas pressure is negligible (zero- $`\beta `$).
* Density scale heights are much larger than the dimensions of flux tube, so that the gravity stratification is negligible.
* The tube geometry is a circular cylinder with coordinates ($`r`$,$`\varphi `$,$`z`$).
* There is a constant magnetic field along the $`z`$ axis, $`๐=B\widehat{z}`$.
* The equilibrium is static.
* There is no initial steady flow inside or outside of the tube.
* Viscous and resistive coefficients, $`\eta `$ and $`\sigma `$ respectively, are constants.
For a variable density, $`\rho (r)`$, a singularity develops wherever the local Alfvรฉn frequency becomes equal to the global frequency of the mode. The relevant radial wave number vanishes and resonant absorption takes place. Let us denote the radius of the tube by $`R`$ and a radius beyond which the resonance occurs by $`R_1<R`$. The thickness of the inhomogeneous layer, $`a=RR_1`$, will be assumed to be small and will be arbitrarily taken to be of the order $`R/10`$. The choice of density profile is also unimportant. We will assume two constant densities, $`\rho _i`$ in $`rR_1`$ and $`\rho _e<\rho _i`$ in $`rR`$, interconnected with a linearly varying profile in $`R_1rR`$.
In the remainder of this section the following steps are taken:
In $`r<R_1`$ and $`r>R`$, dissipative terms are neglected. Solutions of Eqs. (1) and (2) are obtained as per paper I, and their differences across the inhomogeneous layer are calculated.
In $`R_1<r<R`$, within which the resonant layer resides, solutions are found by expanding Eqs. (1) and (2) around the singular point, and the jumps across the resonant layer are found by the prescriptions of Sakurai et al. (1991a).
Interior and exterior solutions are connected by requiring the jumps in (b) to be equal to the differences in (a). This gives an analytical expression for a complex dispersion relation to be solved for the frequencies and the damping rates.
### 2.1 Interior and exterior solutions
From paper I, in the absence of dissipations, all components of $`\delta ๐ฏ`$ and the transverse components of $`\delta ๐`$ are expressible in terms of $`\delta B_z`$ only. The latter, in turn, is the solution of a second order differential equation. Thus,
$`{\displaystyle \frac{k^2}{r}}{\displaystyle \frac{d}{dr}}\left[{\displaystyle \frac{r}{k^2}}{\displaystyle \frac{d\delta B_z}{dr}}\right]+(k^2{\displaystyle \frac{m^2}{r^2}})\delta B_z=0,`$ (3)
$`\delta B_r={\displaystyle \frac{ik_z}{k^2}}{\displaystyle \frac{d\delta B_z}{dr}},\delta B_\varphi ={\displaystyle \frac{mk_z}{k^2}}{\displaystyle \frac{\delta B_z}{r}},`$ (4)
$`\delta v_r={\displaystyle \frac{i}{k^2}}{\displaystyle \frac{\omega }{B}}{\displaystyle \frac{d\delta B_z}{dr}},\delta v_\varphi ={\displaystyle \frac{m}{k^2}}{\displaystyle \frac{\omega }{B}}{\displaystyle \frac{\delta B_z}{r}},\delta v_z=0,`$ (5)
where $`k^2=\omega ^2/v_A^2k_z^2`$, and $`v_A(r)=B/\sqrt{4\pi \rho (r)}`$ is the local Alfvรฉn speed. Here, we have assumed an exponential $`\varphi ,z`$, and $`t`$ dependence, $`exp[i(m\varphi +k_zz\omega t)]`$ for any component of $`\delta ๐ฏ`$ and $`\delta ๐`$.
In the interior region, $`rR_1`$, solutions of Eq. (3) are
$$\delta B_z=\{\begin{array}{cccc}I_m\left(\left|k_i\right|r\right),& k_i^2<0,& \mathrm{surface}\mathrm{waves},& \\ J_m\left(\left|k_i\right|r\right),& k_i^2>0,& \mathrm{body}\mathrm{waves},& \\ & k_i^2=\omega ^2/v_{A_i}^2k_z^2,& & \end{array}$$
(6)
where $`J_m`$ and $`I_m`$ are Bessel and modified Bessel functions of the first kind, respectively. In the exterior region, $`rR`$, the waves should be evanescent. Solutions are
$$\delta B_z=K_m(k_er),k_e^2=k_z^2\omega ^2/v_{A_e}^2>0,$$
(7)
where $`K_m`$ is the modified Bessel function of the second kind.
### 2.2 Dispersion relation and damping
From Sakurai et al. (1991a), Goossens et al. (1992, 1995), and Erdรฉlyi et al. (1995), the jump conditions across the boundary for $`\delta B_z`$ and $`\delta v_r`$ are
$`[\delta B_z]=0,`$ (8)
$`[\delta v_r]=\pi \stackrel{~}{\omega }{\displaystyle \frac{1}{|\mathrm{\Delta }|}}{\displaystyle \frac{m^2}{\rho (r_A)r_A^2}}B_z\delta B_z,`$ (9)
where $`R_1<r_A<R`$ is the radius at which the singularity occurs and $`k^2(r_A)=0`$, $`\stackrel{~}{\omega }=\omega +i\gamma `$, and $`\mathrm{\Delta }=B^2\frac{d}{dr}(\frac{k^2}{\rho })|_{r_A}`$. Substituting the fields of Eqs. (5), (6) and (7) in jump conditions and eliminating the arbitrary amplitudes of the waves, as foreseen initially inside and outside of the boundary layer, gives the dispersion relation
$`d_0(\stackrel{~}{\omega },m)+d_1(\stackrel{~}{\omega },m)=0,`$ (10)
where
$`d_0(\stackrel{~}{\omega },m)={\displaystyle \frac{1}{k_e}}{\displaystyle \frac{K_m^{}(|k_e|R)}{K_m(|k_e|R)}}+{\displaystyle \frac{1}{k_i}}{\displaystyle \frac{J_m^{}(|k_i|R_1)}{J_m(|k_i|R_1)}},`$ (11)
$`d_1(\stackrel{~}{\omega },m)=i\pi {\displaystyle \frac{1}{|\mathrm{\Delta }|}}{\displaystyle \frac{m^2}{\rho (r_A)r_A^2}}.`$ (12)
In principle $`\stackrel{~}{\omega }=\omega +i\gamma `$ is expected to be found as a solution of Eqs. (10-12). In particular for $`\gamma \omega `$, Eq. (10) can be expanded about $`\omega `$ to give
$`\gamma ={\displaystyle \frac{\mathrm{Im}(\mathrm{d}_1(\stackrel{~}{\omega },\mathrm{m}))}{d_0(\stackrel{~}{\omega },m)/\stackrel{~}{\omega }}}|_{\stackrel{~}{\omega }=\omega }.`$ (13)
The results for surface waves are the same as Eqs. (10)- (13), except that $`J_m`$ is replaced by $`I_m`$ everywhere. This completes the formal solutions of $`\stackrel{~}{\omega }`$ and $`\gamma `$. Further analytical progress is still possible when the inhomogeneous layer is thin enough.
## 3 Thin boundary approximation
Here we assume $`(RR_1)/R=a1`$. Equation (10) reduces to $`d_00`$. For body waves, the latter is studied in ample detail in paper I. There, a trio of wave numbers $`(n,m,l)`$, corresponding to $`r`$, $`\varphi `$ and $`z`$ directions, is assigned to each mode. From Eq. (13) the corresponding damping rate becomes
$`\gamma _{nml}`$ $`=`$ $`\left\{{\displaystyle \frac{\pi m^2}{\omega _{nml}^2(\rho _i\rho _e)}}{\displaystyle \frac{a}{R^2}}\right\}/{\displaystyle \frac{d}{d\omega }}\{{\displaystyle \frac{1}{k_e}}{\displaystyle \frac{K_m^{}(|k_e|R)}{K_m(|k_e|R)}}`$ (14)
$`{\displaystyle \frac{1}{k_i}}{\displaystyle \frac{J_m^{}(|k_i|R)}{J_m(|k_i|R)}}\}.`$
Again the results for surface waves are the same as those for body waves except that $`J_m`$ is replaced by $`I_m`$. The two waves exhibit differences, for $`J_m`$ and $`I_m`$ behave differently at the boundary. We also note that each surface mode is designated by only two wave numbers, $`(m,l)`$ corresponding to $`\varphi `$ and $`z`$ directions, respectively.
For $`m=0`$, resonant absorption does not take place, because the jumps in $`\delta B_z`$ and $`\delta v_r`$ vanish and Eqs. (10\- 13) are not valid anymore.
## 4 Numerical results
As typical parameters for a coronal loop, we adopt radius = $`10^3`$ km, length = $`10^5`$ km, $`\rho _i=2\times 10^{14}`$ gr cm<sup>-3</sup>, $`\rho _e/\rho _i=0.1`$, $`B=100`$ G. For these parameters one finds $`v_{A_i}=2000`$ km s<sup>-1</sup>, $`v_{A_e}=6400`$ km s<sup>-1</sup> and $`\omega _A=2`$ rad s<sup>-1</sup>.
In Fig. 1 the magnetic field component, $`\delta B_z(r)`$, for $`n=1`$, $`m=1`$ and $`l=41,50,100`$, are plotted versus $`x=r/R`$. They are normalized to $`\mathrm{max}(\delta B_z,l=41)`$. The segment of the plot in $`0<x<1`$ is from the interior solutions of Eq. (6), while the segment in $`x>1`$ is from the exterior solutions of Eq. (7). Slopes are discontinuous at $`x=1`$. Actually, the correct location of the discontinuities is the point of singularity, $`r_A`$. Its appearance at $`x=1`$ is an artifact caused by extrapolating the interior solutions to cover the boundary layer, rather than using the exact solutions there. As $`l`$ increases, the maximum wave amplitude moves towards the tube axis and away from the inhomogeneous layer.
In Fig. 2, the frequencies and the damping rates are plotted versus $`l`$. The frequencies increase with increasing $`n`$ and/or $`l`$. For $`n=1`$ and $`2`$, damping rates exhibit maxima at $`l50`$ and $`70`$, respectively. For $`n=3`$ there is only a declining branch towards higher $`l`$ values. Our explanation for the occurrence of maxima in damping rates is the localization of the maximum amplitude of a wave within the resonant layer, where the dissipation is expected to be the highest. As Fig. 1 shows, this maximum moves away from that layer at both lower and higher $`l`$ values.
For surface waves, specified by two mode numbers $`(m,l)`$, $`\omega _{ml}`$ and $`\gamma _{ml}`$ are plotted in Figs. 3 and 4. Both frequencies and damping rates show a monotonous increase with $`l`$. Two tube radii, $`R/L=0.01,\&0.02`$, are considered here. The frequencies and damping rates of the thicker tube are almost double those of the thinner tube. This is expected, because in the limit of thin tubes, both $`\omega `$ and $`\gamma `$ are proportional to $`R/L`$ (see, e. g., Van Doorsselaere et al. 2004). The frequencies of the surface waves of Fig. 3 and the $`n=1`$ body frequencies of Fig. 2 both behave similarly. This is because of the similar behavior of $`I_m(kr)`$ and $`J_m(kr)`$ at small arguments. In physical terms this means, at least in the thin tube approximation, the surface and the $`r`$\- fundamental body waves contribute in equal manner to the heating mechanism; they both dump the bulk of their energies near the tube surface, and perhaps both are excited with comparable amplitudes and energies. The $`\omega /2\pi \gamma `$ is the number of oscillations taking place before a wave is completely attenuated. For thick and thin tubes, and $`1l10`$, this number is about $`1.7`$ and $`8`$, respectively. An observed value of Nakariakov et al. (1999) from TRACE data is about $`3.39`$. Those of Goossens et al. (2002) from the same source range from $`1.01`$\- $`3.21`$ corresponding, according to the authors, to various values of $`R/L`$, $`a/R`$, and $`\rho _e/\rho _i`$.
## 5 Summary
We studied the MHD quasi modes of coronal loops. On the assumption that ohmic and viscous dissipations are operative within a thin boundary layer, we obtained an analytic dispersion relation and solved it numerically for the mode frequencies and the damping rates. For realistic values of the initial parameters thickness- to- height ratio of the loop, the density contrast with the background medium, and the equilibrium magnetic field- our numerical values agree with those obtained from observations. As the longitudinal wave number increases, the maximum amplitude of the body eigenmodes shifts away from the resonant layer and causes a decrease in damping rates. Its behavior with increasing radial wave number is not, however, all that straightforward. In additional, we have shown that body and surface modes may contribute equally to the heating of coronal loops.
###### Acknowledgements.
The authors wish to thank Professors Robert Erdรฉlyi and, Marcel Goossens for valuable consultations. Elaborate and meticulous comments of the referee has significantly improved the content and the presentation of the paper. This work was supported by the Institute for Advanced Studies in Basic Sciences (IASBS), Zanjan. |
warning/0506/hep-ph0506086.html | ar5iv | text | # The Meaning of Higgs: ๐โบโข๐โป and ๐พโข๐พ at the Tevatron and the LHC11footnote 1MSUHEP-050608
## 1 Introduction
The origin of electroweak symmetry breaking remains unknown. While the Standard Model (SM) of particle physics is consistent with existing data, theoretical considerations suggest that this theory is only a low-energy effective theory and must be supplanted by a more complete description of the underlying physics at energies above those reached so far by experiment.
The CDF and Dร experiments at the Fermilab Tevatron are currently searching for the Higgs boson of the Standard Model. The production cross-section and decay branching fractions for this state have been predicted in great detail for the mass range accessible to Tevatron Run II. Search strategies have been carefully planned and optimized.
However, if the Tevatron does find evidence for a new scalar state, it may not necessarily be the Standard Higgs. Many alternative models of electroweak symmetry breaking have spectra that include new scalar or pseudoscalar states whose masses could easily lie in the range to which Run II is sensitive. The new scalars tend to have cross-sections and branching fractions that differ from those of the SM Higgs. The potential exists for one of these scalars to be more visible in a standard search than the SM Higgs would be.
In this paper we discuss how to extract information about non-Standard theories of electroweak symmetry breaking from searches for a light SM Higgs at Tevatron Run II and CERN LHC.
The idea of using standard Higgs searches to place limits on new scalar states associated with electroweak symmetry breaking beyond the Standard Model has been applied to LEP results (see e.g. Refs. ). The Tevatron and LHC can potentially access significantly heavier scalars than those to which LEP was sensitive, particularly in models of dynamical symmetry breaking. Ref. studied the potential of Tevatron Run II to augment its search for the SM Higgs boson by considering the process $`ggh_{SM}\tau ^+\tau ^{}`$. While this channel would not suffice as a sole discovery mode,<sup>1</sup><sup>1</sup>1The authors established that discovery of $`h_{SM}`$ in this channel alone (assuming a mass in the range 120 - 140 GeV) would require an integrated luminosity of 14-32 fb<sup>-1</sup>, which is unlikely to be achieved. the authors found that it could usefully be combined with other channels such as $`h_{SM}W^+W^{}`$ or associated Higgs production to enhance the overall visibility of the Higgs. At the same time, the authors determined what additional enhancement of scalar production and branching rate, such as might be provided in a non-standard model like the MSSM, would enable a scalar to become visible in the $`\tau ^+\tau ^{}`$ channel alone at Tevatron Run II. Similar work has been done for $`ggh_{MSSM}\tau ^+\tau ^{}`$ at the LHC and for $`ggh_{SM}\gamma \gamma `$ at the Tevatron and LHC .
Our work builds on these results, considering an additional production mechanism (b-quark annihilation), more decay channels ($`b\overline{b}`$, $`W^+W^{}`$, $`ZZ`$, and $`\gamma \gamma `$), and a wider range of non-standard physics (supersymmetry and dynamical electroweak symmetry breaking) from which rate enhancement may derive. We discuss the possible sizes of the enhancements in the various search channels for each model and pinpoint the model features having the largest influence on the degree of enhancement. We suggest the mass reach of the standard Higgs searches for each kind of non-standard scalar state. We also compare the key signals for the non-standard scalars across models and also with expectations in the SM, to show how one could start to identify which state has actually been found.
Much of our discussion will focus on the degree to which certain standard Higgs search channels are enhanced in non-standard models due to changes in the production rate or branching fractions of the non-standard scalar $`()`$ relative to the values for the standard Higgs boson $`(h_{SM})`$. We define the enhancement factor for the process $`yyxx`$ as the ratio of the products of the width of the (exclusive) production mechanism and the branching ratio of the decay:
$$\kappa _{yy/xx}^{}=\frac{\mathrm{\Gamma }(yy)\times BR(xx)}{\mathrm{\Gamma }(h_{SM}yy)\times BR(h_{SM}xx)}.$$
(1)
Analytic formulas for the decay widths of the SM Higgs boson are taken from , and numerical values are calculated using the HDECAY program .
In Section 2, we introduce supersymmetric and dynamical models of electroweak symmetry breaking and indicate which model features will be particularly relevant to our analysis. In Section 3, we discuss the production and decay of the scalar states of the various models at the Tevatron and LHC and present our results for the enhancement factors. In Section 4, we compare the different models to one another and to the SM. Section 5 holds our conclusions.
## 2 Models of Electroweak Symmetry Breaking
### 2.1 General Remarks
The Standard Higgs Model of particle physics, based on the gauge group $`SU(3)_c\times SU(2)_W\times U(1)_Y`$, accommodates electroweak symmetry breaking by including a fundamental weak doublet of scalar (โHiggsโ) bosons $`\varphi =\left(\genfrac{}{}{0pt}{}{\varphi ^+}{\varphi ^0}\right)`$ with potential function $`V(\varphi )=\lambda \left(\varphi ^{}\varphi \frac{1}{2}v^2\right)^2`$. However the SM does not explain the dynamics responsible for the generation of mass. Furthermore, the scalar sector suffers from two serious problems. The scalar mass is unnaturally sensitive to the presence of physics at any higher scale (e.g. the Planck scale), through contributions of loops of SM particles to the Higgs self-energy. This is known as the gauge hierarchy problem . In addition, if the scalar must provide a good description of physics up to arbitrarily high scale (i.e., be fundamental), the scalarโs self-coupling ($`\lambda `$) is driven to zero at finite energy scales. That is, the scalar field theory is free (or โtrivialโ) . Then the scalar cannot fill its intended role: if $`\lambda =0`$, the electroweak symmetry is not spontaneously broken. The scalars involved in electroweak symmetry breaking must therefore be a party to new physics at some finite energy scale โ e.g., they may be composite or may be part of a larger theory with a UV fixed point. The SM is merely a low-energy effective field theory, and the dynamics responsible for generating mass must lie in physics outside the SM.
In this section, we briefly introduce two classes of physics beyond the standard model that may carry the answer to the puzzle of electroweak symmetry breaking. For a review of supersymmetric models, see ,; for an introduction to dynamical electroweak symmetry breaking, see . In the meantime, we will summarize the aspects of these models which are most germane to our analysis.
### 2.2 Supersymmetry
One interesting possibility for addressing the hierarchy and triviality problems is to introduce supersymmetry. The gauge structure of the minimal supersymmetric SM (MSSM) is identical to that of the SM, but each ordinary fermion (boson) is paired with a new boson (fermion), called its โsuperpartner,โ and two Higgs doublets provide mass to all the ordinary fermions. Each loop of ordinary particles contributing to the Higgs bosonโs mass is now countered by a loop of superpartners. If the masses of the ordinary particles and superpartners are close enough, the gauge hierarchy can be stabilized . Supersymmetry relates the scalar self-coupling to gauge couplings, so that triviality is not a concern.
In order to provide masses to both up-type and down-type quarks, and to ensure anomaly cancellation, the minimal supersymmetric Standard Model (MSSM) contains two Higgs complex-doublet superfields: $`\mathrm{\Phi }_d=(\mathrm{\Phi }_d^0,\mathrm{\Phi }_d^{})`$ and $`\mathrm{\Phi }_u=(\mathrm{\Phi }_u^+,\mathrm{\Phi }_u^0)`$. When electroweak symmetry breaking occurs, the neutral components of the Higgs doublets acquire independent vacuum expectation values (vevs):
$$\mathrm{\Phi }_d=\frac{1}{\sqrt{2}}\left(\begin{array}{c}v_d\\ 0\end{array}\right),\mathrm{\Phi }_u=\frac{1}{\sqrt{2}}\left(\begin{array}{c}0\\ v_u\end{array}\right),$$
(2)
where $`\sqrt{v_d^2+v_u^2}=2M_W/g=246`$ GeV. Out of the original 8 degrees of freedom, 3 serve as Goldstone bosons, absorbed into longitudinal components of the $`W^\pm `$ and $`Z`$, making them massive. The other 5 degrees of freedom remain in the spectrum as distinct scalar states, namely two neutral, CP-even states
$`h`$ $`=`$ $`(\sqrt{2}\text{Re }\mathrm{\Phi }_d^0v_d)\mathrm{sin}\alpha +(\sqrt{2}\text{Re }\mathrm{\Phi }_u^0v_u)\mathrm{cos}\alpha ,`$ (3)
$`H`$ $`=`$ $`(\sqrt{2}\text{Re }\mathrm{\Phi }_d^0v_d)\mathrm{cos}\alpha +(\sqrt{2}\text{Re }\mathrm{\Phi }_u^0v_u)\mathrm{sin}\alpha ,`$ (4)
one neutral, CP-odd state
$$A=\sqrt{2}(\text{Im }\mathrm{\Phi }_d^0\mathrm{sin}\beta +\text{Im }\mathrm{\Phi }_u^0\mathrm{cos}\beta ),$$
(5)
and a charged pair
$$H^\pm =\mathrm{\Phi }_d^\pm \mathrm{sin}\beta +\mathrm{\Phi }_u^\pm \mathrm{cos}\beta .$$
(6)
Here $`\alpha `$ is the mixing angle between $`h`$ and $`H`$ which diagonalizes the neutral boson mass-squared matrix:
$$_0^2=\left(\begin{array}{cc}M_A^2\mathrm{sin}^2\beta +M_Z^2\mathrm{cos}^2\beta & (M_A^2+M_Z^2)\mathrm{sin}\beta \mathrm{cos}\beta \\ (M_A^2+M_Z^2)\mathrm{sin}\beta \mathrm{cos}\beta & M_A^2\mathrm{cos}^2\beta +M_Z^2\mathrm{sin}^2\beta \end{array}\right),$$
(7)
and $`\beta `$ is defined through the ratio $`v_u/v_d`$ (sometimes denoted as $`v_2/v_1`$ )
$$\mathrm{tan}\beta =v_u/v_d.$$
(8)
It is conventional to choose $`\mathrm{tan}\beta `$ and
$$M_A=\sqrt{M_{H^\pm }^2M_W^2}$$
(9)
to define the SUSY Higgs sector. From the above equations one may derive the relations
$`M_{h,H}^2={\displaystyle \frac{1}{2}}\left[(M_A^2+M_Z^2)\sqrt{(M_A^2+M_Z^2)^24M_A^2M_Z^2\mathrm{cos}^22\beta }\right],`$ (10)
$`\mathrm{cos}^2(\beta \alpha )={\displaystyle \frac{M_h^2(M_Z^2M_h^2),}{M_A^2(M_H^2M_h^2)}}`$ (11)
which will be useful for determining when Higgs boson interactions with fermions are enhanced.
The Yukawa interactions of the Higgs fields with the quarks and leptons are given by:
$`_{\mathrm{Yukawa}}=`$ $`h_u\left[\overline{u}P_Lu\mathrm{\Phi }_u^0\overline{u}P_Ld\mathrm{\Phi }_u^+\right]+h_d\left[\overline{d}P_Ld\mathrm{\Phi }_d^0\overline{d}P_Lu\mathrm{\Phi }_d^{}\right]`$ (12)
$`+h_{\mathrm{}}\left[\overline{\mathrm{}}P_L\mathrm{}\mathrm{\Phi }_d^0\overline{\mathrm{}}P_L\nu \mathrm{\Phi }_d^{}\right]+\mathrm{h}.\mathrm{c}.`$
Using Eq. (2) and Eq. (12) we find, for example, for the 3rd generation:
$`h_t=`$ $`{\displaystyle \frac{\sqrt{2}m_t}{v_u}}={\displaystyle \frac{\sqrt{2}m_t}{v\mathrm{sin}\beta }},`$ (13)
$`h_{b,\tau }=`$ $`{\displaystyle \frac{\sqrt{2}m_{b,\tau }}{v_d}}={\displaystyle \frac{\sqrt{2}m_{b,\tau }}{v\mathrm{cos}\beta }}.`$ (14)
To display this in terms of the interactions of the mass eigenstate Higgs bosons with the fermions ($`Y_{f\overline{f}}`$) we may write<sup>2</sup><sup>2</sup>2Note that the interactions of the $`A`$ are pseudoscalar, i.e. it couples to $`\overline{\psi }\gamma _5\psi `$.
$`Y_{ht\overline{t}}/Y_{ht\overline{t}}^{SM}`$ $`=`$ $`\mathrm{cos}\alpha /\mathrm{sin}\beta Y_{hb\overline{b}}/Y_{hb\overline{b}}^{SM}=\mathrm{sin}\alpha /\mathrm{cos}\beta `$
$`Y_{Ht\overline{t}}/Y_{ht\overline{t}}^{SM}`$ $`=`$ $`\mathrm{sin}\alpha /\mathrm{sin}\beta Y_{Hb\overline{b}}/Y_{hb\overline{b}}^{SM}=\mathrm{cos}\alpha /\mathrm{cos}\beta `$ (15)
$`Y_{At\overline{t}}/Y_{ht\overline{t}}^{SM}`$ $`=`$ $`\mathrm{cot}\beta Y_{Ab\overline{b}}/Y_{hb\overline{b}}^{SM}=\mathrm{tan}\beta `$
relative to the Yukawa couplings of the Standard Model ($`Y_{hf\overline{f}}^{SM}=m_f/v)`$. Once again, the same pattern holds for the tau leptonโs Yukawa couplings as for those of the $`b`$ quark.
There are several circumstances under which various Yukawa couplings are enhanced relative to Standard Model values. For high $`\mathrm{tan}\beta `$ (small $`\mathrm{cos}\beta `$), eqns. (15) show that the interactions of all neutral Higgs bosons with the down-type fermions are enhanced by a factor of $`1/\mathrm{cos}\beta `$. In the decoupling limit, where $`M_A\mathrm{}`$, applying eqns. (10) and (11) to eqns. (15) shows that the $`H`$ and $`A`$ Yukawa couplings to down-type fermions are enhanced by a factor of $`\mathrm{tan}\beta `$
$`Y_{Hb\overline{b}}/Y_{hb\overline{b}}^{SM}=Y_{H\tau \overline{\tau }}/Y_{h\tau \overline{\tau }}^{SM}\mathrm{tan}\beta ,`$ (16)
Conversely, for low $`m_Am_h`$, one can check that
$`Y_{hb\overline{b}}/Y_{hb\overline{b}}^{SM}=Y_{h\tau \overline{\tau }}/Y_{h\tau \overline{\tau }}^{SM}\mathrm{tan}\beta `$ (17)
that $`h`$ and $`A`$ Yukawas are enhanced instead. For further details we refer to Ref. where issues of mass-degenerate Higgs bosons in MSSM at large $`\mathrm{tan}\beta `$ have been studied in great detail.
### 2.3 Technicolor
Another intriguing class of theories, dynamical electroweak symmetry breaking (DEWSB), supposes that the scalar states involved in electroweak symmetry breaking could be manifestly composite at scales not much above the electroweak scale $`v250`$ GeV. In these theories, a new asymptotically free strong gauge interaction (technicolor ) breaks the chiral symmetries of massless fermions $`f`$ at a scale $`\mathrm{\Lambda }1`$ TeV. If the fermions carry appropriate electroweak quantum numbers (e.g. left-hand (LH) weak doublets and right-hand (RH) weak singlets), the resulting condensate $`\overline{f}_Lf_R0`$ breaks the electroweak symmetry as desired. Three of the Nambu-Goldstone Bosons (technipions) of the chiral symmetry breaking become the longitudinal modes of the $`W`$ and $`Z`$. The logarithmic running of the strong gauge coupling renders the low value of the electroweak scale natural. The absence of fundamental scalars obviates concerns about triviality.
Many models of DEWSB have additional light neutral pseudo Nambu-Goldstone bosons which could potentially be accessible to a standard Higgs search; these are called โtechnipionsโ in technicolor models. There is not one particular DEWSB model that has been singled out as a benchmark, in the manner of the MSSM among supersymmetric theories. Rather, several different classes of models have been proposed to address various challenges within the DEWSB paradigm of the origins of mass. In this paper, we look at several representative technicolor models. We both evaluate the potential of standard Higgs searches to discover the lightest Pseudo Nambu-Goldstone Boson (PNGB) of each of these models, and also draw some inferences about the characteristics of technicolor models that have the greatest impact on this search potential.
Our analysis will assume, for simplicity, that the lightest PNGB state is significantly lighter than other neutral (pseudo) scalar technipions, so as to heighten the comparison to the SM Higgs boson. The precise spectrum of any technicolor model generally depends on a number of parameters, particularly those related to whatever โextended technicolorโ interaction transmits electroweak symmetry breaking to the ordinary quarks and leptons. Models in which several light neutral PNGBs were nearly degenerate would produce even larger signals than those discussed here.
The specific models we examine are: 1) the traditional one-family model with a full family of techniquarks and technileptons, 2) a variant on the one-family model in which the lightest technipion contains only down-type technifermions and is significantly lighter than the other pseudo Nambu-Goldstone bosons, 3) a multiscale walking technicolor model designed to reduce flavor-changing neutral currents, and 4) a low-scale technciolor model (the Technicolor Straw Man model) with many weak doublets of technifermions, in which the second-lightest technipion $`P^{}`$ is the state relevant for our study (the lightest, being composed of technileptons, lacks the anomalous coupling to gluons required for $`ggP`$ production). For simplicity the lightest relevant neutral technipion of each model will be generically denoted $`P`$; where a specific model is meant, a superscript will be used.
One of the key differences among these models is the value of the technipion decay constant $`F_P`$, which is related to the number $`N_D`$ of weak doublets of technifermions that contribute to electroweak symmetry breaking. In a theory like model 2, in which only a single technifermion condensate breaks the electroweak symmetry, the value of $`F_P`$ is simply the weak scale: $`F_P^{(2)}=v=246`$ GeV. In models where more than one technifermion condensate breaks the EW symmetry, one finds $`v^2=f_P^2+f_2^2+f_3^2+\mathrm{}`$ For example, in the one-family model (model 1), all four technidoublets corresponding to a technifermion โgenerationโ condense, so that the decay constant is fixed to be $`F_P^{(1)}=\frac{v}{2}.`$ In the lowscale model (model 4), the number of condensing technidoublets is much higher, of order 10; setting $`N_D=10`$ yields $`F_P^{(4)}=\frac{v}{\sqrt{10}}`$. In the multiscale model (model 3), the scales at which various technicondensates form are assumed to be significantly different, so that the lowest scale is simply bounded from above. In keeping with and to ensure that the technipion mass will be in the range to which the standard Tevatron Higgs searches are sensitive, we set $`F_P^{(3)}`$ = $`\frac{v}{4}`$.
In section 3, we study the enhancement factors for several production and decay modes of the lightest PNGBs of each technicolor model. Then in section 4, we compare the signatures of these PNGBs to those of a SM Higgs and the Higgs bosons of the MSSM in order to determine how the standard search modes (or additional channels) can help tell these states apart.
## 3 Results For Each Model
In this section, we examine the single production of SUSY Higgses and technicolor PNGBs via the two dominant methods at the Tevatron and LHC: gluon fusion and $`b\overline{b}`$ annihilation. We determine the degree to which these production channels are enhanced relative to production of a SM Higgs, and find which channel dominates for each scalar state. We likewise study the major decay modes: $`b\overline{b}`$, $`\tau ^+\tau ^{}`$, $`\gamma \gamma `$, and $`W^+W^{}`$ in order to determine the branching fractions relative to those of an SM Higgs. We then combine this information to obtain the overall enhancement factor in each channel and the estimated cross-section at each collider.
### 3.1 Supersymmetry
#### 3.1.1 Factors affecting signal strength
Let us consider how the signal of a light Higgs boson could be changed in the MSSM, compared to expectations in the SM. There are several important sources of alterations in the predicted signal, some of which are interconnected.
First, the MSSM includes three neutral Higgs bosons $`=(h,H,A)`$ states. The apparent signal of a single light Higgs could be enhanced if two or three neutral Higgs species are nearly degenerate so that more than one Higgs is actually contributing to the final state being studied. The left-hand frame of Figure 1 illustrates that for Higgs masses around 120 GeV it is possible for several Higgs states to be close in mass. We take advantage of this near-degeneracy by combining the signals of the different neutral Higgs bosons when their masses are closer than the experimental resolution. Specifically, when combining the signal from $`A`$, $`h`$, and $`H`$, we require $`|M_AM_h|`$ and/or $`|M_AM_H|`$ to be less than 0.3$`\sqrt{M_A/GeV}`$ GeV, as compared to the approximate experimental resolution for the Higgs mass of $`\sqrt{M_A/GeV}`$ GeV for $`\tau ^+\tau ^{}`$ or $`b\overline{b}`$ channels. For the Higgs mass range studied here, 0.3$`\sqrt{M_A/GeV}`$ would correspond to a fairly small mass gap of order $`35`$ GeV. For the $`\gamma \gamma `$ channel we do not combine the $`(h,H,A)`$ states but use just one, the $`A\gamma \gamma `$ process, since the experimental mass resolution for this final state could be of the order of one GeV.
Second, the alterations of the couplings between Higgs bosons and ordinary fermions in the MSSM, which were discussed in section 2.2, can change the Higgs decay widths and branching ratios relative to those in the SM. The SM branching fractions are pictured in the right-hand frame of Figure 1 and those in the MSSM (as calculated with the HDECAY<sup>3</sup><sup>3</sup>3 For SUSY HDECAY input we have chosen the squark masses to be 1 TeV, while the trilinear $`A_{t,b,\tau }`$ and $`\mu `$ parameters were taken equal to 200 GeV. program ) are in Figures 2, 3, 4, and the relevant branching ratios for a 130 GeV CP-odd Higgs are given for various $`\mathrm{tan}\beta `$ in Table 1. These changes directly effect the enhancement factor for a given process, as in equation (1). When radiative effects on the masses and couplings are included, the Higgs boson production rate as well as the decay branching fractions can be substantially altered, in a non-universal way. For instance, $`B(h\tau ^+\tau ^{})`$ could be enhanced by up to an order of magnitude due to the suppression of $`B(hb\overline{b})`$ in certain regions of parameter space . However, this gain in branching fraction would be offset<sup>4</sup><sup>4</sup>4 There can be a suppression of $`BR(Hb\overline{b})`$ and $`BR(H\tau \tau )`$ in the parameter region where all Higgs bosons are nearly degenerate . to some degree by a reduction in Higgs production through channels involving $`Y_{b\overline{b}}`$.
Third, recall that SM production of the light Higgs via gluon fusion is dominated by a top-quark loop; the large top quark mass both increases the top-Higgs coupling and suppresses the loop. In the MSSM, a large value of $`\mathrm{tan}\beta `$ enhances the bottom-Higgs coupling (eqns. (16) and (17)), making gluon fusion through a $`b`$-quark loop significant, and possibly even dominant over the top-quark loop contribution.
Fourth, the presence of superpartners in the MSSM gives rise to new squark-loop contributions to Higgs boson production through gluon fusion. Light squarks with masses of order 100 GeV have been argued to lead to a considerable universal enhancement (as much as a factor of five) for MSSM Higgs production compared to the SM.
Finally, enhancement of the $`Y_{b\overline{b}}`$ coupling at moderate to large $`\mathrm{tan}\beta `$ makes $`b\overline{b}`$ a significant means of Higgs production in the MSSM โ in contrast to the SM where it is negligible. To include both production channels when looking for a Higgs decaying as $`xx`$, we define a combined enhancement factor
$`\kappa _{total/xx}^{}`$ $`=`$ $`{\displaystyle \frac{\sigma (ggxx)+\sigma (bbxx)}{\sigma (ggh_{SM}xx)+\sigma (bbh_{SM}xx)}}`$ (18)
$`=`$ $`{\displaystyle \frac{\kappa _{gg/xx}^{}+\sigma (bbxx)/\sigma (ggh_{SM}xx)}{1+\sigma (bbh_{SM}xx)/\sigma (ggh_{SM}xx)}}`$
$`=`$ $`{\displaystyle \frac{\kappa _{gg/xx}^{}+\kappa _{bb/xx}^{}\sigma (bbh_{SM}xx)/\sigma (ggh_{SM}xx)}{1+\sigma (bbh_{SM}xx)/\sigma (ggh_{SM}xx)}}`$
$``$ $`[\kappa _{gg/xx}^{}+\kappa _{bb/xx}^{}R_{bb:gg}]/[1+R_{bb:gg}].`$
Here $`R_{bb:gg}`$ is the ratio of $`b\overline{b}`$ and $`gg`$ initiated Higgs boson production in the Standard Model, which can be calculated using HDECAY.
#### 3.1.2 Enhancement Factors and Cross-sections
Figure 5 (6) presents NLO cross sections at the Tevatron (LHC). For $`b\overline{b}`$ we are using the code of Ref. , <sup>5</sup><sup>5</sup>5 Note that $`b\overline{b}`$ has been recently calculated at NNLO in . while for $`gg`$ we use HIGLU and HDECAY .<sup>6</sup><sup>6</sup>6 Specifically, we use the HIGLU package to calculate the $`ggh_{sm}`$ cross section. We then use the ratio of the Higgs decay widths from HDECAY (which includes a more complete set of one-loop MSSM corrections than HIGLU) to get the MSSM $`gg`$ cross section: $`\sigma ^{MSSM}=\sigma ^{SM}\times \mathrm{\Gamma }(gg)/\mathrm{\Gamma }(h_{SM}gg)`$. Frame (a) shows production of $`h_{SM}`$; frames (b)-(d) show production of the MSSM axial Higgs for several values of $`\mathrm{tan}\beta `$. One can see that in the MSSM the contribution from $`b\overline{b}`$ becomes important even for moderate values of $`\mathrm{tan}\beta 10`$. For $`M_{}<110115`$ GeV the contribution from $`gg`$ process is a bit bigger than that from $`b\overline{b}`$, while for $`M_{}>115`$ GeV $`b`$-quark-initiated production begins to outweigh gluon-initiated production.
Using the Higgs branching fractions from above with these NLO cross sections for $`ggH`$ and $`b\overline{b}H`$ allows us to derive $`\kappa _{total/xx}^{}`$, as presented in Fig. 7 for the Tevatron and LHC. Several comments are in order. In Fig. 7(a) one can see a gap in enhancement factor for $`WW`$ and $`ZZ`$ final states at $`\mathrm{tan}\beta =5`$ for $`M_A`$ between 90-130 GeV. This is related to our procedure of combining signals from $`(A,h,H)`$ bosons. The $`A`$-boson does not couple to $`WW`$ or $`ZZ`$, while the mass gap between $`h`$ and $`H`$ is too big at low values of $`\mathrm{tan}\beta `$ to satisfy our combination criterion ($`|M_AM_{H,h}|<0.3\sqrt{M_A/GeV}`$), so one cannot define an enhancement factor for this parameter region. At higher values of $`\mathrm{tan}\beta `$ there is no corresponding gap for $`WW`$ and $`ZZ`$ final states for $`M_A`$ between 90-130 GeV, however one can observe artificial peaks for $`M_A`$ between 90-130 GeV which are again related to our combination procedure. In addition, there are several โphysicalโ kinks and peaks in the enhancement factor for various Higgs boson final states related to $`WW`$, $`ZZ`$ and top-quark thresholds which can be seen for the respective values of $`M_A`$. At very large values of $`\mathrm{tan}\beta `$ the top-quark threshold effect for the $`\gamma \gamma `$ enhancement factor is almost gone because the b-quark contribution dominates in the loop.
The enhancement factors and cross sections for a 130 GeV CP-odd Higgs are listed, for various values of $`\mathrm{tan}\beta `$, in Table 2. From Table 2 one can see that the enhancement factors at the Tevatron and LHC are very similar. On the other hand, the values of the total rates at the LHC are about two orders of magnitude higher than the corresponding rates at the Tevatron. One should also notice that enhancements of the $`b\overline{b}`$ and $`\tau ^+\tau ^{}`$ signatures are very similar and they rise swiftly by a factor of 200 as $`\mathrm{tan}\beta `$ increases from 5 to 50. In contrast, the $`\gamma \gamma `$ signature is always strongly suppressed! This particular feature of SUSY models, as we will see below, may be important for distinguishing supersymmetric models from models with dynamical symmetry breaking.
It is important to note that combining the signal from the neutral Higgs bosons $`h,A,H`$ in the MSSM turns out to make our results more broadly applicable across SUSY parameter space. As discussed earlier, Figure 1(left) reveals that at moderate-to-high values of $`\mathrm{tan}\beta `$ at least two of the neutral Higgs bosons are degenerate in mass <sup>7</sup><sup>7</sup>7The degenerate pair is either $`(h,A)`$ for $`M_A<M_A^0`$ or $`(H,A)`$ for $`M_A>M_A^0`$ , where the value of $`M_A^0`$ is related to the maximal mass of the light Higgs as a function of $`M_A`$ with other SUSY parameters held fixed. See Figure 1.. The value of the light Higgs mass $`M_h`$ also depends on the degree of mixing between the scalar partners of the top quark; this is parameterized by the variable $`A_t`$. For a given SUSY scale, $`M_S`$, the mass $`M_h`$ takes its maximum value for $`X_tA_t\mu \mathrm{cot}\beta =\sqrt{6}M_S`$ which corresponds to the โmaximal mixing caseโ while for $`X_tA_t\mu \mathrm{cot}\beta =0`$ we have the โminimal mixing caseโ and $`M_h`$ takes on its minimum value. What is interesting is that although the value of $`M_h`$ can differ significantly in the minimal and maximal mixing cases, the combined signal from all 3 Higgses at high $`\mathrm{tan}\beta `$ leads to the nearly the same (within at most few percent) enhancement factor, as shown in Figure 8. Combining the signals from $`A,h,H`$ has the virtue of making the enhancement factor independent of the degree of top squark mixing (for fixed $`M_A`$, $`\mu `$ and $`M_S`$ and medium to high values of $`\mathrm{tan}\beta `$), which greatly reduces the parameter-dependence of our results.
### 3.2 Technicolor
#### 3.2.1 PNGB Production via Gluon Fusion
Single production of a technipion can occur through the axial-vector anomaly which couples the technipion to pairs of gauge bosons. For an $`SU(N_{TC})`$ technicolor group with technipion decay constant $`F_P`$, the anomalous coupling between the technipion and a pair of gauge bosons is given, in direct analogy with the coupling of a QCD pion to photons,<sup>8</sup><sup>8</sup>8Note that the normalization used here differs from that used in by a factor of 4. by
$$N_{TC}๐_{V_1V_2}\frac{g_1g_2}{8\pi ^2F_P}ฯต_{\mu \nu \lambda \sigma }k_1^\mu k_2^\nu ฯต_1^\lambda ฯต_2^\sigma $$
(19)
where $`๐_{V_1V_2}`$ is the anomaly factor, $`g_i`$ are the gauge boson couplings, and the $`k_i`$ and $`ฯต_i`$ are the four-momenta and polarizations of the gauge bosons. The values of the anomaly factors for the lightest PNGB coupling to gluons are given in Table 3 for each model.
The rate of single technipion production in this channel is proportional to the decay width to gluons. In the technicolor models, we have
$$\mathrm{\Gamma }(Pgg)=\frac{m_P^3}{8\pi }\left(\frac{\alpha _sN_{TC}๐_{gg}}{2\pi F_P}\right)^2.$$
(20)
while in the SM, the expression looks like
$$\mathrm{\Gamma }(hgg)=\frac{m_h^3}{8\pi }\left(\frac{\alpha _s}{3\pi v}\right)^2,$$
(21)
in the heavy top-quark approximation. Comparing a PNGB to a SM Higgs boson of the same mass, we find the enhancement in the gluon fusion production rate is
$$\kappa _{ggprod}=\frac{\mathrm{\Gamma }(Pgg)}{\mathrm{\Gamma }(hgg)}=\frac{9}{4}N_{TC}^2๐_{gg}^2\frac{v^2}{F_P^2}$$
(22)
The main factors influencing $`\kappa _{ggprod}`$ for a fixed value of $`N_{TC}`$ are the anomalous coupling to gluons and the technipion decay constant. The value of $`\kappa _{ggprod}`$ for each model (taking $`N_{TC}=4`$) is given in Table 4.
#### 3.2.2 Production via $`b\overline{b}`$ annihilation
The PNGBs couple to b-quarks courtesy of the extended technicolor interactions responsible for producing masses for the ordinary quarks and leptons. The extended technicolor group (of which $`SU(N_{TC})`$ is an unbroken subgroup) includes gauge bosons that couple to both ordinary and technicolored fermions so that the ordinary fermions can interact with the technicondensates that break the electroweak symmetry.
The rate of technipion production via $`b\overline{b}`$ annihilation is proportional to $`\mathrm{\Gamma }(Pb\overline{b})`$. In general, the expression for the decay of a technipion to fermions is
$$\mathrm{\Gamma }(Pf\overline{f})=\frac{N_C\lambda _f^2m_f^2m_P}{8\pi F_P^2}\left(1\frac{4m_f^2}{m_P^2}\right)^{\frac{s}{2}}$$
(23)
where $`N_C`$ is 3 for quarks and 1 for leptons. The phase space exponent, $`s`$, is 3 for scalars and 1 for pseudoscalars; the lightest PNGB in models 1 and 4 is a scalar, while in models 2 and 3 it is assumed to be a pseudoscalar. For the technipion masses considered here, the value of the phase space factor in (23) is so close to one that the value of $`s`$ makes no practical difference. The factor $`\lambda _f`$ is a non-standard Yukawa coupling distinguishing leptons from quarks. Model 2 has $`\lambda _{quark}=\sqrt{\frac{2}{3}}`$ and $`\lambda _{lepton}=\sqrt{6}`$; model 3 also includes a similar factor, but with average value 1; $`\lambda _f=1`$ in models 1 and 4. Finally, it should be noted that model 2 assumes that the lightest technipion is composed only of down-type fermions and cannot decay to $`c\overline{c}`$; since this decay would usually have a small branching ratio and $`c\overline{c}`$ is not a preferred final state for Higgs searches, this has little impact.
For comparison, the decay width of the SM Higgs into b-quarks is:
$$\mathrm{\Gamma }(hb\overline{b})=\frac{3m_b^2m_h}{8\pi v^2}\left(1\frac{4m_b^2}{m_h^2}\right)^{\frac{3}{2}}$$
(24)
The production enhancement for $`b\overline{b}`$ annihilation is (again assuming Higgs and technipion have the same mass):
$$\kappa _{bbprod}=\frac{\mathrm{\Gamma }(Pb\overline{b})}{\mathrm{\Gamma }(hb\overline{b})}=\frac{\lambda _b^2v^2}{F_P^2}\left(1\frac{4m_b^2}{m_h^2}\right)^{\frac{s3}{2}}$$
(25)
The value of $`\kappa _{bbprod}`$ (shown in Table 4) is controlled by the size of the technipion decay constant.
We see from Table 4 that $`\kappa _{bbprod}`$ is at least one order of magnitude smaller than $`\kappa _{ggprod}`$ in each model. Taking the ratio of equations (22) and (25)
$$\frac{\kappa _{ggprod}}{\kappa _{bbprod}}=\frac{9}{4}N_{TC}^2๐_{gg}^2\lambda _b^2\left(1\frac{4m_b^2}{m_h^2}\right)^{\frac{3s}{2}}$$
(26)
we see that the larger size of $`\kappa _{ggprod}`$ is due to the factor of $`N_{TC}^2`$ coming from the fact that gluons couple to a technipion via a techniquark loop. The extended technicolor (ETC) interactions coupling $`b`$-quarks to a technipion have no such enhancement.
In addition, the production cross-section for a SM Higgs boson via $`b\overline{b}`$ annihilation is 2 to 3 orders of magnitude smaller than that for gluon fusion at the Tevatron and LHC . With a smaller SM cross-section and a smaller enhancement factor, it is clear that technipion production via $`b\overline{b}`$ annihilation is essentially negligible at these hadron colliders. Nonetheless, to be conservative, we include the $`b\overline{b}`$ production channel because it tends to slightly reduce the production enhancement factor.
Using the combined enhancement factor definition of eqn. (18), and recalling that $`\kappa _{dec}`$ is the same for both colliders, we find that the small differences due to the values of $`R_{bb:gg}`$ do not give a noticable difference between the values of $`\kappa _{total/xx}^P`$ at the Tevatron and LHC; the production enhancement factors quoted in Table 4 apply to both colliders.
#### 3.2.3 Decays
The decay width of a light technipion into gluons or fermion/anti-fermion pairs has been discussed above. Since the technipions we are studying do not decay to $`W`$ bosons and their decay to $`Z`$ bosons through the axial vector anomaly is negligible in the interesting mass range, the remaining possibility is a decay to photons. Again, this proceeds through the axial vector anomaly (cf. eqn. (19)) and the anomaly factors $`๐_{\gamma \gamma }`$ are shown in Table 3.
We now calculate the technipion branching ratios from the above information, taking $`N_{TC}=4`$. The values are essentially independent of the size of $`M_P`$ within the range 120 GeV - 160 GeV; the branching fractions for $`M_P=130`$ GeV are shown in Table 5. The branching ratios for the SM Higgs at NLO are given for comparison; they were calculated using HDECAY . Note that, in contrast to the technipions, a SM Higgs in this mass range already has a noticeable decay rate to off-shell vector bosons.
Comparing the technicolor and SM branching ratios in Table 5, we see immediately that all decay enhancements, except to the $`gg`$ mode, are generally of order one and therefore much smaller than the production enhancements. Decays to $`b\overline{b}`$ are slightly enhanced, if at all. Decays to $`c\overline{c}`$ are enhanced in our tree-level calculations โ but note that it is higher-order corrections that suppress this mode for the SM Higgs; in any case, this is not a primary discovery channel. Decays to $`\tau `$ leptons are slightly suppressed in general; again, the comparison of tree-level technicolor and loop-level SM Higgs calculations may be a factor here. Model 2 is an exception; its unusual Yukawa couplings yield a decay enhancement in the $`\tau ^+\tau ^{}`$ channel of order the technipionโs (low) production enhancement. In the $`\gamma \gamma `$ channel, the decay enhancement strongly depends on the group-theoretical structure of the model, through the anomaly factor. Table 6 includes the decay enhancements $`\kappa _{dec}^P`$ for the most experimentally promising search channels.
#### 3.2.4 Enhancement Factors and Cross-Sections
Our results for the Tevatron Run II and LHC production enhancements (including both $`gg`$ fusion and $`b\overline{b}`$ annihilation), decay enhancements, and overall enhancements of each technicolor model relative to the SM are shown in Table 6 for a technipion or Higgs mass of 130 GeV. Multiplying $`\kappa _{tot/xx}^P`$ by the cross-section for SM Higgs production via gluon fusion yields an approximate technipion production cross-section, as shown in the right-most column of Table 6.
In each technicolor model, the main enhancement of the possible technipion signal relative to that of an SM Higgs arises at production, making the size of the technipion decay constant the most critical factor in determining the degree of enhancement for fixed $`N_{TC}`$.
Each decay enhancement is in general of order 1, making it significantly smaller than the typical production enhancement. In model 3 where the decay โenhancementโ is actually a suppression, the decay factor is 3 orders of magnitude smaller than the production enhancement. We find that $`Pb\overline{b}`$ is very similar to $`h_{SM}b\overline{b}`$. The decay $`P\tau ^+\tau ^{}`$ generally has a suppressed rate relative to SM expectations; again, this may relate to comparing leading technicolor and NLO SM results. An exception is model 2, where the special structure of the Yukawa coupling leads to a $`\tau ^+\tau ^{}`$ decay enhancement of the same order as the production enhancement. The $`P\gamma \gamma `$ decay enhancement factor depends strongly on the group-theoretic structure of the model through the anomaly factor, ranging from a distinct enhancement in model 4 to a factor-of-10 suppression in model 1.
## 4 Interpretation
We are ready to put our results in context. The large QCD background for $`q\overline{q}`$ states of any flavor makes the tau-lepton-pair and di-photon final states the most promising for exclusion or discovery of the Higgs-like states of the MSSM or technicolor. We now illustrate how the size of the enhancement factors for these two final states vary over the parameter spaces of these theories at the Tevatron and LHC. We use this information to display the likely reach of each experiment in each of these standard Higgs search channels. Then, we compare the signatures of the MSSM Higgs bosons and the various technipions to see how one might tell these states apart from one another.
### 4.1 Visibility of MSSM Higgs Bosons
The left-hand frame of Figure 9 displays contours of enhancement factors of 2, 10, 100 and 1000 for the process $`gg+b\overline{b}h+A+H\tau ^+\tau ^{}`$ in the MSSM at the Tevatron. We see that the enhancement factors grow dramatically as either $`\mathrm{tan}\beta `$ or $`M_A`$ becomes large. These results are consistent with those of . The large increase in the enhancement factor for large values of $`M_A`$ takes place because the Standard Model $`Br[H\tau \tau (b\overline{b})]`$ decreases sharply when the $`WW`$ and $`ZZ`$ decay channels open, while in the MSSM the $`Br[A\tau \tau (b\overline{b})]`$ in the high $`\mathrm{tan}\beta `$ regime hardly changes.
In the right-hand frame of the same figure, we summarize the Tevatronโs ability to explore the MSSM parameter space (in terms of both a $`2\sigma `$ exclusion curve and a $`5\sigma `$ discovery curve) using the process $`gg+b\overline{b}h+A+H\tau ^+\tau ^{}`$. Translating the enhancement factors above into this reach plot draws on the results of . As the $`M_A`$ mass increases up to about 140 GeV, the opening of the $`W^+W^{}`$ decay channel drives the $`\tau ^+\tau ^{}`$ branching fraction down, and increases the $`\mathrm{tan}\beta `$ value required to make Higgses visible in the $`\tau ^+\tau ^{}`$ channel. At still larger $`M_A`$, a very steep drop in the gluon luminosity (and the related $`b`$-quark luminosity) at large $`x`$ reduces the phase space for $``$ production. Therefore for $`M_A>`$170 GeV, Higgs bosons would only be visible at very high values of $`\mathrm{tan}\beta `$.
Figure 10 presents a qualitatively similar picture for LHC, based on the studies of $`h_{SM}\tau ^+\tau ^{}`$ of . The main differences compared to the Tevatron are that the required value of $`\mathrm{tan}\beta `$ at the LHC is lower for a given $`M_A`$ and it does not climb steeply for $`M_A>`$170 GeV because there is much less phase space suppression.
It is important to notice that both, Tevatron and LHC, could observe MSSM Higgs bosons in the $`\tau ^+\tau ^{}`$ channel even for moderate values of $`\mathrm{tan}\beta `$ for $`M_A200`$ GeV, because of significant enhancement of this channel. However the $`\gamma \gamma `$ channel is so suppressed that even the LHC will not be able to observe it in any point of the $`M_A<200`$ GeV parameter space studied in this paper! <sup>9</sup><sup>9</sup>9 In the decoupling limit with large values of $`M_A`$ and low values of $`\mathrm{tan}\beta `$, the lightest MSSM Higgs could be dicovered in the $`\gamma \gamma `$ mode just like the SM model Higgs boson, see e.g. ref.
### 4.2 Visibility of Technipions
In Section 3.2 we found a distinct enhancement of the $`P`$ signal in both the $`\tau ^+\tau ^{}`$ and $`\gamma \gamma `$ search channels for each of the technicolor models studied. As illustrated in the left frame of Figure 11, the available enhancement is well above what is required to render the $`P`$ of any of these models visible in the $`\tau ^+\tau ^{}`$ channel at the Tevatron. Likewise, the right frame of that figure shows that in the $`\gamma \gamma `$ channel at the Tevatron the technipions of models 3 and 4 will be observable at the $`5\sigma `$ level while model 2 is subject to exclusion at the $`2\sigma `$ level. The situation at the LHC is even more promising: Figure 12 shows that all four models could be observable at the $`5\sigma `$ level in both the $`\tau ^+\tau ^{}`$ (left frame) and $`\gamma \gamma `$ (right frame) channels.
### 4.3 Distinguishing the MSSM from Technicolor
In the previous section we have shown that that the Tevatron and LHC have the potential to observe the light (pseudo) scalar states characteristic of both supersymmetry and models of dynamical symmetry breaking. For both classes of models, the $`\tau ^+\tau ^{}`$ channel is enhanced and could be used for discovery of the light Higgs-like states.
Once a supposed light โHiggs bosonโ is observed in a collider experiment, an immediate important task will be to identify the new state more precisely, i.e. to discern โthe meaning of Higgsโ in this context. Comparison of the enhancement factors for different channels will aid in this task. Our study has shown that comparison of the $`\tau ^+\tau ^{}`$ and $`\gamma \gamma `$ channels can be particularly informative in distinguishing supersymmetric from dynamical models. In the case of supersymmetry, when the $`\tau ^+\tau ^{}`$ channel is enhanced, the $`\gamma \gamma `$ channel is suppressed, and this suppression is strong enough that even the LHC would not observe the $`\gamma \gamma `$ signature. In contrast, for the dynamical symmetry breaking models studied we expect simultaneous enhancement of both the $`\tau ^+\tau ^{}`$ and $`\gamma \gamma `$ channels. The enhancement of the $`\gamma \gamma `$ channel is so significant, that even at the Tevatron we may observe technipions via this signature at the $`5\sigma `$ level for Models 3 and 4, while Model 2 could be excluded at 95% CL at the Tevatron. The LHC collider, which will have better sensitivity to the signatures under study, will be able to observe all four models of dynamical symmetry breaking studied here in the $`\gamma \gamma `$ channel, and can therefore distinguish more conclusively between the supersymmetric and dynamical models.
We also would like to stress an important difference between two class of models in their production mechanisms. In supersymmetry the $`b\overline{b}`$ fusion process is likely to be as important as the $`gg`$ fusion mechanism (see Figure 6) in contributing to the total production cross section. In technicolor models, however, the $`b\overline{b}`$ fusion contribution to technipion production is likely to be negligible. This difference could be revealed, in principle, by looking at other (exclusive or semi-exclusive) processes: in case of supersymmetry, for example, one would expect significant enhancement of Higgs boson production associated with $`b`$-quarks.
## 5 Conclusions
In this paper we have shown that searches for a light Standard Model Higgs boson at Tevatron Run II and CERN LHC have the power to provide significant information about important classes of physics beyond the Standard Model. We demonstrated that the new scalar and pseudo-scalar states predicted in both supersymmetric and dynamical models can have enhanced visibility in standard $`\tau ^+\tau ^{}`$ and $`\gamma \gamma `$ search channels, making them potentially discoverable at both the Tevatron Run II and the CERN LHC. The enhancement arises largely from increases in the production rate; we showed that the model parameters exerting the largest influence on the enhancement size are $`\mathrm{tan}\beta `$ in the case of the MSSM and $`N_{TC}`$ and $`F_P`$ in the case of dynamical symmetry breaking. At the same time, the $`W^+W^{}`$ decay pathway is suppressed in the models studied here by at least an order of magnitude, compared to Standard Model expectations. In comparing the key signals for the non-standard scalars across models, we were able to show how one could start to identify which state has actually been found by a standard Higgs search. In particular, we investigated the likely mass reach of the Higgs search in $`pp/p\overline{p}\tau ^+\tau ^{}`$ for each kind of non-standard scalar state, and we demonstrated that $`ppp\overline{p}\gamma \gamma `$ may cleanly distinguish the scalars of supersymmetric models from those of dynamical models.
## 6 Acknowledgments
We thank M. Spira and R. Harlander for discussions, C.-P. Yuan for providing the code from Ref. , and a very thorough referee for a close reading of the manuscript. This work was supported in part by the U.S. National Science Foundation under awards PHY-0354838 (A. Belyaev) and PHY-0354226 (R. S. Chivukula and E. H. Simmons). A. Blum is supported in part by a scholarship from the German National Academic Foundation (Studienstiftung des deutschen Volkes). |
warning/0506/hep-th0506262.html | ar5iv | text | # Defect Structures in Lorentz and CPT Violating Scenarios
## I Introduction
The possibility of breaking Lorentz and CPT symmetries has been considered in several different contexts; see, e.g., Refs. CFJ ; CKL ; cordas . In CFJ the authors modify the usual Maxwell dynamics with the inclusion of a Chern-Simons-like term that violates both Lorentz and CPT symmetries. Other investigations with the addition of contributions that violate Lorentz and CPT symmetries have been done both at low energies, in the standard model CKL , and at higher energies, in string-like models cordas . For models dealing with CPT and Lorentz violating extensions of the standard model, sometimes one modifies the scalar Higgs sector, and this gives room for defect structures of more general profile, which may play important role to describe phase transitions in the earlier universe, due to spontaneous symmetry breaking.
Defects like domain walls, cosmic strings, monopoles and others have been studied in several different aspects Inicio , with applications to Cosmology Vil and Condensed Matter cmat . In particular, kinks are topological defects which in general connect distinct isolated minima in models that develop spontaneous breaking of some discrete symmetry. They appear in two-dimensional space-time, and can be embedded in the four-dimensional space-time, to generate bidimensional structures named domain walls. The role of such defects as seeds for the formation of non-topological structures is interesting lee and has led to several investigations, with the change of the discrete symmetry to an approximate symmetry as , and also when the symmetry is biased to make domains of distinct but degenerate vacua spring unequally b . In two-field models, topological defects may generate other interesting structures, such as defect inside defect did , and junctions of defects junction , and may be of interest in applications concerning conformational structure of polymers and polymer-like chains polymer . They may also induce interesting effects on other fields; for instance, the behavior of fermions in the background of kink-like structures is known to have very significant results jr , and could perhaps be re-examined within the Lorentz-violating scenario.
In this work we study models which combine the two issues, that is, we investigate kink-like structures in scenarios where Lorentz and CPT symmetries may be broken. Our main motivation is related to braneworld, specifically to the Randall-Sundrum scenario brane , because we may follow the lines of bfg and use the Lorentz-violating model described in Sec. III in the context of warped geometry with a single extra dimension. Another motivation is to bring some very well-known results for defect structures in models described by real scalar fields to this new scenario, where Lorentz and CPT symmetries do not play the standard role. In a recent work ludo , kinks were investigated in a model which breaks Lorentz symmetry with the explicit inclusion into the Lagrange density of Lorentz non-invariant higher-order derivative contribution. Our route is different, since we will study Lorentz and CPT breaking without introducing higher derivatives.
To do this, we follow Ref. CKL , in which the new terms arise as modifications in the Higgs sector of the standard model. In the light of the recent understanding of equivalence between non-commutative field theory and Lorentz-violanting extensions involving ordinary fields chkl , the present work is also of interest to non-commutative solitons ns , which has been investigated for a variety of reasons, including self-consistent deformation of the highly constrained structure of local quantum field theory, and the breaking of locality at short distances, which is of direct interest to quantum gravity. Also, the appearing of non-commutativity in field theory in a limit of string theory string provides fresh interest to the subject, in particular on D-branes, specially as non-commutative solitons of tachyon fields of open string theory t .
Our investigations consider static solutions in one spatial dimension. Thus, the static solutions that we consider cannot see effects of non-commutativity. However, we can use the point of view of Ref. vy to investigate how stability modifies the bound states of the model for non-commutative space-time. Moreover, our investigations is also of interest to the non-commutative aspects introduced in Ref. kkl , which investigates kinks and domain walls for non-commutative field theory, directly connected to the tachyon action for unstable brane in open strings; see the recent revision of Ref. s for a variety of motivations on tachyon dynamics.
We organize our work as follows: in the next Sect. II we consider models described by one and by two real scalar fields. There we realize that two-field models lead to richer possibilities, and we show how to extend the Bogomolโnyi bound to the Lorentz and CPT breaking scenario. In Sect. III we investigate an explicit model of two real scalar fields, which can be seen as an extension of a former model, first investigated in Bazeia , which has been used in several other contexts, for instance in Refs. did ; junction ; polymer ; Bazeia1 , engendering broader interest. As we will show, the breaking of both Lorentz and CPT symmetries gives rise to an asymmetry between defects and anti-defects, including the presence of linearly stable solutions that support regions of negative energy density.
We end this work in Sect. IV, where we include our comments and conclusions, pointing some possible extensions of this work.
## II Scalar field models
In this work we investigate defect structures described by scalar fields in models which break Lorentz and CPT symmetries explicitly.
We start with the simplest case, which describes a single real scalar field. In this case, we study models where only the Lorentz symmetry is broken. Next, we deal with two scalar fields, and there we investigate models which break both Lorentz and CPT symmetries.
### II.1 One field
We start with the model
$$=\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{1}{2}\kappa ^{\mu \nu }_\mu \varphi _\nu \varphi V,$$
(1)
where $`V=V(\varphi )`$ is the potential, which controls the way the field self-interacts. We are working in $`(1,1)`$ spacetime dimensions, and the metric is diag$`(g_{\mu \nu })=(1,1),`$ with $`\kappa ^{\mu \nu }`$ being a constant tensor, given by
$$\kappa ^{\mu \nu }=\left(\begin{array}{cc}\beta & \alpha \\ \alpha & \beta \end{array}\right),$$
(2)
where $`\alpha `$ and $`\beta `$ are real parameters. See Ref. KL for other details. For simplicity, we take $`\beta =0`$ for the explicit calculations that follow.
The equation of motion is
$$\ddot{\varphi }\varphi ^{\prime \prime }+2\alpha \dot{\varphi }^{}+\frac{dV}{d\varphi }=0,$$
(3)
where $`\dot{\varphi }=\varphi /t`$ and $`\varphi ^{}=\varphi /x,`$ etc. For static field we get
$$\varphi ^{\prime \prime }=\frac{dV}{d\varphi }.$$
(4)
This is the same equation one gets in the standard situation. Thus, static solutions violate neither Lorentz nor CPT symmetries. However, for time-dependent field, we search for traveling waves and now the equation of motion may have solutions which violate Lorentz and CPT symmetries.
Although our model violates Lorentz symmetry, we can still search for traveling waves in the form $`\varphi =\varphi (u),`$ where $`u=\gamma (xvt),`$ but now $`\gamma =\gamma (v,\alpha )`$ may not have the usual form. We use this into eq. (3) to get to
$$\frac{d^2\varphi }{du^2}=\frac{dV}{d\varphi },$$
(5)
if one sets $`\gamma =1/\sqrt{1v^2+2\alpha v}.`$ This is a general result: it shows that for any static field $`\varphi _s(x)`$ \[topological (kinklike) or nontopological (lumplike)\] which solves eq. (4), there is a traveling wave of the form
$$\varphi (u)=\varphi _s(u),$$
(6)
which solve eq. (5). The traveling wave has the form of a static solution, and it travels with constant velocity $`v,`$ with width $`w=w_0/\gamma ,`$ for $`w_0`$ being the width of the static solution. The velocity is restricted to the interval $`v(\sqrt{1+\alpha ^2}+\alpha ,\sqrt{1+\alpha ^2}+\alpha ).`$ We notice that the limit $`\alpha 0`$ leads to the standard situation, with $`\gamma =\gamma (v,0)=1/\sqrt{1v^2},`$ and $`v(1,1).`$ We also notice that for $`\alpha `$ very small we get $`v(1+\alpha ,1+\alpha ),`$ which shifts by $`\alpha `$ the standard velocity interval.
We consider the model (1) in the absence of potential; this case was recently considered in Ref. amw , with other motivations. The massless excitations now give $`w^2k^22\alpha wk=0,`$ which implies that the velocities should obey $`v_\pm =\pm \sqrt{1+\alpha ^2}+\alpha .`$ They travel with different velocities in the forward and backward directions, showing that the model engenders birefringence. The inclusion of the potential will make the excitations massive, with velocity bounded by the two massless values. This gives an alternative way to understand the bounds in the velocity of the traveling waves that we have just obtained.
The parameter $`\alpha `$ induces an asymmetry for traveling waves with positive and negative velocities, breaking Lorentz invariance. We also see that the time-dependent solutions violate both parity and time reversal, although they are symmetric under PT. Thus, they do not violate CPT, because the scalar field is even under charge conjugation.
We calculate $`\theta ^{\mu \nu }`$ to get the four entries: they are densities which represent energy $`\theta ^{00}`$, energy flux $`\theta ^{10},`$ momentum $`\theta ^{01},`$ and pressure $`\theta ^{11}`$. They are given by
$`\theta ^{00}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{1}{2}}\varphi ^2+V,`$ (7a)
$`\theta ^{01}`$ $`=`$ $`\dot{\varphi }\varphi ^{}\alpha \varphi ^2,`$ (7b)
$`\theta ^{10}`$ $`=`$ $`\dot{\varphi }\varphi ^{}+\alpha \dot{\varphi }^2,`$ (7c)
$`\theta ^{11}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{1}{2}}\varphi ^2V.`$ (7d)
We notice that both equations $`_0\theta ^{00}+_1\theta ^{10}=0`$ and $`_0\theta ^{01}+_1\theta ^{11}=0`$ work on shell. The fact that $`\theta ^{01}\theta ^{10}`$ indicates violation of Lorentz symmetry.
The energy for traveling waves can be written in the form
$$\frac{E_v}{E_0}=\gamma (1+\alpha v),$$
(8)
where $`E_0`$ stands for the energy of the static solutions. We calculate the energy ratio for solutions with opposite velocities to get
$$\frac{E_v}{E_v}=\frac{1+\alpha v}{1\alpha v}\sqrt{\frac{1v^22\alpha v}{1v^2+2\alpha v}},$$
(9)
which is asymmetric, thus violating Lorentz symmetry. We notice that $`E_v<E_v`$ for $`\alpha v>0.`$
This is the general scenario for kinks and lumps in models of the form given by eq. (1). The traveling waves are even under CPT, but they violate Lorentz symmetry.
We illustrate this case with the $`\varphi ^4`$ model. It is described by the potential $`V(\varphi )=(1/2)(1\varphi ^2)^2,`$ where we are using dimensionless field and coordinates. The static kink has the form $`\varphi _s(x)=\mathrm{tanh}x.`$ It has unit width, and we have chosen $`x=0`$ as the center of the solution. The corresponding traveling wave is given by $`\varphi (u)=\mathrm{tanh}\gamma (xvt),`$ which has width $`1/\gamma .`$
We can widen the above investigations using some recent results on deformed defects dd . For the model (1), if one modifies the potential according to
$$V(\varphi )U(\phi )=V(\varphi f(\phi ))/f^2(\phi ),$$
(10)
where $`f=f(\phi )`$ is the deformation function, we can obtain static solution for the modified model in terms of static solution of the starting model. That is, if $`\varphi _s(x)`$ is solution for the potential $`V(\varphi ),`$ then
$$\phi _s(x)=f^1(\varphi _s(x)),$$
(11)
is solution for the modified model with potential $`U(\phi ).`$ Evidently, the presence of the Lorentz breaking term in the model (1) does not modify this result, which shows that the deformation prescriptions introduced in Refs. dd are very naturally extended to traveling waves in the above Lorentz violating scenario.
Before going deeper into Lorentz-violating investigations, some words of caution seem to be necessary. It is important to notice that for the model (1) with $`\kappa ^{\mu \nu }`$ given by (2), we can redefine field and coordinates in order to eliminate Lorentz violation km . This shows that this model is fake Lorentz-violating theory, but we have decided to make the above investigations because it illustrates with simple terms how Lorentz-violating ingredients enter the game for kinks and lumps in $`(1,1)`$ space-time dimensions. Evidently, the procedure suggested to eliminate Lorentz violation indicates that we can extend the energy-momentum tensor (7) in order to make it symmetric and conserved, thus eliminating the presence of Lorentz violation. However, this procedure to eliminate Lorentz violation may not work when we couple the model with more sophisticated fields.
Another issue concerns the need to make the classical solutions time-dependent to make them feel the presence of Lorentz violation. This fact reminds us very much of the investigations done in Ref. vy , in which non-commutativity is only seen by the fluctuations around classical static kinks in $`(1,1)`$ non-commutative space-time. This point will be further explored in a forthcoming investigation, in which we deal with stability of the Lorentz-violating solutions that appear in this work.
### II.2 Two fields
We now turn attention to two-field models. Firstly, we consider the class of models
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi +{\displaystyle \frac{1}{2}}\kappa ^{\mu \nu }_\mu \varphi _\nu \varphi +`$ (12)
$`{\displaystyle \frac{1}{2}}_\mu \chi ^\mu \chi +{\displaystyle \frac{1}{2}}\kappa ^{\mu \nu }_\mu \chi _\nu \chi V(\varphi ,\chi ).`$
This class of models can be seen as an extension for two fields of the class introduced in the former Sect. II.A. Thus, it also suffers from the same problem of being fake Lorentz-violating theory km , but we explore some peculiarities before introducing a genuine Lorentz-violation family of models. Our point is that these models may be seen as effective portions of some more sophisticated models, involving coupling with other more complex fields.
The equations of motion are given by
$`\ddot{\varphi }\varphi ^{\prime \prime }+2\alpha \dot{\varphi }^{}+{\displaystyle \frac{V}{\varphi }}=0,`$ (13a)
$`\ddot{\chi }\chi ^{\prime \prime }+2\alpha \dot{\chi }^{}+{\displaystyle \frac{V}{\chi }}=0.`$ (13b)
Thus, for static solutions we get
$`\varphi ^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{V}{\varphi }},`$ (14a)
$`\chi ^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{V}{\chi }},`$ (14b)
which do not depend on $`\alpha ,`$ and so they correspond to standard models. The case
$$V(\varphi ,\chi )=\frac{1}{2}W_\varphi ^2+\frac{1}{2}W_\chi ^2,$$
(15)
where $`W_\varphi =W/\varphi `$ and $`W_\chi =W/\chi ,`$ leads to models of the form considered in Refs. Bazeia ; Bazeia1 and in other works.
We consider traveling waves in the form $`\varphi =\varphi (u)`$ and $`\chi =\chi (u)`$ with $`u=\gamma (xvt),`$ as before. The equations of motion change to
$`{\displaystyle \frac{d^2\varphi }{du^2}}`$ $`=`$ $`{\displaystyle \frac{V}{\varphi }},`$ (16a)
$`{\displaystyle \frac{d^2\chi }{du^2}}`$ $`=`$ $`{\displaystyle \frac{V}{\chi }},`$ (16b)
where we have set $`\gamma =1/\sqrt{1v^2+2\alpha v}.`$ For this reason, if the model supports static solutions $`\varphi _s(x)`$ and $`\chi _s(x),`$ it also supports traveling waves in the form
$$\varphi =\varphi _s(u),\chi =\chi _s(u),$$
(17)
which travels with constant velocity $`v,`$ and with width $`w=w_0/\gamma ,`$ as before.
This class of models is similar to the former one, and it may support traveling waves which preserve CPT, although they violate Lorentz symmetry. We notice that extensions to a set of $`N`$ real scalar fields works straightforwardly.
Another class of models can be considered. In this case the Lagrange density has the form
$$=\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{1}{2}_\mu \chi ^\mu \chi +\kappa ^\mu \varphi _\mu \chi V(\varphi ,\chi ).$$
(18)
The presence of the vector $`\kappa ^\mu =(a,b),`$ $`a`$ and $`b`$ being real parameters, leads to both Lorentz and CPT violation; see Ref. CKL for other details. The model may support kinks and lumps, if the potential $`V=V(\varphi ,\chi )`$ is chosen properly. This class of models may support defect structures which violate both Lorentz and CPT symmetries, leading to richer scenarios. In particular, we are now dealing with a genuine Lorentz-violating family of models, since it is not possible to remove the Lorentz-violating $`\kappa `$-dependent term from the theory anymore; see CKL and, in particular K , in connection with a varying coupling.
For the model at hand, the equations of motion have the form
$`_\mu ^\mu \varphi \kappa ^\mu _\mu \chi +{\displaystyle \frac{V}{\varphi }}=0,`$ (19a)
$`_\mu ^\mu \chi +\kappa ^\mu _\mu \varphi +{\displaystyle \frac{V}{\chi }}=0.`$ (19b)
The energy-momentum tensor has the four entries:
$`\theta ^{00}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\dot{\varphi }^2+\dot{\chi }^2+\varphi ^2+\chi ^2)b\varphi \chi ^{}+V,`$ (20a)
$`\theta ^{10}`$ $`=`$ $`\varphi ^{}\dot{\varphi }\chi ^{}\dot{\chi }+b\varphi \dot{\chi },`$ (20b)
$`\theta ^{01}`$ $`=`$ $`\varphi ^{}\dot{\varphi }\chi ^{}\dot{\chi }a\varphi \chi ^{},`$ (20c)
$`\theta ^{11}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\dot{\varphi }^2+\dot{\chi }^2+\varphi ^2+\chi ^2)+a\varphi \dot{\chi }V.`$ (20d)
We notice that the equations $`_\mu \theta ^{\mu \nu }=0`$ work on shell. Also, $`\theta ^{01}\theta ^{10}`$ shows that the model volates Lorentz symmetry. In this case, it is not possible to improve the energy-momentum tensor to make it symmetric and conserved; this is a true manifestation of Lorentz violation for this family of models CKL ; K .
For static fields, that is, for field configurations that only depend on the space coordinate $`x,`$ the equations of motion become
$`\varphi ^{\prime \prime }+b\chi ^{}`$ $`=`$ $`{\displaystyle \frac{V}{\varphi }},`$ (21a)
$`\chi ^{\prime \prime }b\varphi ^{}`$ $`=`$ $`{\displaystyle \frac{V}{\chi }}.`$ (21b)
These equations do not depend on $`a;`$ thus, if one chooses $`b`$ equal to zero, the static solutions are not affected by Lorentz and CPT symmetries. However, they may be affected by the motion of traveling waves, as we have already shown in the former case.
For nonzero $`b,`$ we see that the above equations violate both Lorentz and CPT symmetries. They do not respect parity transformation, although they are even under T and C. The absence of parity symmetry breaks the kink $``$ antikink exchange scenario, which in general appears in models that do not violate parity. However, we notice that the substitutions $`xx`$ and $`bb`$ do not change the equations of motion (21) for static fields, if the potential is even under $`bb.`$ In this case, kinks for the model with $`b`$ positive would become antikinks for the model with $`b`$ negative.
The presence of $`b`$ in the equations of motion and energy density changes the standard scenario. To attain a Bogomolโnyi bound BPS we modify the potential in Eq. (15). We consider a new class of models, identified by
$$V_s(\varphi ,\chi )=\frac{1}{2}(W_\varphi +s_1\chi )^2+\frac{1}{2}(W_\chi +s_2\varphi )^2,$$
(22)
where $`W=W(\varphi ,\chi )`$ is a smooth function of the two fields, with $`s_1`$ and $`s_2`$ being real constants, which obey $`s_2s_1=b`$. This potential is an extension of the potential considered in Bazeia ; it gets to its original form in the limit $`b0.`$ This modification is introduced to attain a Bogomolโnyi bound BPS , but it changes the way the fields interact, since the potential now depends explicitly on $`b,`$ the parameter which breaks Lorentz and CPT symmetries.
This class of models can be further investigated for the presence of topological solutions. We consider static fields, $`\varphi =\varphi (x)`$ and $`\chi =\chi (x)`$. We write the energy density for static solutions in the form
$$\theta ^{00}=\frac{dW}{dx}+\frac{1}{2}\left(\varphi ^{}W_\varphi s_1\chi \right)^2+\frac{1}{2}\left(\chi ^{}W_\chi s_2\varphi \right)^2.$$
(23)
The energy is minimized to the value $`E^{ij}=\mathrm{\Delta }W_{ij},`$ with $`\mathrm{\Delta }W_{ij}=W_iW_j,`$ for $`W_i=W(\overline{\varphi }_i,\overline{\chi }_i),`$ and $`v_i=(\overline{\varphi }_i,\overline{\chi }_i)`$ being a minimum of the potential, obeying $`V(\overline{\varphi }_i,\overline{\chi }_i)=0.`$ This bound is attained for field configurations which obey the first-order equations
$`\varphi ^{}`$ $`=`$ $`W_\varphi +s_1\chi ,`$ (24a)
$`\chi ^{}`$ $`=`$ $`W_\chi +s_2\varphi ,`$ (24b)
with the boundary conditions: the pair $`(\varphi ,\chi )`$ goes to $`(\overline{\varphi }_i,\overline{\chi }_i)`$ for $`x\mathrm{},`$ and to $`(\overline{\varphi }_j,\overline{\chi }_j)`$ for $`x\mathrm{}.`$ This is the Bogomolโnyi bound BPS , now extended to the above class of models, which violate Lorentz and CPT symmetries.
We can see that solutions of the above first-order equations solve the equations of motion. Also, despite the modification in the model, the static solutions satisfy
$$\frac{1}{2}\varphi ^2+\frac{1}{2}\chi ^2=V,$$
(25)
which shows that the gradient and potential portions of the energy contribute equally.
We remark that since $`b`$ changes the energy density in eq. (20a), the form (23) is only obtained when we consider the potential in the specific form (22), with $`s_2s_1=b`$. We compare this with the case which preserves Lorentz and CPT symmetries to see that the Bogomolโnyi bound requires the inclusions of new terms into the potential.
The asymmetry that appears for $`b0`$ may contribute to destabilize the defect solutions. However, we can show that solutions to the above first-order equations are linearly stable. The calculation follows the standard route BS . The full investigation will be done in another work, and here we show the main steps of the calculation. This investigation is important, because we will show below that there are models which support kinks of unusual profile. We introduce general fluctuations for the two fields in the form: $`\varphi (x,t)=\varphi (x)+\eta (x,t)`$ and $`\chi (x,t)=\chi (x)+\xi (x,t).`$ We use these fields in the equations of motion to get to the Schrรถdinger-like equation, $`H\mathrm{\Psi }_n(x)=\omega _n^2\mathrm{\Psi }_n(x),`$ where $`\mathrm{\Psi }_n(x)`$ is a two-component wave function and the Hamiltonian has the form
$$H=\frac{d^2}{dx^2}ib\sigma _2\frac{d}{dx}+U,$$
(26)
where $`\sigma _2`$ is a Pauli matrix and
$`U=\left(\begin{array}{ccc}^2V_s/\varphi ^2& ^2V_s/\varphi \chi & \\ ^2V_s/\chi \varphi & ^2V_s/\chi ^2& \end{array}\right).`$ (29)
We use $`V_s`$ as in Eq. (22) to write $`H=S^{}S,`$ where $`S`$ is the first-order operator
$$S=\frac{d}{dx}+\left(\begin{array}{ccc}W_{\varphi \varphi }\hfill & W_{\varphi \chi }+s_1\hfill & \\ W_{\chi \varphi }+s_2\hfill & W_{\chi \chi }\hfill & \end{array}\right).$$
(30)
This shows that $`H`$ is non-negative, and so the corresponding eigenvalues must obey $`w_n^20.`$ This result is general; it extends the result of Ref. BS to the above model, and it shows that the solutions of the first-order equations (24) are linearly stable.
## III Example
The last class of models deserves further attention. We illustrate this case with an example. We consider $`s_1=0`$ and $`s_2=b,`$ and the following function Bazeia
$$W(\varphi ,\chi )=\varphi \frac{1}{3}\varphi ^3r\varphi \chi ^2,$$
(31)
where $`r`$ is a real parameter. This gives the potential
$$V(\varphi ,\chi )=\frac{1}{2}(1\varphi ^2r\chi ^2)^2+\frac{1}{2}\left(2r\varphi \chi b\varphi \right)^2.$$
(32)
The model may support several minima, depending on the values of $`r`$ and $`b.`$ We consider $`r`$ and $`b`$ positive, and $`b^2/4r(0,1)`$ to write
$$v_{h\pm }=(\pm Q,b/2r),v_{v\pm }=(0,\pm \sqrt{1/r}),$$
(33)
where $`Q=\sqrt{1b^2/4r}.`$ There are four minima, two horizontally aligned, and two vertically aligned, as the subscripts indicate. The limit $`b0`$ implies $`Q1,`$ bringing the minima $`v_{h\pm }`$ back to $`(\pm 1,0),`$ to the $`\varphi `$ axis, as expected Bazeia .
There are five topological sectors, for solutions that solve the first-order equations, one with energy or tension $`t_1=(4/3)Q^3,`$ and four with tensions degenerate to the value $`t_2=(2/3)Q^3.`$ As one knows, in the absence of Lorentz and CPT violation, the standard situation engenders BPS and anti-BPS solutions, which connect the minima in the two possible senses. However, parity violation breaks this symmetry, excluding one of the two possibilities. In the model under investigation, for instance, in the more energetic sector, there is only one solution, connecting $`v_hv_{h+}.`$ The same for the other sectors, where there are solutions connecting $`v_hv_{v+},`$ $`v_{v+}v_{h+},`$ $`v_hv_v,`$ and $`v_vv_{h+}.`$ In Fig. we illustrate how the orbits appear connecting the minima of the potential.
The model may admit another sector, connecting the minima $`v_{v\pm }.`$ This sector cannot have solutions that obey the first-order equations. Although in this case we have been unable to find any explicit solution connecting the two minima asymptotically, we could verify that the straight line orbit which solves the model for $`b=0`$ does not exist in the present case, for $`b0.`$
It is interesting to notice that in a string theory scenario with the above realization of Lorentz and CPT violations, the asymmetry between defects and anti-defects prevent the presence of anti-defects. If this persists in the string theory, it would certainly prevent the presence of open strings ending on a pair braneโanti-brane, and this would certainly change the way tachyon condensation could appear.
For the model under investigation, the first-order equations are
$`\varphi ^{}`$ $`=`$ $`1\varphi ^2r\chi ^2,`$ (34a)
$`\chi ^{}`$ $`=`$ $`b\varphi 2r\chi \varphi .`$ (34b)
It is not hard to see that these equations admit the integrating factor $`f(\chi )=1/(\chi b/2r)^{1+1/r}.`$ Thus, we use $`\stackrel{~}{\chi }=\chi b/2r`$ to write the orbits, for $`r1/2`$ and $`r1,`$
$$\varphi ^2=\frac{r}{2r1}\stackrel{~}{\chi }^2+\frac{b}{r1}\stackrel{~}{\chi }+C\stackrel{~}{\chi }^{\frac{1}{r}}+Q^2,$$
(35)
where $`C`$ is an integration constant. The limit $`b0`$ changes this result to the orbits first obtained in Ref. Izquierdo:2002sz .
The specific cases $`r=1`$ and $`r=1/2`$ need particular attention. They have orbits given by, respectively,
$`\varphi ^2`$ $`=`$ $`\stackrel{~}{\chi }^2+{\displaystyle \frac{b}{r}}\stackrel{~}{\chi }\mathrm{ln}\stackrel{~}{\chi }+C\stackrel{~}{\chi }+Q^2,`$ (36a)
$`\varphi ^2`$ $`=`$ $`C\stackrel{~}{\chi }^2+\stackrel{~}{\chi }^2\mathrm{ln}\stackrel{~}{\chi }{\displaystyle \frac{b}{r}}\stackrel{~}{\chi }+Q^2.`$ (36b)
We have being unable to solve the first-order equations analytically for $`r`$ and $`C`$ arbitrary. For this reason, we have used some specific values for $`C`$: firstly, we take the limit $`C\mathrm{},`$ to see that in this case the orbit is a straight line segment joining $`v_{h+}`$ and $`v_h`$ with $`\chi =b/2r.`$ This limit reduces the first-order equations (34) to the single equation $`\varphi ^{}=Q^2\varphi ^2,`$ which is solved by
$$\varphi (x)=Q\mathrm{tanh}(Qx),$$
(37)
where we are using $`x=0`$ as the center of the solution. The corresponding energy density is given by $`ฯต=Q^4\mathrm{sech}^4(Qx)`$.
Another interesting value for the integration constant is $`C=0`$. This choice leads to the solutions
$`\varphi _\pm (x)`$ $`=`$ $`{\displaystyle \frac{Q\mathrm{sinh}(2rQx)}{\pm B+\mathrm{cosh}(2rQx)}},`$ (38a)
$`\chi _\pm (x)`$ $`=`$ $`{\displaystyle \frac{b}{2r}}\pm {\displaystyle \frac{A}{\pm B+\mathrm{cosh}(2rQx)}},`$ (38b)
where we have used $`A=(1r)Q^2K`$ and $`B=bK/4,`$ where
$$K=\sqrt{\frac{12r}{r(12r+r^2Q^2)}},$$
(39)
with $`r(0,1/2).`$
We notice that the limit $`b0`$ changes the solutions (38) to the simpler form
$`\varphi ^0(x)`$ $`=`$ $`\mathrm{tanh}(2rx),`$ (40a)
$`\chi _\pm ^0(x)`$ $`=`$ $`\pm \sqrt{\left({\displaystyle \frac{1}{r}}2\right)}\mathrm{sech}(2rx),`$ (40b)
which are solutions of the model first investigated in Ref. Bazeia . We recall that the above solutions were found with the elliptic orbits
$$\varphi ^2+\frac{r}{12r}\chi ^2=1,$$
(41)
which are good orbits for $`r(0,1/2).`$ We notice that these orbits are exactly the orbits obtained in eq. (35) in the limit $`b0`$ for the value $`C=0.`$
The energy density corresponding to the above solutions can be written as
$$\theta ^{00}=\varphi ^2+\chi ^2b\varphi \chi ^{},$$
(42)
and for the non-trivial solutions with $`C=0`$ we use Eqs. (38) and (39) to obtain
$`\theta _\pm ^{00}(x)`$ $`=`$ $`{\displaystyle \frac{4r^2Q^4}{[B\pm \mathrm{cosh}(2rQx)]^4}}[1+B^2\pm `$ (43)
$`\mathrm{cosh}(2rQx)\left(2B+{\displaystyle \frac{bA}{2rQ^2}}\mathrm{sinh}^2(2rQx)\right)+`$
$`({\displaystyle \frac{A^2}{Q^2}}+{\displaystyle \frac{bAB}{2rQ^2}}+B^2)\mathrm{sinh}^2(2rQx)].`$
The orbits and solutions for $`C=0`$ are shown in Fig. and , respectively, and in Fig. we plot the corresponding energy densities. These figures are shown for $`r=1/4`$ and $`b=1/3.`$ We see that the upper orbit gives standard defect structures. However, the lower orbit gives unusual defects, making the topological solution non monotonic, a fact due to the breaking of Lorentz invariance, which also responds for the presence of regions of negative energy density, as shown in Fig. .
To introduce specific results, we notice that in the defect solution for lower orbit, the behavior of the $`\varphi `$ field, which ensures the topological profile of the solution, shows two critical points, at the values $`x_c^\pm =\pm (1/2rQ)\mathrm{arcsech}(\mathrm{B}),`$ for which $`\varphi (x_c^\pm )=\pm Q/\sqrt{1B^2}.`$ For these values, the energy density is given by
$$\theta _{}^{00}(x_c)=rb^2\frac{(12r)(12r+r^2Q^2)}{(1r)^4},$$
(44)
and it is always negative, for the range of values that we are considering. For $`r=1/4`$ and $`b=1/3`$ we get $`x_c^\pm =\pm 3.8575.`$ Although the energy is positive, the energy density is negative in the two regions $`|x|3.0625,`$ which include the critical points of $`\varphi ;`$ see Fig. . These regions of negative energy densities form the outer side of the defect, and they disappear in the limit $`b0,`$ in the absence of Lorentz and CPT breaking. The core of the defect changes insignificantly for $`b`$ small, and so it may entrap another defect in the same way it used to do in the standard situation did . The appearance of negative energy density is an unusual behavior, which leads us to think that such solutions are unstable, but we have already show that they are linearly stable in general. We will further investigate stability in another work, to examine how to find stable solutions for specific models which violate both Lorentz and CPT symmetries.
The value $`b=1/3`$ is not small. Since $`b`$ measures how the model deviates from the standard situation, it should be very small. Former studies on bounds in the Higgs sector for extensions of the standard model suggest the order of magnitude of $`b.`$ The constraint is very tight in more realistic situations bound . In our toy model, however, we have used $`b=1/3`$ to highlight the effects the breaking of Lorentz and CPT symmetries may induce in the defect structures that appear in the model under consideration. Moreover, the present investigations may be of some use in applications to condensed matter โ see, e.g., Refs.cmat ; polymer โ and there violation of Lorentz invariance should have another interpretation. Indeed, in condensed matter we have found interesting investigations m in which one deals with very similar solutions, engendering profiles of almost the same type of the kink-like solutions that appear for a non-vanishing $`b,`$ not that much small. We can also mimic Lorentz-violating models in condensed matter with materials which naturally select preferable directions in space, which can be described with continuum version of the Dzyaloshinkii-Moriya model DM .
## IV Comments and Conclusions
In this work we have investigated models described by real scalar fields, in scenarios which violate both the Lorentz and CPT symmetries. We first dealt with models described by a single real field, and there we have shown that the addition of the Lorentz breaking term changes no static sector of the model. However, traveling waves see the Lorentz breaking, and we have constructed the way the traveling waves appear. Moreover, we have extended this result to deformed defects, and to models described by two or more real scalar fields.
In the case of two fields, we have invented another class of models, and we have investigated an explicit example, which generalizes former results to the Lorentz and CPT breaking scenario. These models do not support the usual defect and anti-defect structures simultaneously, and there are solutions that engender unusual profile, making the energy density negative in the outer side of the defect. The asymmetry for defect and anti-defect that we have found may perhaps be of some use to build string theory scenarios where open strings ending on a braneโanti-brane system are suppressed by CPT violation.
The present investigations will continue in another work, where we study linear stability of the solutions that we have just found in this paper. There we will show explicitly how to construct stable defect structures which violate Lorentz and CPT symmetries. We will also investigate supersymmetric extensions rc ; bbc ; prl of the above models, to see how the solutions of the first-order equations behave as BPS states.
We believe that the idea that the geodesic motion in moduli space can be used to describe the low energy dynamics of defect structures manton may be extended to the present context. Eventually, it may change the scenario constructed in es for the standard model, which preserves both Lorentz and CPT symmetries.
The suggestion that the models here studied may mimic features of more realistic systems, can also be extended to the case of heterotic M theory, following the recent work ant , which has investigated the effects of collision of scalar field kinks with boundaries, motivated from its cousin, the five dimensional heterotic M theory. The investigation shows that kink-boundary effects appears as direct application of the moduli space evolution.
Other lines of investigations concern the presence of junctions of defects, in Lorentz and CPT violating scenarios. Work on this is now in progress, in models which follow the lines of Ref. junction . We are also exploring similar models, with focus on tachyon kinks, motivated by ideas present in Refs. t ; kkl ; s ; tk . Furthermore, the inclusion of fermions is important not only for supersymmetry, but also to allow investigations concerning the behavior of fermions jr in the background of these Lorentz-violating kink-like structures. Another issue concerns the use of defect structures in scalar field theory to generate brane in warped geometry with a single extra dimension, as motivated by Ref. brane . Practical possibilities have already been examined in Ref.bfg , and we are now searching for brane within the present Lorentz-violating scenario. Evidently, the presence of Lorentz violation requires that we somehow modify the standard scenario, with the addition of extra terms to compensate the asymmetry of the energy-momentum tensor. Similar recent investigation was done in jp , where a Chern-Simons modification of General Relativity has been considered, which may help us enlighten the issue.
The authors would like to thank F.A. Brito, A.R. Gomes, L. Losano, J.R. Nascimento and V.M. Pereira for discussions, and CAPES, CNPq, PADCT/CNPq, PROCAD/CAPES, and PRONEX/CNPq/FAPESQ for financial support. |
warning/0506/cond-mat0506255.html | ar5iv | text | # Rolling friction and bistability of rolling motion
## 1 INTRODUCTION
Rolling friction belongs to the most important dissipation mechanisms which may sensitively affect the dynamical behavior of granular systems. Since 1785 when Vince described systematic experiments to determine the nature of friction laws , rolling friction has been investigated by many scientists due to its great importance in engineering and natural sciences, e.g. \[Czichos (1978), Barwell (1979), Szeri (1980), Dowson et al. (1994), Shpenkov (1995), Rabinowicz (1965)\].
Experimentally it was observed that the rolling friction coefficient depends non-monotonically on the velocity \[Olson et al. (2003), Yi et al. (2002), Gustafsson (1997), Ray (1997), Bakker et al. (1987), Burckhardt (1993), Yi & Jeong (1998)\]: For small velocities the rolling friction force increases with velocity, while for fast motion it decays. In the case of a soft cylinder rolling on a hard plane \[Brilliantov & Pรถschel (1998)\], the contact surface between the bodies is flat which allows for the application of Hertzโ contact theory extended to the contact of viscoelastic particles \[Brilliantov et al. (1996)\]. In the opposite case, this assumption is not justified since the plane shape follows the shape of the rolling body in the contact area. Here we investigate rolling of a hard cylinder on a soft plane and bistability of rolling motion as a consequence of the resulting friction law.
## 2 MODEL
Consider a cylinder of radius $`R`$, mass $`M`$ and moment of inertia $`I`$ which rolls at velocity $`v`$ on an inclined plane (quantities such as $`M`$, $`I`$, as well as the forces $`F_R`$, $`F_{\mathrm{ex}}`$ and similar quantities, are given as per unit length of the cylinder, i.e. as line densities), see Fig. 1.
For certain materials it has been shown that surface effects such as adhesion may have significant influence on rolling friction (e.g. \[Barquins et al. (1978), Deryaguin & Toporov (1994), Kendall (1975), Fuller & Roberts (1981), Roberts & Thomas (1975)\]). On the other hand, for viscoelastic materials it was reported that rolling friction is due very little to surface interactions, i.e., the major part is due to deformation losses within the bulk of the material \[Tabor (1955), Tabor (1952), Atack & Tabor (1958), Drutowski (1959)\]. We will, therefore, neglect surface effects. The surface is modeled by non-interacting springs (for the justification of this approximation see \[Pรถschel et al. (1999)\]). The properties of the plane are described by the coefficients $`k`$, $`\gamma `$, and $`m`$: $`k\mathrm{d}\xi `$ and $`\gamma \mathrm{d}\xi `$ are the elastic and dissipative forces and $`m\mathrm{d}\xi `$ is the mass of the springs in the region $`\mathrm{d}\xi `$. For $`x_{}xx_+`$ the surface is deformed by the cylinder. Assuming Hookeโs law, the dynamics of the surface is described by
$$m\ddot{y}\left(x\right)+\gamma \dot{y}\left(x\right)+ky\left(x\right)=f(x,t),$$
(1)
where $`f(x,t)`$ is the density of the force which acts onto the springs in the region of contact. This force has to be provided by the cylinder. In a wide range of material parameters, the condition $`k<\gamma ^2/4k^2`$ assures that the contact area is continuous, i.e., the cylinder contacts the surface everywhere in the interval $`(x_{},x_+)`$, see \[Pรถschel et al. (1999)\] for details.
From geometry we find the deformation $`y\left(x\right)`$ in the contact area $`(x_{},x_+)`$:
$$y\left(x\right)=Rh\sqrt{R^2\left(xx_c\right)^2}\frac{\left(xx_c\right)^2}{2R}h,$$
(2)
where $`x_c(t)`$ is the position of the cylinder, $`hy_{\mathrm{min}}=y\left(x_c\right)`$ is the penetration depth of the cylinder and $`R\left|xx_c\right|`$ was assumed. Thus, in the stationary state, $`x_c=vt`$, we obtain
$$\dot{y}\left(x\right)=v\frac{xx_c}{R},\ddot{y}(x)=\frac{v^2}{R}$$
(3)
and, therefore,
$$f\left(\xi \right)=\frac{k}{2R}\xi ^2\frac{\gamma v}{R}\xi +\frac{mv^2}{R}hk.$$
(4)
with $`\xi xx_c`$. The boundary of the contact area at the front side of the cylinder in the direction of motion follows from geometry (in the co-moving frame),
$$\xi _+=\sqrt{2Rh}$$
(5)
while the other boundary follows from the contact condition $`f\left(\xi _{}\right)=0`$:
$$\xi _{}=\frac{\gamma v}{k}\sqrt{2hR+\left(\frac{\gamma ^2}{k^2}2\frac{m}{k}\right)v^2}.$$
(6)
The springs at $`x=x_+`$ are special: at time $`t\delta `$ ($`\delta +0`$) their velocity is zero, while an infinitesimal time later it is finite. This singularity may be attributed to a finite force $`F_N^{}`$ acting at the point $`x_+`$. During the time $`\mathrm{d}t`$ the cylinder moves by $`v\mathrm{d}t`$ and accelerates springs of mass $`mv\mathrm{d}t`$, therefore
$$F_N^{}=\frac{\mathrm{d}p}{\mathrm{d}t}=\dot{y}\left(x_+\right)mv=\sqrt{\frac{2h}{R}}mv^2.$$
(7)
The dissipation of energy due to the plane deformation and the instantaneous acceleration of the plane material at $`\xi _+`$ in the instant of first contact reads
$$\dot{E}=_\xi _{}^{\xi _+}d\xi f\left(\xi \right)\dot{y}\left(\xi \right)\frac{m}{2}\dot{y}^2\left(\xi _+\right)vvF_R,$$
(8)
defining the force $`F_R`$ which acts against rolling
$$F_R=\frac{1}{R}_\xi _{}^{\xi _+}d\xi \xi f\left(\xi \right)+mv^2\frac{h}{R}.$$
(9)
To evaluate $`F_R`$ we have to determine the penetration depth $`h(v)`$. It is given implicitly by equilibrating the total force exerted by the plane to the cylinder and the weight of the cylinder,
$$_\xi _{}^{\xi _+}f\left(\xi \right)d\xi +F_N^{}=Mg\mathrm{cos}\alpha ,$$
(10)
where $`f\left(\xi \right)`$ and $`F_N^{}`$ are given by Eqs. (4) and (7). For an analytic evaluation of $`F_R`$ we assume $`\xi _{}=0`$. This approximation means that the surface recovers so slowly, that the recovering surface does not transmit energy back to the cylinder\[Greenwood et al. (1961)\]. Since for realistic parameters (see below) the effects of interest, namely bistability of the rolling velocity and noise controlled velocity, occur at rather large velocity, this approximation is justified. Then Eq. (10) reduces to $`Mg\mathrm{cos}\alpha =\gamma vh+\frac{2}{3}\sqrt{2R}kh^{3/2}`$, which, for large damping, $`v\gamma k\sqrt{2Rh}`$, yields
$$h(v)\frac{Mg\mathrm{cos}\alpha }{\gamma v+\frac{2}{3}k\sqrt{\frac{2RMg\mathrm{cos}\alpha }{\gamma v}}}.$$
(11)
With the same arguments from Eq. (9) we obtain the rolling friction force
$$F_R=\frac{kh^2}{2}+\frac{mv^2h}{R}+\frac{2\gamma vh}{3R}\sqrt{2Rh}.$$
(12)
Figure 2 shows $`F_R(v)`$ as given by Eq. (12) together with the numerical solution of Eqs. (5,6,9,10) for realistic material parameters, corresponding to soft rubber ($`m=`$100 kg/m<sup>2</sup>, $`k=10^7`$ kg/m<sup>2</sup>/s<sup>2</sup> , $`\gamma =5\times 10^5`$ kg/m<sup>2</sup>/s). The faster the cylinder moves, the more efficiently energy is dissipated due to increasing deformation rate. On the other hand, the faster it moves the smaller is the contact area, i.e. the amount of deformed material decreases. Consequently, $`F_R`$ depends non-monotonically on the velocity, which agrees with experimental results \[Olson et al. (2003), Yi et al. (2002), Gustafsson (1997), Ray (1997), Bakker et al. (1987), Burckhardt (1993), Yi & Jeong (1998)\].
## 3 CYLINDER ROLLING DOWN AN INCLINE
The rolling friction force $`F_R`$ is now applied to describe the motion of a cylinder rolling down a plane inclined by the angle $`\alpha `$. The cylinder is subjected to an external driving force $`F_{\mathrm{ex}}=Mg\mathrm{sin}\alpha `$ ($`g`$ is the acceleration of gravity). The rolling friction force $`F_R`$ and the viscous drag force $`F_D`$ due to the surrounding air counteract this motion. We also assume that the tangential force acting between the cylinder and the surface at the contact area is strong enough to keep it from sliding. Newtonโs equation for the cylinder reads
$$\left(M+I/R\right)\dot{v}=F_D(v)F_R(v)+F_{\mathrm{ex}}+\zeta (t).$$
(13)
The viscous drag force is given by
$$F_D=Av+Bv^2.$$
(14)
Detailed analysis \[Pรถschel et al. (2005)\] suggests $`A=0`$ and $`B=0.2`$ kg/m<sup>2</sup>. Figure 3 shows the total force acting on the cylinder, $`F(v)=F_DF_R+F_{\mathrm{ex}}`$.
The steady state condition, $`\dot{v}=0`$, may be fulfilled either for only one velocity (top and bottom curves in Fig. 3) or for three different velocities (middle curves). The former case implies a unique stationary velocity, while the latter allows for three stationary velocities. Only two of them, the smallest and the largest, correspond to stable motion ($`F/v<0`$). Modifying the inclination $`\alpha `$, the system may transit from one stable solution to the other. Figure 4 shows the bifurcation diagram.
As shown in Figs. 2-4 the analytical theory (dashed curves) agrees well with the numerical results if the damping parameter $`\gamma `$ is large.
## 4 NOISE-INDUCED JUMPS
Up to now, we did not consider noise, which is always present in realistic systems. The stochastic force $`\zeta (t)`$ describing fluctuations in the media is modeled by Gaussian white noise of zero average, $`\zeta (t)=0`$, and intensity $`\sigma `$: $`\zeta (t)\zeta (t^{})=\sigma ^2\delta (tt^{})`$. If the system has only one stable velocity, the addition of noise does not change the motion of the cylinder qualitatively. In this case the velocity fluctuates around the average value given by the steady state condition $`\dot{v}=0`$. Figure 5 (top row) shows the velocity of the cylinder as obtained from numerical integration of the stochastic equation (13). The corresponding velocity distribution reveals a single peak (Fig. 5, bottom row).
For parameters corresponding to bistable velocity, the presence of noise changes the system qualitatively due to stochastic jumps between the meta-stable velocities, Fig. 5 (middle column). Consequently, the velocity distribution has two well separated peaks.
Due to the nonlinear dependence of rolling friction on the velocity, the average velocity may increase or decrease with increasing noise level, Fig. 6.
This phenomenon may be understood by analyzing the potential $`U(v)\mathrm{d}F/\mathrm{d}v`$. Depending on the parameters, $`U(v)`$ has a double- or single-well shape \[Pรถschel et al. (2005)\], corresponding to the stationary state(s) of the velocity, see Fig. 3. In both cases, $`U(v)`$ is an asymmetric function in the surrounding of its minima where $`F(v)=0`$. In the absence of noise the velocity remains in its (meta-)stable state. Subjected to noise, however, the velocity of the cylinder fluctuates around this minimum. If $`U(v)`$ close to the minimum is steeper in the direction of lower velocities than in the direction of higher velocities, the average velocity will be shifted towards higher velocities and, thus, increases with increasing noise level (Fig. 6a). In the opposite case, large noise impedes rolling (Fig. 6b).
## 5 SUMMARY
We studied the rolling motion of a hard cylinder on an inclined viscoelastic plane in the presence of a surrounding medium. For certain realistic parameters the stationary velocity of the cylinder is bistable. For large damping of the planeโs viscous deformation the numerical results agree well with theory.
In the presence of noise as is unavoidable in any realistic experiment, by means of numerical simulations we found noise-induced transitions between the meta-stable velocities. Depending on the system parameters, increasing noise level may accelerate or decelerate the rolling motion.
The described effects may be important for technical systems where the presence of noise may lead to an effective increase of the mobility of a rolling body which is driven by an external force. |
warning/0506/hep-ph0506044.html | ar5iv | text | # Prompt photon photoproduction at HERA in the ๐_๐-factorization approach
## 1 Introduction
The prompt photon production in $`ep`$ collisions at HERA is a subject of the intensive studies \[1โ4\]. The theoretical and experimental investigations of the such processes have provided a direct probe of the hard subprocess dynamics, since produced photons are largely insensitive to the effects of final-state hadronization. Usually photons are called โpromptโ if they are coupled to the interacting quarks. From the theoretical point, these photons in $`ep`$ collisions can be produced via direct $`\gamma q\gamma q`$ and resolved production mechanisms. In resolved events, the photon emitted by the electron fluctuate into a hadronic state and a gluon and/or a quark of this hadronic fluctuation takes part in the hard interactions. Also observed final state photon may arise from fragmentation process , where a quark or gluon decays into $`\gamma `$. The cross section of such processes involves relative poorly known parton-to-photon fragmentation functions . However, the isolation criterion which introduced in experimental analyses substantially reduces the fragmentation component. In any case, for the theoretical description of prompt photon production at modern (HERA, Tevatron) and future (LHC) colliders the detailed knowledge of parton (quark and gluon) distributions in a proton and in a photon is necessary.
Usually quark and gluon densities are described by the Dokshitzer-Gribov-Lipatov-Altarelli-Parizi (DGLAP) evolution equation where large logarithmic terms proportional to $`\mathrm{ln}\mu ^2`$ are taken into account only. The cross sections can be rewritten in terms of process-dependent hard matrix elements convoluted with quark or gluon density functions. In this way the dominant contributions come from diagrams where parton emissions in initial state are strongly ordered in virtuality. This is called collinear factorization, as the strong ordering means that the virtuality of the parton entering the hard scattering matrix elements can be neglected compared to the large scale $`\mu `$.
However, at high energies (or small $`x\mu ^2/s1`$) effects of finite virtualities and transverse momenta of the incoming partons may become more and more important. These effects can be systematically accounted for in a $`k_T`$-factorization formalism \[9โ12\]. Just as for DGLAP, in this way it is possible to factorize an observable into a convolution of process-dependent hard matrix elements with universal parton distributions. But as the virtualities and transverse momenta of the emitted partons are no longer ordered, the matrix elements have to be taken off-shell and the convolution made also over transverse momentum $`๐ค_T`$ with the unintegrated (i.e. $`k_T`$-dependent) parton distributions. The unintegrated parton distribution $`f_a(x,๐ค_T^2)`$ determines the probability to find a type $`a`$ parton carrying the longitudinal momentum fraction $`x`$ and the transverse momentum $`๐ค_T`$. In particular, usage of the unintegrated parton densities have the advantage that it takes into account true kinematics of the process under consideration even at leading order.
The unintegrated parton distributions $`f_a(x,๐ค_T^2)`$ are a subject of intensive studies . Various approaches to investigate these quantities has been proposed. It is believed that at assymptotically large energies (or very small $`x`$) the theoretically correct description is given by the Balitsky-Fadin-Kuraev-Lipatov (BFKL) evolution equation where large terms proportional to $`\mathrm{ln}1/x`$ are taken into account. Another approach, valid for both small and large $`x`$, have been developed by Ciafaloni, Catani, Fiorani and Marchesini, and is known as the CCFM model . It introduces angular ordering of emissions to correctly treat gluon coherence effects. In the limit of asymptotic energies, it almost equivalent to BFKL \[17โ19\], but also similar to the DGLAP evolution for large $`x1`$. The resulting unintegrated gluon distribution depends on two scales, the additional scale $`\overline{q}`$ is a variable related to the maximum angle allowed in the emission and plays the role of the evolution scale $`\mu `$ in the collinear parton densities.
Also it is possible to obtain the two-scale involved unintegrated parton distributions from the conventional ones using the Kimber-Martin-Ryskin (KMR) prescription . In this way the $`\mu `$ dependence in the unintegrated parton distribution enters only in last step of the evolution, and single scale evolution equations can be used up to this step. Such procedure can be applied to a proton as well as photon and is expected to account for the main part of the collinear higher-order QCD corrections. The KMR-constructed parton densities were used, in particular, to describe the heavy quark production in $`\gamma \gamma `$ collisions at CERN LEP2 and prompt photon hadroproduction at fixed target experiments and at Fermilab Tevatron collider (in the double logarithmic approximation).
In the present paper we will apply the KMR method to obtain the unintegrated quark and gluon distributions in a proton $`f_a(x,๐ค_T^2,\mu ^2)`$ and in a photon $`f_a^\gamma (x,๐ค_T^2,\mu ^2)`$ independently from other authors. After that, we calculate the inclusive prompt photon photoproduction at HERA energies. Such calculations in the $`k_T`$-factorization approach will be performed for the first time. We will investigate the transverse energy $`E_T^\gamma `$ and pseudo-rapidity $`\eta ^\gamma `$ distributions of the produced prompt photons and compare our theoretical results with the recent experimental data taken by the H1 and ZEUS collaborations. In order to estimate the theoretical uncertainties of our predictions we will study the renormalization and factorization scale dependences of the calculated cross sections. Next we calculate the associated prompt photon and jet production rates using some physically motivated approximation. In order to investigate the underlying dynamics more detail, we will study the angular correlations between the prompt photon and jet in the transverse momentum plane. It was shown that theoretical and experimental studying of such quantities is a direct probe of the non-collinear parton evolution.
The additional motivation of our investigations within the $`k_T`$-factorization approach is the fact that the next-to-leading order (NLO) collinear QCD calculations are $`3040`$% below the data, especially in rear pseudo-rapidity (electron direction) region. So, one of the main goals of this paper is to investigate whether the $`k_T`$-factorization formalism could give a better description of the HERA data than collinear NLO QCD calculations.
The our paper is organized as follows. In Section 2 the KMR unintegrated parton densities in a proton and in a photon are obtained and their properties are discussed. In particular, we compare the KMR gluon distributions with ones taken from the full CCFM equation. In Section 3 we present the basic formulas with a brief review of calculation steps. In Section 4 we present the numerical results of our calculations. Finally, in Section 5, we give some conclusions. The compact analytic expressions for the off-mass shell matrix elements of all the subprocesses under consideration are given in Appendix. These formulas may be useful for the subsequent applications.
## 2 The KMR unintegrated partons
The Kimber-Martin-Ryskin approach is the formalism to construct parton distributions $`f_a(x,๐ค_T^2,\mu ^2)`$ unintegrated over the parton transverse momenta $`๐ค_T^2`$ from the known conventional parton distributions $`a(x,\mu ^2)`$, where $`a=xg`$ or $`a=xq`$. This formalism is valid for a proton as well as photon and can embody both DGLAP and BFKL contributions. It also accounts for the angular ordering which comes from coherence effects in gluon emission.
We start from parton distributions in a proton. The key observation here is that the $`\mu `$ dependence of the unintegrated distributions $`f_a(x,๐ค_T^2,\mu ^2)`$ enters at the last step of the evolution, and therefore single scale evolution equations (DGLAP or unified DGLAP-BFKL ) can be used up to this step. It was shown that the unintegrated distributions obtained via unified DGLAP-BFKL evolution are rather similar to those based on the pure DGLAP equations. It is because the imposion of the angular ordering constraint is more important than including the BFKL effects. Based on this point, in our calculations we will use much more simpler DGLAP equation up to the last evolution step. In this approximation, the unintegrated quark and gluon distributions are given by
$$\genfrac{}{}{0pt}{}{f_q(x,๐ค_T^2,\mu ^2)=T_q(๐ค_T^2,\mu ^2){\displaystyle \frac{\alpha _s(๐ค_T^2)}{2\pi }}\times }{\times {\displaystyle \underset{x}{\overset{1}{}}}dz[P_{qq}(z){\displaystyle \frac{x}{z}}q({\displaystyle \frac{x}{z}},๐ค_T^2)\mathrm{\Theta }(\mathrm{\Delta }z)+P_{qg}(z){\displaystyle \frac{x}{z}}g({\displaystyle \frac{x}{z}},๐ค_T^2)],}$$
$`(1)`$
$$\genfrac{}{}{0pt}{}{f_g(x,๐ค_T^2,\mu ^2)=T_g(๐ค_T^2,\mu ^2){\displaystyle \frac{\alpha _s(๐ค_T^2)}{2\pi }}\times }{\times {\displaystyle \underset{x}{\overset{1}{}}}dz[{\displaystyle \underset{q}{}}P_{gq}(z){\displaystyle \frac{x}{z}}q({\displaystyle \frac{x}{z}},๐ค_T^2)+P_{gg}(z){\displaystyle \frac{x}{z}}g({\displaystyle \frac{x}{z}},๐ค_T^2)\mathrm{\Theta }(\mathrm{\Delta }z)],}$$
$`(2)`$
where $`P_{ab}(z)`$ are the usual unregulated leading order DGLAP splitting functions, and $`q(x,\mu ^2)`$ and $`g(x,\mu ^2)`$ are the conventional quark and gluon densities. The theta functions which appear in (1) and (2) imply the angular-ordering constraint $`\mathrm{\Delta }=\mu /(\mu +|๐ค_T|)`$ specifically to the last evolution step to regulate the soft gluon singularities. For other evolution steps, the strong ordering in transverse momentum within the DGLAP equations automatically ensures angular ordering. It is important that parton distributions $`f_a(x,๐ค_T^2,\mu ^2)`$ extended now into the $`๐ค_T^2>\mu ^2`$ region. This fact is in the clear contrast with the usual DGLAP evolution<sup>1</sup><sup>1</sup>1We would like to note that cut-off $`\mathrm{\Delta }`$ can be also taken $`\mathrm{\Delta }=|๐ค_T|/\mu `$ . In this case the unintegrated parton distributions given by (1) โ (2) vanish for $`๐ค_T^2>\mu ^2`$ in accordance with the DGLAP strong ordering in $`๐ค_T^2`$..
The virtual (loop) contributions may be resummed to all orders by the quark and gluon Sudakov form factors
$$\mathrm{ln}T_q(๐ค_T^2,\mu ^2)=\underset{๐ค_T^2}{\overset{\mu ^2}{}}\frac{d๐ฉ_T^2}{๐ฉ_T^2}\frac{\alpha _s(๐ฉ_T^2)}{2\pi }\underset{0}{\overset{z_{\mathrm{max}}}{}}๐zP_{qq}(z),$$
$`(3)`$
$$\mathrm{ln}T_g(๐ค_T^2,\mu ^2)=\underset{๐ค_T^2}{\overset{\mu ^2}{}}\frac{d๐ฉ_T^2}{๐ฉ_T^2}\frac{\alpha _s(๐ฉ_T^2)}{2\pi }\left[n_f\underset{0}{\overset{1}{}}๐zP_{qg}(z)+\underset{z_{\mathrm{min}}}{\overset{z_{\mathrm{max}}}{}}๐zzP_{gg}(z)\right],$$
$`(4)`$
where $`z_{\mathrm{max}}=1z_{\mathrm{min}}=\mu /(\mu +|๐ฉ_T|)`$. The form factors $`T_a(๐ค_T^2,\mu ^2)`$ give the probability of evolving from a scale $`๐ค_T^2`$ to a scale $`\mu ^2`$ without parton emission. According to (3) and (4) $`T_a(๐ค_T^2,\mu ^2)=1`$ in the $`๐ค_T^2>\mu ^2`$ region.
We would like to note that such definition of the $`f_a(x,๐ค_T^2,\mu ^2)`$ is correct for $`๐ค_T^2>\mu _0^2`$ only, where $`\mu _01`$ GeV is the minimum scale for which DGLAP evolution of the collinear parton densities is valid. Everywhere in our numerical calculations we set the starting scale $`\mu _0`$ to be equal $`\mu _0=1`$ GeV. Since the starting point of our derivation is the leading order DGLAP equations, the unintegrated parton distributions must satisfy the normalisation condition
$$a(x,\mu ^2)=\underset{0}{\overset{\mu ^2}{}}f_a(x,๐ค_T^2,\mu ^2)๐๐ค_T^2.$$
$`(5)`$
This relation will be exactly satisfied if we define
$$f_a(x,๐ค_T^2,\mu ^2)|_{๐ค_T^2<\mu _0^2}=a(x,\mu _0^2)T_a(\mu _0^2,\mu ^2).$$
$`(6)`$
Then, we have obtained the unintegrated parton distributions $`f_a(x,๐ค_T^2,\mu ^2)`$ in a proton. In order to obtain the unintegrated parton distribution $`f_a^\gamma (x,๐ค_T^2,\mu ^2)`$ in a photon the same formulas (1) โ (4) can be also used . In the last case the conventional parton distributions $`a(x,\mu ^2)`$ in a proton should be replaced by the corresponding parton densities $`a^\gamma (x,\mu ^2)`$ in a photon.
In Figure 1 we show unintegrated parton densities $`f_a(x,๐ค_T^2,\mu ^2)`$ in a proton at scale $`\mu ^2=100\mathrm{GeV}^2`$ as a function of $`x`$ for different values of $`๐ค_T^2`$, namely $`๐ค_T^2=2\mathrm{GeV}^2`$ (a), $`๐ค_T^2=5\mathrm{GeV}^2`$ (b), $`๐ค_T^2=10\mathrm{GeV}^2`$ (c) and $`๐ค_T^2=20\mathrm{GeV}^2`$ (d). The solid, dashed, short dashed, dotted and dash-dotted lines correspond to the unintegrated $`u+\overline{u}`$, $`d+\overline{d}`$, $`s`$, $`c`$ and gluon (divided by factor $`10`$) distributions, respectively. We have used here the standard GRV (LO) parametrizations of the collinear quark and gluon densities $`a(x,\mu ^2)`$. In order to be sure that normalization condition (5) is correctly satisfied we have performed the numerical integration of the parton densities $`f_a(x,๐ค_T^2,\mu ^2)`$ over transverse momenta $`๐ค_T^2`$. So, in Figure 2 we show our result for effective $`u+\overline{u}`$ quark and gluon (also divided by factor 10) distributions for different scales $`\mu ^2=2\mathrm{GeV}^2`$ (a), $`\mu ^2=5\mathrm{GeV}^2`$ (b), $`\mu ^2=10\mathrm{GeV}^2`$ (c), $`\mu ^2=20\mathrm{GeV}^2`$ (d). The solid lines correspond to the effective densities obtained from the unintegrated ones using relation (5). The dashed lines correspond to the collinear GRV (LO) parton distributions. One can see that normalization condition (5) is exactly satisfied practically for all $`x`$ and $`\mu ^2`$ values. There are only rather small (less then few percent) violations of (5) in the case of the $`u+\overline{u}`$ quark distributions at large $`x>0.2`$. We have checked numerically that the expression (5) is true also for other parton densities in a proton and in a photon.
For comparison we plot in Figure 2 (as dash-dotted lines) the corresponding LO parton distributions obtained by the CTEQ collaboration (CTEQ5L set). It is clear that there are some differences in both normalization and shape between GRV and CTEQ parametrizations. In general, the CTEQ curves lie below the GRV ones by about $`10`$%. This difference tends to be small when scale $`\mu ^2`$ is large. However, the CTEQ collaboration does not provide a set of the parton distributions in a photon, which are necessary for calculation of the resolved contributions. Therefore everywhere in our numerical analysis we will use only the GRV parametrizations.
It is interesting to compare the KMR-constructed unintegrated parton densities with the distributions obtained in other approaches. Recently the full CCFM equation in a proton and in a photon was solved numerically using a Monte Carlo method, and new fits of the unintegrated gluon distributions (J2003 set 1 โ 3) have been presented . The input parameters were fitted to describe the proton structure function $`F_2(x,Q^2)`$. These unintegrated gluon densities were used also in description of the forward jet production at HERA, charm and bottom production at Tevatron , and charm and $`J/\psi `$ production at LEP2 energies . In Figure 3 we plot KMR (as solid lines) and J2003 set 1 (as dashed lines) gluon distributions in a proton at scale $`\mu ^2=100\mathrm{GeV}^2`$ as a function of $`x`$ for different values of $`๐ค_T^2`$, namely $`๐ค_T^2=2\mathrm{GeV}^2`$ (a), $`๐ค_T^2=10\mathrm{GeV}^2`$ (b), $`๐ค_T^2=20\mathrm{GeV}^2`$ (c) and $`๐ค_T^2=50\mathrm{GeV}^2`$ (d). One can see that J2003 set 1 gluon density is less steep at small $`x`$ compared to the KMR one. The KMR gluon lie below J2003 set 1 at small $`๐ค_T^2`$ region for $`x>310^3`$. Typically the difference between solid and dashed lines is about 30% โ 40% at $`x=0.01`$. This fact results to the some underestimation of the calculated cross sections in the KMR approach. This underestimation is about 30% at HERA and 50% at Tevatron<sup>2</sup><sup>2</sup>2See Ref. for more details.. Therefore we can expect a rather large sensitivity of our predictions to the parton evolution scheme.
We would like to point out again that behaviour of different unintegrated parton distributions in a proton in the small $`๐ค_T^2`$ region (which essentially drives the total cross sections) is very different, as it is clear shown on Figure 3. However, the CCFM evolution does not include quark-initiated chains and therefore can not be used in our analysis since prompt photon production at HERA strongly depends on the quark distributions (see Section 3). Therefore in our following investigations we will use only the KMR unintegrated parton densities. But one should remember that dependence of our results on the evolution scheme may be rather large, and further theoretical attempts (in order to investigate the unintegrated quark distributions more detail) are necessary to reduce this uncertainty.
## 3 Calculation details
### 3.1 The subprocesses under consideration
In $`ep`$ collisions at HERA prompt photons can be produced by one of three mechanisms: a direct production, a single resolved production and via parton-to-photon fragmentation processes . The direct contribution to the $`\gamma p\gamma +X`$ process is the Deep Inelastic Compton (DIC) scattering on the quark (antiquark)
$$\gamma (k_1)+q(k_2)\gamma (p_\gamma )+q(p^{}),$$
$`(7)`$
where the particles four-momenta are given in parentheses. It gives the $`๐ช(\alpha _{em}^2)`$ order contribution to the hadronic cross section. Here $`\alpha _{em}`$ is Sommerfeldโs fine structure constant. The single resolved processes are
$$q(k_1)+g(k_2)\gamma (p_\gamma )+q(p^{}),$$
$`(8)`$
$$g(k_1)+q(k_2)\gamma (p_\gamma )+q(p^{}),$$
$`(9)`$
$$q(k_1)+q(k_2)\gamma (p_\gamma )+g(p^{}).$$
$`(10)`$
Since the parton distributions in a photon $`a_\gamma (x,\mu ^2)`$ have a leading behavior proportional to $`\alpha _{em}\mathrm{ln}\mu ^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\alpha _{em}/\alpha _s`$, these subprocesses give also the $`๐ช(\alpha _{em}^2)`$ contributions and therefore should be taken into account in our analysis.
In addition to the direct and resolved production, photons can be also produced through the fragmentation of a hadronic jet into a single photon carrying a large fraction $`z`$ of the jet energy . These processes are described in terms of quark-to-photon $`D_{q\gamma }(z,\mu ^2)`$ and gluon-to-photon $`D_{g\gamma }(z,\mu ^2)`$ fragmentation functions . The main feature of the fragmentation contribution in leading order is fact that produced photon balanced by a jet on the opposite side of the event and accompanied by collinear hadrons on the same side of the event.
It is very important that in order to reduce the huge background from the secondary photons produced by the decays of $`\pi ^0`$, $`\eta `$ and $`\omega `$ mesons the isolation criterion is introduced in the experimental analyses. This criterion is the following. A photon is isolated if the amount of hadronic transverse energy $`E_T^{\mathrm{had}}`$, deposited inside a cone with aperture $`R`$ centered around the photon direction in the pseudo-rapidity and azimuthal angle plane, is smaller than some value $`E_T^{\mathrm{max}}`$:
$$\genfrac{}{}{0pt}{}{E_T^{\mathrm{had}}E_T^{\mathrm{max}},}{(\eta \eta ^\gamma )^2+(\varphi \varphi ^\gamma )^2R^2.}$$
$`(11)`$
The both H1 and ZEUS collaborations take $`R=1`$, $`E_T^{\mathrm{max}}=ฯตE_T^\gamma `$ with $`ฯต=0.1`$ in the experiment \[1โ4\]. Isolation not only reduces the background but also significantly reduces the fragmentation components. It was shown that after applying the isolation cut (11) the contribution from the fragmentation subprocesses is about 5 โ 6% of the total cross section. Since the dependence of our results on the non-collinear parton evolution scheme may be rather large (as it was demonstrated in Section 2), in our further analysis we will neglect the relative small fragmentation contribution and consider only the direct and resolved production (7) โ (10). We note that photon produced in these processes is automatically isolated from the quark or gluon jet by requiring a non-zero transverse momentum of a photon or jet in the $`\gamma p`$ center-of-mass frame.
It was claimed that direct box diagram $`\gamma g\gamma g`$, which is formally of the next-to-next-to-leading order (NNLO), gives approximately 6% contribution to the total NLO cross section. In the present paper we will not take into account this diagram also.
### 3.2 Cross section for prompt photon production
Let $`p_e`$ and $`p_p`$ be the four-momenta of the initial electron and proton. The direct contribution (7) to the $`\gamma p\gamma +X`$ process in the $`k_T`$-factorization approach can be written as
$$d\sigma ^{(\mathrm{dir})}(\gamma p\gamma +X)=\underset{q}{}\frac{dx_2}{x_2}f_q(x_2,๐ค_{2T}^2,\mu ^2)๐๐ค_T^2\frac{d\varphi _2}{2\pi }๐\widehat{\sigma }(\gamma q\gamma q),$$
$`(12)`$
where $`\widehat{\sigma }(\gamma q\gamma q)`$ is the hard subprocess cross section via quark or antiquark having fraction $`x_2`$ of a initial proton longitudinal momentum, non-zero transverse momentum $`๐ค_{2T}`$ ($`๐ค_{2T}^2=k_{2T}^20`$) and azimuthal angle $`\varphi _2`$. The expression (12) can be easily rewritten in the form
$$\genfrac{}{}{0pt}{}{\sigma ^{(\mathrm{dir})}(\gamma p\gamma +X)={\displaystyle \underset{q}{}}{\displaystyle }{\displaystyle \frac{E_T^\gamma }{8\pi (x_2s)^2(1\alpha )}}|\overline{}|^2(\gamma q\gamma q)\times }{\times f_q(x_2,๐ค_{2T}^2,\mu ^2)dy^\gamma dE_T^\gamma d๐ค_{2T}^2{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle \frac{d\varphi ^\gamma }{2\pi }},}$$
$`(13)`$
where $`|\overline{}|^2(\gamma q\gamma q)`$ is the hard matrix element which depends on the transverse momentum $`๐ค_{2T}^2`$, $`s=(k_1+p_p)^2`$ is the total energy of the subprocess under consideration, $`y^\gamma `$, $`E_T^\gamma `$ and $`\varphi ^\gamma `$ are the rapidity, transverse energy and azimuthal angle of the produced photon in the $`\gamma p`$ center-of-mass frame, and $`\alpha =E_T^\gamma \mathrm{exp}y^\gamma /\sqrt{s}`$.
The formula for the resolved contribution to the prompt photon photoproduction in the $`k_T`$-factorization approach can be obtained by the similar way. But one should keep in mind that convolution in (12) should be made also with the unintegrated parton distributions $`f_a^\gamma (x,๐ค_T^2,\mu ^2)`$ in a photon, i.e.
$$\genfrac{}{}{0pt}{}{d\sigma ^{(\mathrm{res})}(\gamma p\gamma +X)={\displaystyle \underset{a,b}{}}{\displaystyle }{\displaystyle \frac{dx_1}{x_1}}f_a^\gamma (x_1,๐ค_{1T}^2,\mu ^2)d๐ค_{1T}^2{\displaystyle \frac{d\varphi _1}{2\pi }}\times }{\times {\displaystyle }{\displaystyle \frac{dx_2}{x_2}}f_b(x_2,๐ค_{2T}^2,\mu ^2)d๐ค_{2T}^2{\displaystyle \frac{d\varphi _2}{2\pi }}d\widehat{\sigma }(ab\gamma c),}$$
$`(14)`$
where $`a,b,c=q`$ and/or $`g`$, $`\widehat{\sigma }(ab\gamma c)`$ is the cross section of the photon production in the corresponding parton-parton interaction (8) โ (10). Here parton $`a`$ has fraction $`x_1`$ of a initial photon longitudinal momentum, non-zero transverse momentum $`๐ค_{1T}`$ ($`๐ค_{1T}^2=k_{1T}^20`$) and azimuthal angle $`\varphi _1`$. We can easily obtain the final expression from equation (14). It has the form
$$\genfrac{}{}{0pt}{}{\sigma ^{(\mathrm{res})}(\gamma p\gamma +X)={\displaystyle \underset{a,b}{}}{\displaystyle }{\displaystyle \frac{E_T^\gamma }{8\pi (x_1x_2s)^2}}|\overline{}|^2(ab\gamma c)\times }{\times f_a^\gamma (x_1,๐ค_{1T}^2,\mu ^2)f_b(x_2,๐ค_{2T}^2,\mu ^2)d๐ค_{1T}^2d๐ค_{2T}^2dE_T^\gamma dy^\gamma dy^c{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle \frac{d\varphi ^\gamma }{2\pi }},}$$
$`(15)`$
where $`y^c`$ is the rapidity of the parton $`c`$ in the $`\gamma p`$ center-of-mass frame. It is important that hard matrix elements $`|\overline{}|^2(ab\gamma c)`$ depend on the transverse momenta $`๐ค_{1T}^2`$ and $`๐ค_{2T}^2`$. We would like to note that if we average the expressions (13) and (15) over $`๐ค_{1T}`$ and $`๐ค_{2T}`$ and take the limit $`๐ค_{1T}^20`$ and $`๐ค_{2T}^20`$, then we obtain well-known expressions for the prompt photon production in leading-order (LO) perturbative QCD.
The experimental data taken by the H1 and ZEUS collaborations refer to the prompt photon production in the $`ep`$ collisions, where electron is scattered at small angle and the mediating photon is almost real ($`Q^20`$). Therefore $`\gamma p`$ cross sections (13) and (15) needs to be weighted with the photon flux in the electron:
$$d\sigma (epe^{}+\gamma +X)=f_{\gamma /e}(y)๐\sigma (\gamma p\gamma +X)๐y,$$
$`(16)`$
where $`y`$ is a fraction of the initial electron energy taken by the photon in the laboratory frame, and we use the Weizacker-Williams approximation for the bremsstrahlung photon distribution from an electron:
$$f_{\gamma /e}(y)=\frac{\alpha _{em}}{2\pi }\left(\frac{1+(1y)^2}{y}\mathrm{ln}\frac{Q_{\mathrm{max}}^2}{Q_{\mathrm{min}}^2}+2m_e^2y\left(\frac{1}{Q_{\mathrm{max}}^2}\frac{1}{Q_{\mathrm{min}}^2}\right)\right).$$
$`(17)`$
Here $`m_e`$ is the electron mass, $`Q_{\mathrm{min}}^2=m_e^2y^2/(1y)^2`$ and $`Q_{\mathrm{max}}^2=1\mathrm{GeV}^2`$, which is a typical value for the recent photoproduction measurements at the HERA collider.
The multidimensional integration in (13), (15) and (16) has been performed by means of the Monte Carlo technique, using the routine VEGAS . The full C$`++`$ code is available from the authors on request<sup>3</sup><sup>3</sup>3lipatov@theory.sinp.msu.ru. For readerโs convenience, we collect the analytic expressions for the off-shell matrix elements which correspond to all partonic subprocesses under consideration (7) โ (10) in the Appendix. These formulas may be useful for the subsequent applications.
## 4 Numerical results
We now are in a position to present our numerical results. First we describe our theoretical input and the kinematical conditions. After we fixed the unintegrated parton distributions in a proton $`f_a(x,๐ค_T^2,\mu ^2)`$ and in a photon $`f_a^\gamma (x,๐ค_T^2,\mu ^2)`$, the cross sections (13) and (15) depend on the energy scale $`\mu `$. As it often done for prompt photon production, we choose the renormalization and factorization scales to be $`\mu =\xi E_T^\gamma `$. In order to estimate the theoretical uncertainties of our calculations we will vary the scale parameter $`\xi `$ between 1/2 and 2 about the default value $`\xi =1`$. Also we use LO formula for the strong coupling constant $`\alpha _s(\mu ^2)`$ with $`n_f=3`$ active (massless) quark flavours and $`\mathrm{\Lambda }_{\mathrm{QCD}}=232`$ MeV, such that $`\alpha _s(M_Z^2)=0.1169`$. In our analysis we will not neglect the charm quark mass and set it to be $`m_c=1.4`$ GeV.
### 4.1 Inclusive prompt photon production
Experimental data for the inclusive prompt photon production $`epe^{}+\gamma +X`$ comes from both ZEUS and H1 collaborations. Two differential cross section are determined: first as a function of the transverse energy $`E_T^\gamma `$, and second as a function of pseudo-rapidity $`\eta ^\gamma `$. The ZEUS data refer to the kinematic region<sup>4</sup><sup>4</sup>4Here and in the following all kinematic quantities are given in the laboratory frame where positive OZ axis direction is given by the proton beam. defined by $`5<E_T^\gamma <10`$ GeV and $`0.7<\eta ^\gamma <0.9`$ with electron energy $`E_e=27.5`$ GeV and proton energy $`E_p=820`$ GeV. The fraction $`y`$ of the electron energy trasferred to the photon is restricted to the range $`0.2<y<0.9`$. Additionally the available ZEUS data for the prompt photon pseudo-rapidity distributions have been given also for three subdivisons of the $`y`$ range, namely $`0.2<y<0.32`$ ($`134<W<170`$ GeV), $`0.32<y<0.5`$ ($`170<W<212`$ GeV) and $`0.5<y<0.9`$ ($`212<W<285`$ GeV). The more recent H1 data refer to the kinematic region defined by $`5<E_T^\gamma <10`$ GeV, $`1<\eta ^\gamma <0.9`$ and $`0.2<y<0.7`$ with electron energy $`E_e=27.6`$ GeV and proton energy $`E_p=920`$ GeV.
The transverse energy distributions of the inclusive prompt photon for different kinematical region are shown in Figures 4 and 5 in comparison to the HERA data. Instead of presenting our theoretical predictions as continuous lines, we adopt the binning pattern encoded in the experimental data. The solid histograms obtained by fixing both the factorization and normalization scales at the default value $`\mu =E_T^\gamma `$, whereas upper and lower dashed histograms correspond to the $`\mu =E_T^\gamma /2`$ and $`\mu =2E_T^\gamma `$ scales, respectively. One can see that predicted cross sections agree well with the experimental data except the moderate $`E_T^\gamma `$ region. We would like to note that overall agreement with data can be improved when unintegrated quark and gluon distributions in a proton and in a photon will be studied more detail. It is because the KMR approach tends to underestimate the calculated cross sections, as it was discussed in Section 2. The collinear NLO QCD calculations give the similar description of the transverse energy distributions measured by the ZEUS collaboration. At the same time, according to the analysis which was done by the H1 collaboration , in order to obtain a realistic comparison of their data and theory the corrections for hadronisation and multiple interactions should be taken into account in the predictions<sup>5</sup><sup>5</sup>5See Ref. for more details.. The correction factors are typically 0.7 โ 0.9 depending on a bin. The NLO calculations are approximately 30% โ 40% below the H1 data if the corrections for hadronisation and multiple interactions are applied. We would like to note that these corrections are not accounted for in our analysis. The effect of scale variations in transverse energy distributions is rather large: the relative difference between results for $`\mu =E_T^\gamma `$ and results for $`\mu =E_T^\gamma /2`$ or $`\mu =2E_T^\gamma `$ is about 15%.
The pseudo-rapidity distributions of the inclusive prompt photon production compared with the HERA data in different kinematical region are shown in Figures 6 and 7. All histograms here are the same as in Figure 4. One can see that measured distributions are reasonably well described in the pseudo-rapidity region $`0.4\eta ^\gamma 0.9`$ only. For $`1\eta ^\gamma 0.4`$ our predictions lie mostly below the experimental points<sup>6</sup><sup>6</sup>6Note that such disagreement between predicted and measured cross sections is observed for collinear NLO QCD calculations also.. The discrepancy between data and theory at negative $`\eta ^\gamma `$ is found to be relative strong at low values of the initial photon fractional momentum $`y`$. So, in Figures 8, 9 and 10 we show the inclusive cross sections $`d\sigma /d\eta ^\gamma `$ evaluated for the three $`y`$ ranges $`0.2<y<0.32`$ ($`134<W<170`$ GeV), $`0.32<y<0.5`$ ($`170<W<212`$ GeV) and $`0.5<y<0.9`$ ($`212<W<285`$ GeV), respectively. All histograms here are the same as in Figure 4. In the lowest $`y`$ range, both our predictions and experimental data show a peaking at negative $`\eta ^\gamma `$, but it is stronger in the data. In the high $`y`$ region, $`0.5<y<0.9`$, a good agreement is obtained. This fact allows to establish that the above discussed discrepancy between the data and theory at $`1\eta ^\gamma 0.4`$ is coming from the low ($`0.2<y<0.32`$) and medium ($`0.32<y<0.5`$) $`y`$ region. The scale variation changes the estimated cross sections by about 15%. The collinear NLO QCD calculations give the similar description of the pseudo-rapidity distributions measured by the ZEUS collaboration. At the same time, after corrections for hadronisation and multiple interactions (not accounted for in our analysis) the NLO predictions are 30% โ 40% below the H1 data.
As it was already mentioned above, the dependence of the our results on a renormalization/factorization scale $`\mu `$ is rather large, about 10% โ 15% in the wide kinematic range. There are also additional uncertainties come from the unintegrated parton densities, as it was discussed in Section 2. The theoretical uncertainties of the collinear NLO QCD calculations are about 3% . This fact indicates that contribution from NNLO and high order terms is not significant. At the same time the strong scale dependence of our results demonstrates the necessarity of reducing of uncertainties in the non-collinear parton evolution.
The individual contributions from the direct and resolved production mechanisms to the total cross section in the $`k_T`$-factorization approach is about 47% and 53%, respectively. In the resolved conribution, the channels (8), (9) and (10) account for 80%, 14% and 6%. Additionally, using the Duke-Owens (DO) parton-to-photon fragmentation functions, we perform the estimation of the fragmentation component (not shown in Figures). We find that after applying isolation cut it give only a very small (about few percent) contribution.
### 4.2 Prompt photon production in association with jet
Now we demonstrate how $`k_T`$-factorization approach can be used to calculate the semi-inclusive prompt photon production rates. The produced photon is accompanied by a number of partons radiated in the course of the parton evolution. As it has been noted in Ref. , on the average the parton transverse momentum decreases from the hard interaction block towards the proton. As an approximation, we assume that the parton $`k^{}`$ emitted in the last evolution step compensates the whole transverse momentum of the parton participating in the hard subprocess, i.e. $`๐ค_{}^{}{}_{T}{}^{}๐ค_T`$. All the other emitted partons are collected together in the proton remnant, which is assumed to carry only a negligible transverse momentum compared to $`๐ค_{}^{}{}_{T}{}^{}`$. This parton gives rise to a final hadron jet with $`E_T^{\mathrm{jet}}=|๐ค_{}^{}{}_{T}{}^{}|`$ in addition to the jet produced in the hard subprocess. From these hadron jets we choose the one carrying the largest transverse energy, and then compute prompt photon with an associated jet cross sections.
Experimental data for such processes were obtained very recently by the H1 collaboration. The cross sections measured differentially as a function of $`E_T^\gamma `$, $`E_T^{\mathrm{jet}}`$, and the pseudo-rapidities $`\eta ^\gamma `$ and $`\eta ^{\mathrm{jet}}`$ in the kinematic region defined by $`5<E_T^\gamma <10`$ GeV, $`E_T^{\mathrm{jet}}>4.5`$ GeV, $`1<\eta ^\gamma <0.9`$, $`1<\eta ^{\mathrm{jet}}<2.3`$ and $`0.2<y<0.7`$ with electron energy $`E_e=27.6`$ GeV and proton energy $`E_p=920`$ GeV. There are no ZEUS data for the prompt photon plus jet production, although some data for distribution of events, not corrected for the detector effects, were presented .
The transverse energy $`E_T^\gamma `$ and pseudo-rapidity $`\eta ^\gamma `$ distributions of the prompt photon plus jet production are shown in Figures 11 and 12 in comparison with H1 data. All histograms here are the same as in Figure 4. In contrast to the inclusive case, one can see that our predictions are consistent with the data in most bins, although some discrepancies are present. The scale variation as it was described above changes the estimated cross sections by about 10%. The results of the collinear NLO calculations which include corrections for hadronisation and multiple interactions give the similar results and consistent with data also.
In Figures 13 and 14 we show our predictions for the transverse energy $`E_T^{\mathrm{jet}}`$ and pseudo-rapidity $`\eta ^{\mathrm{jet}}`$ distributions in comparison with H1 data. All histograms here are the same as in Figure 4. A rather good agreement between our results and data is obtained again. It is interesting to note that shape of the predicted pseudo-rapidity $`\eta ^{\mathrm{jet}}`$ distribution coincide with the one obtained in the collinear NLO calculations . At the same time the shape of this distribution is not reproduced by the leading-order QCD calculations . This fact can demonstrate again that the main part of the collinear high-order corrections is already included at LO level in $`k_T`$-factorization formalism. The scale dependence of our predictions is about 10%.
The most important variables for testing the structure of colliding proton and photon are the longitudinal fractional momenta of partons in these particles. In order to reconstruct the momentum fractions of the initial partons from measured quantities the observables $`x_\gamma `$ and $`x_p`$ are introduced :
$$x_\gamma =\frac{E_T^\gamma (e^{\eta ^\gamma }+e^{\eta ^{\mathrm{jet}}})}{2yE_e},x_p=\frac{E_T^\gamma (e^{\eta ^\gamma }+e^{\eta ^{\mathrm{jet}}})}{2E_p}.$$
$`(18)`$
These observables make explicit use only of the photon energy, which is better measured than the jet energy. The $`x_\gamma `$ distribution is particularly sensitive to the photon structure function. At large $`x_\gamma `$ region ($`x_\gamma >0.85`$) the cross section is dominated by the contribution of processes with direct initial photons, whereas at $`x_\gamma <0.85`$ the resolved photon contributions dominate .
So, in Figures 15 and 16 the $`x_\gamma `$ and $`x_p`$ distributions are shown in comparison with H1 data. One can see that our predictions reasonable agree with experimental data. The NLO calculations without corrections for hadronisation and multiple interactions give the similar results. However, NLO calculations tend to underestimate the H1 data if these corrections are taken into account. The hadronic and multiple interaction corrections improve the description of the data at $`x_\gamma <0.6`$ only .
Further understanding of the process dynamics and in particular of the high-order correction effects may be obtained from the transverse correlation between the produced prompt photon and the jet. The H1 collaboration has measured the distribution on the component of the prompt photonโs momentum perpendicular to the jet direction in the transverse plane, i.e.
$$p_T=|๐ฉ_T^\gamma \times ๐ฉ_T^{\mathrm{jet}}|/|๐ฉ_T^{\mathrm{jet}}|=E_T^\gamma \mathrm{sin}\mathrm{\Delta }\varphi ,$$
$`(19)`$
where $`\mathrm{\Delta }\varphi `$ is the difference in azimuth between the photon and the jet. In the collinear leading order approximation, the distribution over $`p_T`$ must be simply a delta function $`\delta (p_T)`$, since the produced photon and the jet are back-to-back in the transverse plane. Taking into account the non-vanishing initial parton transverse momenta $`๐ค_{1T}`$ and $`๐ค_{2T}`$ leads to the violation of this back-to-back kinematics in the $`k_T`$-factorization approach.
The normalised $`p_T`$ distributions are shown in Figures 17 and 18 separately for the regions $`x_\gamma <0.85`$ and $`x_\gamma >0.85`$, where direct and resolved photon induced processes dominate, respectively. All histograms here are the same as in Figure 4. Our predictions are consistent with the H1 data for all $`x_\gamma `$ values except large $`p_T`$ region. So, at $`p_T>5`$ GeV the results of our calculations lie slightly below the data at $`x_\gamma <0.85`$ and above the data at $`x_\gamma >0.85`$. At the same time the NLO QCD prediction gives a better description of the $`p_T`$ distributions at $`x_\gamma <0.85`$ than another one . It is because in this region the cross section is dominated by $`๐ช(\alpha _s)`$ corrections to the processes with resolved photons, which are not included in the NLO calculations . In general, we can conclude that our results lie between the predictions and the predictions in the whole $`x_\gamma `$ range. This fact indicates again that the main part of the high-order collinear corrections is effectively included in our calculations.
Finally, we would like to note that there are, of course, still rather large theoretical uncertainties in our results connected with unintegrated parton distributions, and it is necessary to work hard until these uncertainties will be reduced. However, it was shown that the properties of different unintegrated parton distributions clear manifest themselves in the azimuthal correlation between transverse momenta of the final state particles. Therefore we can expect that further theoretical and experimental studying of these correlations will give important information about non-collinear parton evolution dynamics in a proton and in a photon.
## 5 Conclusions
We have investigated the prompt photon photoproduction at the HERA collider in the $`k_T`$-factorization approach. In order to obtain the unintegrated quark and gluon distributions in a proton and in a photon we used the Kimber-Martin-Ryskin prescription. We have investigated both inclusive and associated with jet prompt photon production rates. Such calculations in the $`k_T`$-factorization approach were performed for the first time.
We took into account both the direct and resolved contributions and investigated the sensitivity of the our results to renormalization and factorization scales. There are, of course, also theoretical uncertainties due to non-collinear evolution scheme. However, much more work needs to be done before these uncertainties will be reduced.
We have found that our predictions for the inclusive prompt photon production are in reasonable agreement with the H1 and ZEUS data except rear (electron direction) pseudo-rapidity region. In contrast, our results for prompt photon associated with jet are consistent with data in the whole kinematical range. However, the scale dependence of our results is rather large compared to the collinear NLO QCD calculations. At the same time we demonstrate that main part of the standard high-order corrections is already included in the $`k_T`$-factorization formalism at LO level.
Note that in our analysis we neglect the contribution from the fragmentation processes and from the direct box diagram ($`\gamma g\gamma g`$). Since the relative large box contribution (about 6% of the total NLO cross section) is mainly due to large gluonic content of the proton at small $`x`$, studying of this subprocess should be also very interesting in the $`k_T`$-factorization approach. We plan to investigate it in detail in the forthcoming publications.
## 6 Acknowledgements
The authors are very grateful also to S.P. Baranov for encouraging interest and helpful discussions. This research was supported in part by the FASI of Russian Federation.
## 7 Appendix
Here we present the compact analytic expressions for the hard matrix elements which appear in (13) and (15). In the following, $`\widehat{s}`$, $`\widehat{t}`$ and $`\widehat{u}`$ are usual Mandelstam variables for corresponding $`22`$ subprocesses and $`e_q`$ is the fractional electric charge of quark $`q`$.
We start from the direct subprocess (7). The corresponding squared matrix element summed over final polarization states and averaged over initial ones read
$$|\overline{}|^2(\gamma q\gamma q)=\frac{2(4\pi )^2\alpha _{em}^2e_q^4}{(\widehat{s}m^2)^2(\widehat{u}m^2)^2}F_{\gamma q}(๐ค_{2T}^2),$$
$`(A.1)`$
where $`m`$ is the quark mass, and
$$\genfrac{}{}{0pt}{}{F_{\gamma q}(๐ค_T^2)=6m^8(3\widehat{s}^2+14\widehat{s}\widehat{u}+3\widehat{u}^2)m^4+(\widehat{s}^3+7\widehat{s}^2\widehat{u}+}{7\widehat{u}^2\widehat{s}+\widehat{u}^3)m^2(\widehat{s}^2+\widehat{u}^2)\widehat{s}\widehat{u}.}$$
$`(A.2)`$
It is important to note that when we calculate the Diracโs traces we set the incoming quark four-momentum to be equal $`k_2=x_2p_p`$. Therefore these formulas formally are the same as in the usual leading-order collinear approach and there is no obvious dependence on the parton transverse momentum $`๐ค_{2T}`$. However, this dependence is present because we have used true off-shell kinematics in order to estimate the cross section (13). It is in the clear contrast with the collinear calculations.
The squared matrix elements of the resolved photon contributions (8) โ (10) summed over final polarization states and averaged over initial ones read
$$|\overline{}|^2(qg\gamma q)=\frac{(4\pi )^2\alpha _{em}\alpha _se_q^2}{3(\widehat{t}m^2)^2(\widehat{s}m^2)^2}F_{qg}(๐ค_{1T}^2,๐ค_{2T}^2),$$
$`(A.3)`$
$$|\overline{}|^2(gq\gamma q)=\frac{(4\pi )^2\alpha _{em}\alpha _se_q^2}{3(\widehat{s}m^2)^2(\widehat{u}m^2)^2}F_{gq}(๐ค_{1T}^2,๐ค_{2T}^2),$$
$`(A.4)`$
$$|\overline{}|^2(qq\gamma g)=\frac{8(4\pi )^2\alpha _{em}\alpha _se_q^2}{9(\widehat{t}m^2)^2(\widehat{u}m^2)^2}F_{qq}(๐ค_{1T}^2,๐ค_{2T}^2),$$
$`(A.5)`$
where functions $`F_{qg}(๐ค_{1T}^2,๐ค_{2T}^2)`$, $`F_{gq}(๐ค_{1T}^2,๐ค_{2T}^2)`$ and $`F_{qq}(๐ค_{1T}^2,๐ค_{2T}^2)`$ are given by
$$F_{qg}(๐ค_{1T}^2,๐ค_{2T}^2)=6m^8(2๐ค_{2T}^4+2(\widehat{s}+\widehat{t})๐ค_{2T}^2+3\widehat{s}^2+3\widehat{t}^2+14\widehat{s}\widehat{t})m^4+$$
$$(2(\widehat{s}+\widehat{t})๐ค_{2T}^4+8\widehat{s}\widehat{t}๐ค_{2T}^2+\widehat{s}^3+\widehat{t}^3+7\widehat{s}\widehat{t}^2+7\widehat{s}^2\widehat{t})m^2$$
$$\widehat{s}\widehat{t}(2๐ค_{2T}^4+2(\widehat{s}+\widehat{t})๐ค_{2T}^2+\widehat{s}^2+\widehat{t}^2),$$
$`(A.6)`$
$$F_{gq}(๐ค_{1T}^2,๐ค_{2T}^2)=6m^8(2๐ค_{1T}^4+2(\widehat{s}+\widehat{u})๐ค_{1T}^2+3\widehat{s}^2+3\widehat{u}^2+14\widehat{s}\widehat{u})m^4+$$
$$(2(\widehat{s}+\widehat{u})๐ค_{1T}^4+8\widehat{s}\widehat{u}๐ค_{1T}^2+\widehat{s}^3+\widehat{u}^3+7\widehat{s}\widehat{u}^2+7\widehat{s}^2\widehat{u})m^2$$
$$\widehat{s}\widehat{u}(2๐ค_{1T}^4+2(\widehat{s}+\widehat{u})๐ค_{1T}^2+\widehat{s}^2+\widehat{u}^2),$$
$`(A.7)`$
$$\genfrac{}{}{0pt}{}{F_{qq}(๐ค_{1T}^2,๐ค_{2T}^2)=6m^8(3\widehat{t}^2+3\widehat{u}^2+14\widehat{t}\widehat{u})m^4+(\widehat{t}^3+\widehat{u}^3+}{7\widehat{t}\widehat{u}^2+7\widehat{t}^2\widehat{u})m^2\widehat{t}\widehat{u}(\widehat{t}^2+\widehat{u}^2).}$$
$`(A.8)`$
Since we take into account the $`๐ค_T`$ depencence of the incoming virtual gluon polarization tensor, the functions $`F_{qg}(๐ค_{1T}^2,๐ค_{2T}^2)`$ and $`F_{gq}(๐ค_{1T}^2,๐ค_{2T}^2)`$ also depend obviously on the gluon transverse momentum. It is clear that if we take the limit $`๐ค_{1T}^20`$, $`๐ค_{2T}^20`$ in (A.1) โ (A.8) we easily obtain the corresponding collinear formulas.
Finally, we would like to point out again that in numerical computations we use precise off-shell kinematics and therefore all expressions (A.1) โ (A.8) depends on the parton transverse momentum. In particular, the incident parton momentum fractions $`x_1`$ and $`x_2`$ in (13) and (15) have some $`๐ค_T`$ dependence. In the limit $`๐ค_{1T}0`$, $`๐ค_{2T}0`$ we reproduce standard leading-order QCD collinear results. |
warning/0506/astro-ph0506416.html | ar5iv | text | # An Analysis of the Shapes of Ultraviolet Extinction Curves. IV. Extinction without Standards (To appear in the September 2005 Astronomical Journal)
## 1 INTRODUCTION
A detailed determination of the wavelength dependence of interstellar extinction, i.e., the absorption and scattering of light by dust grains, is important for two very different reasons. First, since it is a product of the optical properties of dust grains, extinction provides critical diagnostic information about interstellar grain populations (including size distribution, grain structure, and composition), providing guidance for interstellar grain models. Second, the accuracy to which the intrinsic spectral energy distributions (SEDs) of most astronomical objects can be determined depends on how well the effects of extinction can be removed from observations. In both cases, a fundamental issue is how accurately the wavelength dependence of extinction can be measured. Consequently, it is essential to have a firm grasp of how measurement errors can affect the determination of this wavelength dependence.
This paper is the culmination of a series of โtechniquesโ papers published over the past five years (Fitzpatrick & Massa 1999; Massa & Fitzpatrick 2000; and Fitzpatrick & Massa 2005; hereafter FM99, MF00, and FM05, respectively) whose aim has been to develop a technique to simultaneously determine the wavelength dependence of extinction (to higher accuracy than previously possible) and the physical properties of a reddened star. It represents a continuation of our earlier series on the properties of UV extinction (Fitzpatrick & Massa 1986; 1988; and 1990, hereafter FM90), and provides a detailed description of a new technique for deriving interstellar extinction curves which does not rely on observations of standard stars, virtually eliminates the effects of โmismatchโ error, and yields an accurate assessment of the uncertainties. This โextinction-without-standardsโ technique opens the door to a new class of extinction studies, including regions heretofore inaccessible. For example, errors in the traditional โpair methodโ approach to extinction strongly limit our ability to study extinction in two important regimes. The first is low-$`E(BV)`$ sightlines, where one might hope to relate extinction properties to the environmental properties of specific physical regions. The second is extinction derived from mid- to late B stars. These stars are plentiful, and often constitute the bulk of the stars available for extinction measurements in intrinsically interesting regions, such as the Pleiades. We will demonstrate how the extinction-without-standards approach overcomes both of these problems and present examples of each. Some first results from this program were illustrated by Fitzpatrick (2004, hereafter F04). In addition to the extinction results, we will demonstrate that our approach simultaneously provides accurate stellar parameters which are also astrophysically interesting.
In ยง2, we provide a broad overview of the basic problem of measuring an extinction curve, and compare the merits of curves derived by the pair method and curves derived by using stellar atmosphere models. In ยง3, we describe our new model-based technique in detail, and list the basic ingredients needed to determine an extinction curve using this approach. In ยง4, we describe the data used in the current study. In ยง5 we provide a number of sample extinction curves derived using model atmospheres. We also demonstrate the high precision of the new curves and the reliability of the error analysis employed. Finally in ยง6 we describe some to the scientific advantages of this new technique and our plans to exploit them.
## 2 MEASURING EXTINCTION
To understand how interstellar extinction is measured and to appreciate how different measurement techniques can affect the outcome, we begin with the intrinsic elements of an uncalibrated observation of a spectral energy distribution (SED) of a reddened star obtained at the earth, $`f_\lambda `$. This can be expressed as
$$f_\lambda =F_\lambda r_\lambda \theta _R^2\mathrm{\hspace{0.17em}10}^{0.4A_\lambda }$$
(1)
where $`F_\lambda `$ is the intrinsic surface flux of the star at wavelength $`\lambda `$, $`r_\lambda `$ is the response function of the instrument, $`\theta _R(R/d)^2`$ is the angular radius of the star (where $`d`$ is the stellar distance and $`R`$ is the stellar radius), and $`A_\lambda `$ is the absolute attenuation of the stellar flux by intervening dust (i.e., the total extinction) at $`\lambda `$. Alternatively, the observed SED can be expressed in terms of magnitudes $`m_\lambda `$ by
$`m_\lambda `$ $`=`$ $`2.5\mathrm{log}F_\lambda \theta _R^2+A_\lambda +C_\lambda \mathrm{or}`$ (2)
$`m_\lambda `$ $`=`$ $`M_\lambda +5\mathrm{log}d5+A_\lambda +C_\lambda `$ (3)
where $`C_\lambda =2.5\mathrm{log}r_\lambda `$ is a term which transforms between the observed magnitude system and absolute flux units, and $`M_\lambda `$ is the traditional definition of the absolute magnitude of the star at $`\lambda `$.
The difficulty in measuring the total extinction $`A_\lambda `$ can be seen by rearranging these equations to solve for the extinction term. Equation 1 yields
$$A_\lambda =2.5\mathrm{log}(r_\lambda \theta _R^2\frac{F_\lambda }{f_\lambda }),$$
(4)
while Equation 3 yields
$$A_\lambda =m_\lambda M_\lambda 5\mathrm{log}d+5+C_\lambda .$$
(5)
In either case, a true measurement of $`A_\lambda `$ would require calibrated SED observations, knowledge of the intrinsic SED of the star, and measurements of both the stellar distance and radius, or their ratio $`\theta _R`$. Unfortunately, there are no early-type stars for which both of these latter quantities are known to sufficient accuracy to allow a meaningful measurement of $`A_\lambda `$. As a result, indirect methods must be employed, and $`A_\lambda `$ is always a derived quantity, subject to assumptions.
In virtually all extinction studies, the actual measured quantity is a โcolor excessโ which describes the extinction at a wavelength $`\lambda `$ relative to that at a fiducial wavelength. The traditional approach is to adopt the $`V`$ band as the fiducial, since $`V`$ magnitudes are widely available, accurately calibrated, and typically of high quality. This color excess can be expressed as either
$$E(\lambda V)A_\lambda A_V=2.5\mathrm{log}(\frac{r_\lambda }{r_V}\frac{F_\lambda }{F_V}\frac{f_V}{f_\lambda }),$$
(6)
as based on Equation 4, or
$`E\left(\lambda V\right)`$ $``$ $`A_\lambda A_V`$ (7)
$`=`$ $`\left(m_\lambda m_V\right)\left(M_\lambda M_V\right)+\left(C_\lambda C_V\right)`$
$`=`$ $`m\left(\lambda V\right)M\left(\lambda V\right)+C\left(\lambda V\right),`$
as based on equation 5. Thus, the determination of the color excess requires only the measurement of the observed SED and a knowledge of the shape of the intrinsic SED. There are two basic approaches to determining color excesses, based on the use of either unreddened stars or stellar atmosphere models to represent the intrinsic SEDs of reddened stars. These two techniques, and the issue of the normalization of extinction curves, are discussed in the three subsections to follow.
### 2.1 The Pair Method
The first approach is the โpair method.โ A pair method curve is constructed by comparing the fluxes of a reddened star and an (ideally) identical unreddened โstandard star.โ Essentially, the absolute magnitudes in Eq. 7 are replaced by the observed magnitudes of the standard star and the curve is usually expressed in the form
$$k\left(\lambda V\right)\frac{E\left(\lambda V\right)}{E\left(BV\right)}\frac{m\left(\lambda V\right)m\left(\lambda V\right)_0}{\left(BV\right)\left(BV\right)_0},$$
(8)
where the color excesses $`E(\lambda V)`$ are normalized by $`E`$(B V) and the subscripted quantities refer to unreddened indices for the standard star. (The issue of normalization will be discussed below.) When the two stars are observed using the same instrument, this method has the advantage that calibration terms cancel, eliminating any dependence on the absolute flux calibration, $`r_\lambda `$ or $`C_\lambda `$. There are, however, two disadvantages to this technique. The first is that the grid of unreddened standard stars is necessarily limited, so that some mismatch in the SEDs of the reddened and unreddened stars, termed โmismatch errorโ, is inevitable. The second is that there are very few truly unreddened early-type stars, so usually the โunreddenedโ standard must be corrected for some small amount of extinction whose exact magnitude and wavelength dependence are uncertain. This creates an error that can propagate into the resulting extinction curves.
Massa, Savage, & Fitzpatrick (1983) presented a detailed study of the uncertainties affecting pair method extinction curves and showed that mismatch effects dominate the error budget. If one uses a single unreddened spectral standard for each spectral class, and assumes that all spectral classifications are perfect, then as a result of spectral binning, mismatch errors will be equal to or less than half a spectral class. Figure 1 shows how such mismatches can affect extinction curves derived from main sequence B stars. In each of the four groups of curves in Figure 1, the true shape of an extinction curve affecting a B1, B2, B5, or B9 star is indicated by the solid curve. The extinction curves which would be derived via the pair method in the presence of a $`\pm \frac{1}{2}`$ spectral class mismatch error are shown by the dash-dot curves for the case of $`E(BV)=0.15`$ and by the dashed curves for $`E(BV)=0.30`$. (The derived curves fall below the true curve in the UV and above the true curve in the IR when the standard star is cooler than the reddened star, and vice versa.) In practice, spectral classifications are not perfect and the available unreddened standard stars do not necessarily lie in the middle of the range of properties within a single spectral class. Both these effects exacerbate the spectral mismatch problem and thus the uncertainties shown in Figure 1 are likely closer to typical errors, rather than extremes.
Although it may appear from Figure 1 that a discontinuity at the Balmer jump at $`\lambda ^12.7\mu \mathrm{m}^1`$ would provide an obvious indicator of the presence of spectral mismatch in an extinction curve, this is almost never practical. The spectrophotometric data available for constructing extinction curves are generally limited to UV wavelengths ($`\lambda ^1>3.3\mu \mathrm{m}^1`$), which are highlighted in Figure 1. Typically, the only data available in the optical and near-UV are photometric indices which straddle the Balmer jump, such as the Johnson $`UB`$ color, and these cannot be used to distinguish the effects of mismatch from intrinsic curve shape.
Mismatch error clearly can have a profound effect on the shapes of curves derived from stars with low color excesses, particularly in the mid-to-late B spectral range and most particularly in the UV spectral region. In fact, it is mismatch error which provides the low temperature cutoff to the spectral range of stars from which useful UV extinction measurements can be made. Mismatch also severely limits extinction studies based on stars hotter and/or more luminous than the main sequence B stars. For the O stars, very few unreddened standard stars exist and extreme mismatching of spectral types is often necessary to derive a curve โ although, since the intrinsic UV/optical SEDs of the O stars are not well known, it is not clear how large an effect this introduces in the derived curves. The use of luminous, evolved stars for extinction studies is particularly problematic since unreddened standards are rare and the sensitivity of the intrinsic SEDs to both temperature and surface gravity can lead to very severe mismatch effects in the resultant curves (although see the discussion in Cardelli, Sembach, & Mathis 1992 for results in the early-B spectral range).
### 2.2 Model Atmosphere Techniques
The second approach to deriving extinction curves is to model the intrinsic SED of the reddened star in order to isolate the effects of extinction. This technique was first used by Whiteoak (1966) to analyze optical spectrophotometry and has been applied in one form or another many times since. We refer to it as โextinction-without-standardsโ, since it does not rely upon a set of unreddened standard stars to determine the extinction curve. In this method, a model atmosphere of the reddened star is determined from its photometric or spectral properties. The advantage of this approach is that, in principle, a perfect, unreddened match can be determined for the intrinsic SED of the reddened star, eliminating the mismatch error which plagues pair method curves. The apparent disadvantages of the approach are not actual disadvantages, but rather requirements, which can limit its accuracy and range of applicability. The first requirement is that a set of models must exist whose accuracy can be quantified and validated. The second is that the observations must contain adequate โreddening freeโ information to accurately determine the intrinsic SED of a reddened star. The third is that the fluxes must be precisely calibrated, i.e., $`r_\lambda `$ must well determined.
We have been developing the necessary constituents of this method over the past five years. We began by demonstrating that the Kurucz (1991) ATLAS9 models provide faithful representations of the observed UV and optical SEDs of near main sequence B stars (FM99). We then showed that a combination of an observed UV SED and optical photometry provides adequate reddening independent information to determine both the appropriate ATLAS9 model for a reddened star and a set of parameters (defined by FM90) which determine the shape of its interstellar extinction curve. Subsequently, we refined the calibration of the IUE data, in order to improve the quality of the fits and the robustness of the physical information derived from them (MF00). Next, we verified the physical parameters derived from the models through applications of the models to eclipsing binary data, where the results must agree with other constraints (see, Fitzpatrick et al. 2003 and references therein). Finally, we have used Hipparcos data of unreddened B stars to derive a consistent recalibration of optical and NIR photometry (FM05) and, in the process, once again demonstrated the internal consistency of the models for near main sequence B stars. As a result of these efforts, we now have internally consistent $`r_\lambda `$ for IUE and optical and NIR photometry, and are in a position to apply the results and to quantify the associated errors. These are the objectives of the current paper, and are discussed fully in the next section.
### 2.3 Normalizing The Curves
Once color excesses have been determined, we are faced with the problem of how to compare excesses derived for lines of sight with different amounts of extinction. After all, we are interested in the โshapesโ of the curves, since these may reveal important clues about the size distribution and composition of the dust. This desire to compare shapes, brings us to the normalization problem.
From a purely mathematical point of view, a straightforward approach would be to normalize the curves by their norm,
$$e(\lambda V)=\frac{E(\lambda V)}{\sqrt{_\lambda E(\lambda V)^2}}$$
(9)
With appropriate weighting, this normalization could, on average, minimize the observational error in the normalization factor and reduce systematic effects. However, such a normalization would be sensitive to the strength of narrow features, such as the 2175ร
bump, and this could mask the overall agreement between curves over the majority of the wavelength range. A second approach is to search for some immutable feature of the curve and normalize by that. The idea behind this procedure is that if some aspect of all curves is always the same, then all curves can be normalized by the strength of this property and all of the resulting curves will be directly comparable. This is the motivation for the $`A_V`$ normalization adopted by Cardelli, Clayton, & Mathis (1989). They assume that all extinction curves have a very similar (although not identical) form for $`\lambda V`$. However, there are problems with this normalization as well. In particular, as pointed out above, $`A_V`$ is not a directly measured quantity, but must be derived from IR photometry and requires assumptions about the shape of extinction curves at very long wavelengths. The shape of this extinction is often considered to have a universal form, but this has been demonstrated for only a relatively small number of sightlines (Reike & Lebofsky 1985; Martin & Whittet 1990). In addition, measurements of extinction in the IR can be compromised because the stellar SEDs become increasingly faint at long wavelengths and other sources of light, such as circumstellar emission or scattering from dust in the near stellar environment, can contaminate the SED measurements. Furthermore, IR color excesses are usually small for stars that are detectable in the UV, so UV curves normalized by quantities derived from IR photometry may be affected by large normalization errors. Consequently, the absolute level of such curves can be poorly defined. Finally, even with the advent of the 2MASS data base, there are still many stars which do not have IR photometry available.
As a result of the complications mentioned above, and since our intent is to demonstrate how precisely curves can be measured while making a minimum number of assumptions, we have opted to use the conventional $`E(BV)`$ normalization, as shown in Eq. 8. While the interpretation of $`E`$(B V) as a measure of the โamountโ of interstellar dust is not straightforward, its widespread availability, observational precision, and lack of requisite assumptions make it the best choice for this study. Nevertheless, we note that it is a simple matter to transform from one normalization to another, and emphasize that $`E(\lambda V)`$ is actually the basic measurable quantity.
## 3 EXTINCTION-WITHOUT-STANDARDS
### 3.1 Formulation of the Problem
Our earlier studies (FM99, FM05) have shown that the observed SEDs, $`f_\lambda `$, of lightly- or unreddened main sequence B stars can be modeled very successfully by using a modified form of Eq. 1, namely,
$$f_\lambda =F_\lambda \theta _R^2\mathrm{\hspace{0.17em}10}^{0.4E(BV)[k(\lambda V)+R(V)]}.$$
(10)
The use of absolutely calibrated datasets eliminates the calibration term $`r_\lambda `$ and the total extinction $`A_\lambda `$ has been broken down into a normalized shape term ($`k(\lambda V)`$), a normalized zeropoint ($`R(V)A_V/`$E$`(\mathrm{BโV})`$), and a scale factor ($`E`$(B V)). Providing that the righthand side of the equation can be represented in a parameterized form, the equation can be treated as a non-linear least squares problem and the optimal values of the parameters โ which provide the best fit to the observed SED $`f_\lambda `$ โ can be derived along with error estimates. Because the stars under study were lightly reddened, we could replace the extinction curve $`k(\lambda V)`$ and the offset term $`R(V)`$ with average Galactic values without loss of accuracy. Using the Kurucz ATLAS9 stellar atmosphere models to represent $`F_\lambda `$, the results of the fitting procedure were estimates of 6 parameters: $`E`$(B V), $`\theta _R`$, and the four parameters that define the best-fitting model, i.e., $`T_{\mathrm{eff}}`$, $`\mathrm{log}g`$, the metallicity \[m/H\], and the microturbulence velocity $`v_t`$. We performed the fits using the MPFIT procedure developed by Craig Markwardt <sup>1</sup><sup>1</sup>1Markwardt IDL Library available at http://astrog.physics.wisc.edu/$``$craigm/idl/idl.html.
The fitting process described above begins to break down when the color excess $`E`$(B V) of the target stars exceed $``$0.05 mag. By this we mean that large systematic residuals begin to appear, which greatly exceed the measurement errors. The reason is simple: the wavelength dependence of interstellar extinction curves varies greatly from sightline-to-sightline, and once $`E(BV)0.05`$ mag the differences between the true shapes of the curves and the assumed mean form begin to exceed the measurement error. However, FM99 noted that the SEDs of significantly reddened stars could still be modeled using Equation (10) (and the non-linear least squares approach) if the wavelength dependence of the extinction curve could be represented in a flexible form whose shape could be adjusted parametrically to achieve a best fit to the observations, and if these parameters were determined simultaneously with the stellar parameters. This is the essence of our โextinction-without-standardsโ approach.
Successfully modeling the shape of reddened stellar SEDs requires four principal ingredients: 1) an observed SED that spans as large a wavelength range as possible, 2) an accurate absolute flux calibration ($`r_\lambda `$ or $`C_\lambda `$), 3) an extinction curve whose shape can be described by a manageable set of parameters, and 4) a grid of stellar surface fluxes, $`F_\lambda `$, whose defining parameters can be determined from the observational data. In ยง4 we will describe the particular datasets used in this paper to demonstrate our technique. We have already discussed how MF00 and FM05 have determined the necessary calibrations. In the remainder of this section, we describe our flexible form for the interstellar extinction curve (ยง3.2) and the grid of stellar surface fluxes with which we have developed our approach (ยง3.3).
### 3.2 A Flexible Representation of the Interstellar Extinction Curve
We adopt a flexible and adjustable form for the UV-through-IR extinction curve, whose shape can be optimized to fit the SED of a reddened star through the adjustment of a specific set of parameters in the least-squares minimization procedure. This curve is illustrated in Figure 2. It consists of two main regions: 1) the UV ($`\lambda <2700`$ ร
; solid curve) where the parameterized form of FM90 is adopted and 2) the near-UV/optical/IR ($`\lambda >2700`$ ร
; dashed line) where we use a cubic spline interpolation through a set of UV (U<sub>1</sub>, U<sub>2</sub>), optical (O<sub>1</sub>, O<sub>2</sub>, O<sub>3</sub>, O<sub>4</sub>), and IR (I<sub>1</sub>, I<sub>2</sub>, I<sub>3</sub>, I<sub>4</sub>) anchor points to represent the curve. The interpolation is performed using the Interactive Data Language (IDL) procedure SPLINE. We adopt a spline representation for the near-UV/optical/IR curve simply because we do not have reliable, detailed information on the wavelength dependence of the extinction in the near-IR through near-UV region (1$`\mu `$m โ 3000 ร
). It is ironic that the portion of the curve that is accessible from the ground is more poorly characterized than the portion accessible only from space. As a result, we do not know whether the optical to near-IR region of the curve can be represented by a compact analytical formula. Our hope is that, by applying our procedure to a large sample of sightlines, we will ultimately be able to characterize the shape of the extinction law in this region by simple relations and determine whether sightline-to-sightline variations are correlated with other aspects of the curve or with interstellar environment. The placement of the spline points resulted from considerable experimentation, but certainly cannot be represented as an objectively determined optimal result. The current arrangement does, however, allow us to model the major available datasets to a level consistent with the observational errors.
The FM90 parameterization scheme contains 6 free parameters to represent 3 functionally separate features that are summed to produce the UV curve. An underlying linear component, indicated by the dotted line in Figure 2, is specified by an intercept $`c_1`$ and a slope $`c_2`$. The Lorentzian-like 2175 ร
bump is fit by a Drude profile $`D(x,x_0,\gamma )`$, where $`x_0`$ and $`\gamma `$ specify the position and FWHM of the bump, respectively, whose strength is determined by a scale factor $`c_3`$. Finally, the degree of departure of the curve in the far-UV from the underlying linear component is specified by a single parameter $`c_4`$. Defining $`x\lambda ^1`$, the complete UV function is given by:
$$k(\lambda V)=c_1+c_2x+c_3D(x,x_0,\gamma )+c_4F(x),$$
(11)
where
$$D(x,x_0,\gamma )=\frac{x^2}{(x^2x_0^2)^2+x^2\gamma ^2},$$
(12)
and
$`F\left(x;x>5.9\right)`$ $`=`$ $`0.5392\left(x5.9\right)^2+0.05644\left(x5.9\right)^3,`$ (13)
$`F\left(x;x5.9\right)`$ $`=`$ $`0.`$ (14)
It is worth emphasizing that the FM90 parameterization is a mathematical scheme only, which allows us to reproduce UV extinction curves in a shorthand form (and makes the current program possible). However, it is not to be assumed that the functional components of the scheme represent actual separate extinction components arising in distinct dust grain populations. This parameterization has proven very useful and to the best of our knowledge is able to reproduce all known UV extinction curves to the level of observational error. Its flexibility will be illustrated in ยง5 below. Note that we terminate the FM90 formula at 2700 ร
, although we originally utilized UV data extending to 3000 ร
. Additional experience with the UV data has suggested that real extinction curves begin to exhibit departures from the FM90 formula near 2700 ร
, due to the unrealistic extrapolation of the linear component into the blue-violet region.
The 10 spline anchor points which characterize the near-UV/optical/IR portion of the curve are determined by a least squares fit to the IUE data longward of 2700 ร
and the available optical and IR photometry. Although there are 10 anchor points, there are also 6 constraints, so the fit actually introduces only 4 additional degrees of freedom. We now discuss these constraints in detail.
The two UV anchor points, U<sub>1</sub> and U<sub>2</sub> at 2700 and 2600 ร
, respectively, are fixed at the values of the FM90 UV fitting function at their wavelengths and are not adjustable. Together with O<sub>1</sub> at 3300ร
, these points guarantee that the curve which passes through the IUE data between 2700 and 3000 ร
will join both the UV and optical portions of the curve smoothly.
The four optical anchor points, O<sub>1</sub>, O<sub>2</sub>, O<sub>3</sub>, and O<sub>4</sub> at 3300, 4000, 5530, and 7000 ร
, respectively, are fit under two constraints: that the interpolated curve produces a value of $`k(\lambda V)=0`$ in the V band, and that the curve be normalized to unity in $`E`$(B V). Thus, only two free parameters emerge from this region.
The four IR points, I<sub>1</sub>, I<sub>2</sub>, I<sub>3</sub>, and I<sub>4</sub> are located at 0.25, 0.50, 0.75, and 1.0 $`\mu `$m<sup>-1</sup>, respectively. These four points are constrained to satisfy the formula
$$I_nk(\lambda V)=k_{IR}\lambda _n^{1.84}R(V)$$
(15)
where the scale factor, $`k_{IR}`$, and the intercept, $`R(V)`$, are the only free parameters. This is the power-law form usually attributed to IR extinction, with a value for its exponent from Martin & Whittet (1990). The exponent of the power-law could, potentially, be included as a free parameter in the fitting procedure, and we will investigate this in the future. However, our impression is that the IR data available to us (primarily 2MASS $`JHK`$ photometry) are insufficient to determine this quantity accurately. All results presented in this paper assume an exponent of -1.84 in Equation 15.
### 3.3 The Stellar Surface Fluxes
To represent the intrinsic surface fluxes, $`F_\lambda `$, of reddened stars we utilize R.L. Kuruczโs line-blanketed, hydrostatic, LTE, plane-parallel ATLAS9 models, computed in units of erg cm<sup>-2</sup> sec<sup>-1</sup> ร
<sup>-1</sup> and the synthetic photometry derived from the models by Fitzpatrick & Massa (2005). These models are functions of four parameters: $`T_{\mathrm{eff}}`$, $`\mathrm{log}g`$, \[m/H\], and $`v_t`$. All of these parameters can be determined in the fitting process although, because of data quality, it is sometimes necessary to constrain one or more to a reasonable value and solve for the others.
The general technique of deriving extinction curves via stellar atmospheres is, of course, not dependent on the specific set of models used. In the present case, the most important consideration for our adoption of the ATLAS9 models is that they work โ at least within a specific spectral domain. FM99 and FM05 have shown that these models provide excellent fits to the observed SEDs for lightly- or un-reddened main sequence B stars throughout the UV through near-IR spectral region. In addition, experience with eclipsing binaries (see Fitzpatrick et al. 2003 and references within) has shown that the good SED fits also yield accurate estimates of the physical properties of the stars. Because of the physical ingredients of the models (specifically LTE and plane-parallelism) we currently restrict our attention to the main sequence B stars. We plan to investigate how well these models reproduce the SEDs and properties of somewhat more luminous B stars and also the later O types. Also, we will take advantage of more complex models (e.g., the non-LTE TLUSTY models) as the available grids expand their parameter ranges.
### 3.4 Summary
To summarize, we model the observed SEDs of reddened near main sequence B stars by treating equation (10) as a non-linear least squares minimization problem. As a result we can simultaneously obtain estimates of the physical properties of a reddened star and the shape of the interstellar extinction curve distorting the starโs SED. A total of 16 parameters specify the righthand side of the equation, including $`\theta _R`$, $`E(BV)`$, four to define $`F_\lambda `$, and ten to define the shape of the extinction curve. Depending on limitations of the available data, and known properties of the stars or interstellar medium, any subset of the parameters can be constrained to predetermined values.
## 4 THE DATA
In the following section, we will illustrate the potential of our extinction-without-standards technique, utilizing a set of 27 lightly-to-heavily reddened stars. For this demonstration, and indeed for extinction determinations in general, the ideal SED dataset would consist of absolutely-calibrated spectrophotometry spanning the UV-through-IR spectral regions. Such data would allow a straightforward comparison between observations and stellar atmosphere models (since both are presented in simple flux units) and would provide the most detailed view of the wavelength dependence of interstellar extinction. While a small amount of such data is available (see, e.g., Fitzpatrick at al. 2003), the largest existing database of absolutely-calibrated spectrophotometry is the low-resolution archive of the International Ultraviolet Explorer satellite (IUE) which covers the UV region only (1150โ3000 ร
). In the optical and near-IR, the largest SED databases are photometric in nature, consisting of Johnson, Strรถmgren, and Geneva photometry in the optical and 2MASS JHK photometry in the near IR. Utilizing these resources, we can examine the UV region for small scale features but can only study the broad scale wavelength-dependence of extinction in the optical and near-IR regions.
We use NEWSIPS IUE data (Nichols & Linsky 1996) obtained from the MAST archive at STScI. These data were corrected for residual systematic errors and placed onto the HST/FOS flux scale of Bohlin (1996) using the corrections and algorithms described by MF00. This step is absolutely essential for our program since our โcomparison starsโ are stellar atmosphere models and systematic errors in the absolute calibration of the data do not cancel out as in the case of the pair method. (The NEWSIPS database is also contaminated by thermally- and temporally-dependent errors, which would not generally cancel out in the pair method โ see MF00.) Multiple spectra from each of IUEโs wavelength ranges (SWP or LWR and LWP) were combined using the NEWSIPS error arrays as weights. Small aperture data were scaled to the large aperture data and both trailed and point source data were included. Short and long wavelength data were joined at 1978 ร
to form a complete spectrum covering the wavelength range $`1150\lambda 3000`$ ร
. Data longward of 3000 ร
were ignored because they are typically of low quality and subject to residual systematic effects. The IUE data were resampled to match the wavelength binning of the ATLAS9 model atmosphere calculations in the wavelength regions of interest.
Mean values of the Johnson UBV, Strรถmgren uvby$`\beta `$, and Geneva UBB<sub>1</sub>B<sub>2</sub>VV<sub>1</sub>G photometric magnitudes, colors, and indices for the program stars were acquired from the Mermilliod et al. (1997) archive. 2MASS JHK magnitudes for all stars, along with their associated errors, were obtained from the 2MASS All-Sky Point Source Catalog at the NASA/IPAC Infrared Science Archive. Johnson $`VR`$, $`RI`$, and JHK data are also available for some of the stars, and were obtained from the Mermilliod et al. archive.
## 5 SOME INITIAL RESULTS
In this section, we demonstrate the potential of our extinction-without-standards technique, utilizing a set of the 27 reddened stars, listed in Table 1, which fall into three groups: 1) stars in the open cluster IC 4665; 2) stars with moderate-to-heavy reddening; and 3) lightly reddened stars in a specific region of the sky. These representative examples illustrate the advantages of our approach and highlight several scientific applications which will be pursed in future studies, using expanded samples of stars. In addition, they provide confirmation of the error analysis incorporated in our approach.
For this demonstration sample, the SED fitting procedure was applied as described above, with the following additional details:
* The SED data modeled in the fitting procedure include IUE UV spectrophotometry, the Johnson $`V`$, $`BV`$, and $`UB`$ indices, the Strรถmgren $`by`$, $`m_1`$, $`c_1`$, and $`\beta `$ indices, the Geneva $`UB`$, $`VB`$, $`B_1B`$, $`B_2B`$, $`V_1B`$, and $`GB`$ indices, and the 2MASS JHK magnitudes. Johnson $`VR`$, $`RI`$, and JHK data are also available for a few of the stars.
* The optical extinction spline point O<sub>4</sub> at 7000 ร
is only well-determined when optical $`R`$ and $`I`$ band photometry are available. Therefore in the examples below, we only solve for O<sub>4</sub> in such cases. For the other stars, the optical portion of the extinction curve is determined only by the spline points O<sub>1</sub>, O<sub>2</sub>, and O<sub>3</sub>.
* During the $`\chi ^2`$ minimization, a reddened and distance-attenuated model was created from the current set of input parameters, and then synthetic photometry was performed on this model to produce the photometric indices, which were then compared with observations. The synthetic photometry was calibrated as described by FM05, with the calibration extended to redder colors by us. The UV model fluxes and recalibrated IUE fluxes, both in units of erg cm<sup>-2</sup> sec<sup>-1</sup> ร
<sup>-1</sup>, were compared directly.
* The initial weighting factors for the various datasets in the $`\chi ^2`$ minimization were determined from their observational uncertainties, i.e., $`weight\sigma ^2`$. We then scaled the weights of the optical/near-IR photometry so that their total weight was equal to that of the IUE UV spectrophotometry, thus balancing the fit between the two datasets. This is the procedure adopted by FM05, except that we include the Strรถmgren $`\beta `$ index along with the rest of the optical/near-IR photometry, rather than assigning it its own (high) weight. FM05 weighted $`\beta `$ heavily in recognition of its value as a surface gravity indicator. However, we have found that the temperature-sensitivity of $`\beta `$ (particularly in the later B stars) combined with observational errors, can lead to very unsatisfactory fits when $`\beta `$ is over-emphasized. Treating $`\beta `$ in the same manner as the rest of the photometric indices seems to be the simplest and most reasonable approach.
* Because interstellar H i Lyman $`\alpha `$ absorption in reddened stars can have a significant impact on the starโs apparent continuum level far from line center at 1215 ร
, we convolve the profile of this heavily damped line with the model atmosphere SEDs before comparing with observations. Along heavily reddened sightlines, where the H i column density $`N(`$i$`)`$ is high and the signature of the atomic absorption strong, the value of $`N(`$i$`)`$ can actually be incorporated into the fitting procedure as a free parameter to be optimized. We will utilize this capability in future studies. For the present, mostly lightly-reddened, sample, however, the $`N(`$i$`)`$ values used to construct the line profiles were taken from the survey of Diplas & Savage (1994) or else computed from the general relation $`N(`$i$`)=4.8\times 10^{21}E(BV)\mathrm{cm}^2`$ from Bohlin, Savage, & Drake (1978). The inclusion of the Lyman $`\alpha `$ line insures that we distinguish the effects of dust extinction from atomic absorption in the far UV region.
* The uncertainties in the best-fit parameters were determined by running 50 Monte Carlo simulations for each star, during which the input data were randomly varied assuming a Gaussian distribution of observational uncertainties and a new fit performed. The zero points and random photometric uncertainties of the short-wavelength and long-wavelength IUE fluxes were varied as described in FM04; the assumed observational errors in the Johnson, Strรถmgren, and Geneva indices were as given in Table 7 of FM04; and the uncertainties in the 2MASS data were as obtained from the 2MASS archive. In addition, the $`V`$ magnitude was assumed to have a 1-$`\sigma `$ uncertainty of 0.015 mag. The adopted 1-$`\sigma `$ uncertainties for each parameter were taken as the standard deviation of the values produced by the 50 simulation.
### 5.1 The Open Cluster IC 4665
We begin our examples by examining the extinction towards an open cluster. While multiple scientific rewards can result from the study of extinction towards cluster stars (see the discussion in ยง6), our primary interest in IC 4665 is to demonstrate the โtechnicalโ advantages of our approach. Namely, the use of cluster extinction curves to help evaluate the magnitude of the uncertainties in the measurement of extinction curves, as discussed in detail by Massa & Fitzpatrick (1986). In particular, and because of its low $`E(BV)`$ and late-B stellar population, IC 4665 extinction curves provide an especially sensitive test of the precision and range of our extinction-without-standards approach.
The wavelength dependence of extinction towards IC 4655 was first examined by Hackwell, Hecht, & Tapia (1991; hereafter HHT) for the purpose of studying the relationship between extinction and IR emission, as measured by the Infrared Astronomical Satellite (IRAS). This remains one of the most challenging extinction studies yet performed for two reasons: 1) the mean reddening in the cluster is very low, $`E(BV)<0.2`$ mag, and 2) the spectral types of the available targets run from mid- to late B. Both facts exacerbate errors in the standard pair method approach, as has been shown in Figure 1. HHT recognized these uncertainties and ultimately concluded that the wavelength dependence of extinction among the cluster stars is uniform to within their ability to measure it.
Figure 3 shows the results of the SED fits for the nine IC 4665 stars considered here. The SEDโs of the best-fitting, reddened models are shown by the histogram-style curves. In the UV, the binned IUE fluxes are shown by the small circles. In the optical region, Johnson UBVRI magnitudes (converted to flux) are indicated by circles, Strรถmgren uvby magnitudes by triangles, and Geneva UB<sub>1</sub>B<sub>2</sub>VV<sub>1</sub>G magnitudes by diamonds. In the near-IR, 2MASS and Johnson JHK magnitudes are shown by the large filled and open circles, respectively. Note that the photometric data have been converted to flux units for display purposes only. The comparison between models and observations was performed in the native photometric format (i.e., in magnitudes or colors as noted above).
Figure 4 shows our extinction-without-standards curves for the IC 4665 stars. The symbols show the actual normalized ratios between the models and the stellar SEDs, while the thick solid curves show the flexible UV-through-IR extinction curves whose shapes were determined by the fitting procedure. The curves have been arbitrarily shifted vertically for clarity, but all are shown compared with a similarly-shifted estimate of the average Galactic extinction curve for reference (thick dash-dot curves corresponding to $`R(V)=3.1`$, from Fitzpatrick 1999). As will be discussed further below, we assumed a value of $`R(V)=3.1`$ towards the cluster and did not include the IR JHK data in the fitting procedure. Thus, only the average Galactic curve is shown for wavelengths longward of 6000 ร
.
The various parameters determined by the fits are listed in Tables 1 and 2. The flexible extinction curves themselves, in the form $`E(\lambda V)/E(BV)`$ can be reconstructed from the parameters given in Table 2. The 1-$`\sigma `$ uncertainties of the extinction curves are indicated in Figure 4 by the grey shaded regions. The regions are centered on the means of the 50 Monte Carlo simulations with which we performed our error analysis and their thickness shows the standard deviation of the individual simulations.
Figure 5 compares our new curves in the UV with those derived by HHT using the standard pair method. HHTโs curves were reconstructed from the data in their Table 5. Note that the HHT study originally included 17 stars. We have eliminated 4 A-type stars, 3 chemically peculiar B-type stars, and 1 B-type shell system from consideration here since their extinction curves are particularly uncertain. Thus, Figures 4 and 5 show only the best-determined curves in the cluster, from the point of view of both the pair method and our technique. The remarkable aspect of Figure 5 is not the scatter among HHTโs curves โ it is exactly what should be expected given the limitations of the pair method โ but rather the impressive improvement in precision gained by the extinction-without-standards technique, as evidenced by the decrease in the curve-to-curve scatter. Clearly, this higher level of precision allows much firmer conclusions to be drawn about the degree of intrinsic variation of extinction across the face of the cluster, and also affords an improved potential to detect correlated behavior between the extinction properties and other aspects of the interstellar environment.
While we will present a full scientific analysis of the IC 4665 results in a future paper, the issue of the intrinsic variation of extinction among the cluster stars is important to consider here, since it can shed light on the accuracy of our error analysis. The top curve in Figure 4 shows the mean of the nine individual IC 4665 curves. The error bars show the actual sample standard deviation and the average Galactic curve is again presented for comparison. The critical result is that, over most of its wavelength range, the standard deviation of the sample about the mean curve is comparable to the predicted uncertainties in the individual curves, as shown by the shaded regions. This indicates that the (small) curve-to-curve variations seen are at a level consistent with our expected uncertainties and that the intrinsic level of variation among the cluster stars must be very low.
Another way to approach the issue of variability is to look at the scatter among the various parameters which define the extinction curves. These are shown in Table 3, where columns 2 and 3 list the weighted mean values of the parameters and their observed standard deviations, respectively. The predicted uncertainties, i.e., the RMS of the Monto Carlo-based errors for the individual stars, are listed in column 4. If our error analysis is reasonable, then the observed scatter should be the quadratic sum of the expected errors plus any intrinsic variability. The value of examining cluster extinction curves is that one might reasonably suppose (at least as a starting point) that the individual curves, derived for nearly coincident lines of sight, are actually independent measurements of a single โcluster curve,โ i.e., no intrinsic scatter. In such a case, the measured scatter actually reflects the real measurement errors and provides an important test of the error analysis. Examination of Table 3 shows that the predicted and observed scatters are indeed generally very similar, supporting the position that any intrinsic variations among the IC 4665 curves are close to the level of our ability to measure them and that we have accurately assessed the uncertainties in our results. The final column of Table 3 shows the implied values of the possible small intrinsic variations.
Several of the individual parameters merit some additional comment. Both the 2175 ร
bump FWHM ($`\gamma `$) and its strength ($`c_3`$ or $`A_{bump}`$) show evidence for some small level of variability within the sample, above our expected measurement errors. However, this is not clearcut because these measurements involve the region of the IUE spectra which typically has the lowest quality data and it is possible that weak systematic effects in the data themselves โ which are not accounted for in the error analysis โ could produce the small level of variability seen. One such systematic is a โreciprocity effectโ in long wavelength data which is not corrected by the MF00 algorithms (see Figures 12 and 13 of MF00). We conclude conservatively that there is marginal evidence for bump variations within the IC 4665 sample, but will ultimately rely on studies of several open clusters to determine whether our error analysis faithfully predict the real measurement errors in the bump region.
The case of the far-UV curvature (parameter $`c_4`$), for which the implied intrinsic scatter is three times greater than our measurement errors, is more interesting. It is clear from Figures 4 and 5 and Table 3 that most of this apparent variability arises from two sightlines, namely, those towards HD 161165 and HD 161184. In fact, the observed scatter in $`c_4`$ towards the other seven sightlines is essentially identical to the expected measurement error. Because HD 161165 and HD 161184 are the two coolest stars in the sample, and the only stars with effective temperatures less than 12,000 K, we suspected that the high $`c_4`$ values in their extinction curves might be artifacts of the analysis, resulting from a failure of the models to accurately portray the far-UV SEDs of these stars. The investigation presented below suggests that this could well be the case, with the extinction curve anomalies possibly arising from a difference between the chemical composition profiles of the stars and the atmosphere models.
The data in Table 1 show a very uniform โmetallicityโ \[m/H\] for the cluster, with a weighted mean of -0.50 and a standard deviation of 0.04. This is close to, and actually slightly smaller than, our estimate of the measurement errors, simultaneously confirming the accuracy of our error analysis and imposing a small upper limit on the intrinsic compositional variability within the cluster. Although our \[m/H\] values are simple scale factors which apply to a template set of ATLAS9 solar abundances, FM99 showed that the \[m/H\] derived in our analysis is most sensitive to the abundance of Fe โ due to a very strong opacity signature in the mid-UV โ and is most analogous to \[Fe/H\]. Since the Fe abundance in the ATLAS9 models is $``$0.2 dex larger than the currently accepted solar value of 7.45 (where H = 12.00; Asplund, Grevesse, & Sauval 2005), our results suggest that the IC 4665 stars are deficient by about a factor of two in Fe as compared to the Sun. The small scatter observed in \[m/H\] is both satisfying and expected, since the surface composition of Fe is not subject to evolutionary modification in these young stars, which presumably all formed from the same parent material.
We experimented with forcing \[m/H\] to a more solar-like level for the IC 4665 stars and immediately found two effects: 1) the $`\chi ^2`$ values all increased significantly because the Fe features in the mid-UV could not be fit as well, and 2) most of the extinction curves remained unchanged, but the anomaly in the far-UV region for HD 161165 and HD 161184 decreased dramatically. The first result was expected. The second was a surprise, but is understandable. If the cluster stars (or at minimum, the two coolest stars HD 161165 and HD 161184) have a non-solar ratio of Fe to the light metals, e.g., \[C/Fe\] $`>`$ 0, then our best fit models โ biased towards the Fe abundance โ would underestimate the opacity due to the light metals. For the cooler stars, such elements provide significant opacity in the far-UV region and the ATLAS9 models would not be able to account for such a specific opacity difference. The fitting procedure would respond by finding a higher far-UV extinction curve. The curves for the hotter stars are less affected by changing \[m/H\] since the light metal opacity is less significant. We tentatively conclude that the chemical composition of the B stars in IC 4665 may deviate strongly from that of the Sun, with subsolar Fe but a more โnormalโ level of the light metals such as C. This suggestion is easily tested with a fine analysis of high resolution stellar spectra and we will pursue this in the future. On the positive side, this result suggests that the UV continua of late B and early A stars might be used to determine both a scaled \[m/H\] and a mean \[CNO/Fe\] index, assuming a grid of models parameterized by both these composition indices is available. On the more sober side, it is a reminder that our technique is susceptible to stellar abnormalities. As with the pair method, all available data should be consulted to determine whether a particular star is suitable for deriving an extinction curve.
Our final comment on the IC 4665 results concerns the behavior of the IR data. In Figure 4 we show two sets of JHK photometry for each star and for the average cluster curve. The solid symbols indicate 2MASS measurements and the open symbols show HHTโs data. The two sets of data systematically differ, with the 2MASS results suggesting a value of $`R(V)`$ somewhat smaller then the average Galactic value of 3.1 and the HHT data suggesting a slightly larger value (HHT derive a mean of $`R(V)=3.25`$). This latter result is more consistent with expectations, given that the mean far-UV curve is lower than the Galactic average (see, e.g., Cardelli et al. 1989), but this is insufficient evidence for rejecting the 2MASS data. For now, our solution for this quandary has been to ignore both sets of data in fitting the IC 4665 SEDs and adopt a default value of $`R(V)=3.1`$. However, the discrepancy for measurements in this very complex region bears further investigation, as does that fact that, in both datasets, the mean curves actually show more extinction in the 2.2 $`\mu `$m $`K`$ band than in the 1.65 $`\mu `$m $`H`$ band.
The discussion above leads us to three primary conclusions: 1) The extinction towards IC 4665 shows at most only a small degree of spatial variability โ comparable with our measurement errors โ and that the cluster extinction curve would best be represented by averaging the results for the seven hottest cluster stars studied here; 2) we are able to determine accurate \[m/H\] values from the observations; and 3) our analysis yields reliable estimates of the (small) uncertainties in the extinction-without-standards curves and parameters, barring the presence of unusual systematic anomalies in the data sample. The first and second conclusions are scientific issues which we will pursue further in the future. The third provides a formidable demonstration of the superiority of the extinction-without-standards technique over the classical pair method for deriving extinction curves in both the precision of the results and the quantification of the uncertainties.
### 5.2 Moderately-Reddened Stars
With our error estimates verified, we now examine how curves derived by the current approach compare to pair method curves. Figure 6 shows the UV-through-IR fits to the SEDs of a set of 9 moderately-reddened early-B stars, most of which are well known for their extinction properties, and Figure 7 shows the corresponding extinction curves. These stars were selected because they illustrate the wide range that exists in the wavelength dependence of both UV and IR extinction. For these stars the IR data allow us to determine the values of $`R(V)`$ and so the flexible extinction curve fits are shown throughout the IR-to-UV domain. As for the IC 4665 stars, all the parameters describing the fits are given in Tables 1 and 2.
Also shown in Figure 7 are UV curves based on the pair method technique. The curve for HD 294264 is from Valencic, Clayton, & Gordon (2004); the curves for HD 210121 and HD 27778 are from unpublished measurements by us; and the others are from the catalog of FM90. The agreement between the model-based and pair method curves is reasonable, and much better than seen for IC 4665. This is consistent with spectral mismatch as the prime cause of the existing discrepancies, since the nine stars have both higher $`E(BV)`$ values and earlier spectral types than the IC 4665 stars, both of which tend to reduce the influence of mismatch errors. Note also that the best agreement between the pair method and the model-based curves occurs for the star HD 204827, for which identical results are found. Again, this is as it should be, since its $`E(BV)`$ is the largest of any star in the sample, minimizing the impact mismatch effects. The 1-$`\sigma `$ uncertainties for our flexible extinction curve fits are indicated by shaded regions around the curves, as in Figure 4.
The curves shown in Figure 7 demonstrate the ability of the parameterized, flexible extinction curve (see ยง3.2) to conform itself to the wide range of extinction curve shapes encountered in interstellar space.
### 5.3 Lightly-Reddened Stars
Figure 8 shows a set of nine SED fits for low color excess stars (0.10 $`E(BV)`$ 0.21) located along sightlines bounded by the Galactic coordinates $`347\mathrm{ยฐ}<l<355\mathrm{ยฐ}`$ and $`18\mathrm{ยฐ}<b<26\mathrm{ยฐ}`$. Figure 9 shows the corresponding extinction curves. These stars are all mid-to-late B members of the Upper Scorpius complex (Garrison 1967) and we encountered them while testing a procedure for scanning the IUE archives and automatically generating model-based extinction curves for stars in the appropriate spectral range. When examining the results for this pilot program which sampled high-latitude stars, it became obvious that curves derived from stars in this specific region demonstrated similar curve morphology which is distinctly different from the Galactic average. Given the low reddening and preponderance of late B stars, this strong systematic behavior would be missed by a pair method survey, with its large inherent mismatch errors.
Strong regional signatures are important in extinction studies because they may highlight the effects of specific physical processes on dust grain populations. Although we will examine this specific region in the future, two points are worth mentioning. First, it is important that similar curves result from stars with spectral types ranging from B2.5 to B9, verifying that the curves are not the result of some temperature-dependent systematic effect in fitting process. It is also worth noting that the region is near the $`\rho `$ Oph dark cloud. The star $`\rho `$ Oph A (HD 147933, whose extinction curve is shown in Figure 7) is located at a comparable distance and just south of the region, at Galactic coordinates of ($`l`$, $`b`$) = ($`353.7\mathrm{ยฐ},17.7\mathrm{ยฐ}`$). The mean curve for the nine lightly-reddened stars is shown at the top of Figure 9, along with the average Galactic curve (dash-dot curve) and the HD 1479433 curve (dotted curve) for comparison. In the UV, the regional curve is seen to be almost identical to that for the more heavily-reddened HD 147933 ($`E(BV)0.5`$), suggesting that the sightlines sample similar dust populations. Interestingly, however, the curves are not identical in the IR. The mean 2MASS JHK data for the nine star sample implies a value of $`R(V)3.4`$, slightly larger than the Galactic mean of 3.1, while HD 147933 has a value of $`R(V)4.3`$. This contrast between UV and IR behavior may indicate different physical processes at work along the higher density sightline towards HD 147933, or perhaps different timescales in the response of UV and IR extinction to modifications of dust grain properties.
In producing the extinction curves for these lightly reddened stars, we found that, in some cases, the value of $`k_{IR}`$ was very poorly determined and produced (presumably) spurious โbumps and wigglesโ in the IR portion of the curves. To eliminate this effect, we imposed a constraint on $`k_{IR}`$ for the whole sample, namely $`k_{IR}=0.63R(V)0.84`$. This is taken from F04, who found a very strong relationship between $`R(V)`$ and $`k_{IR}`$ from a larger sample of more heavily reddened stars (see Figure 6 of F04). The ability to impose scientifically reasonable constraints on the fitting procedure is a major advantage of the extinction-without-standards approach, and can potentially allow high quality extinction results to be derived from stars with even lower reddenings than those shown in Figure 9.
## 6 DISCUSSION
In the previous sections we first provided an overview of the process used to determine an extinction curve, clarifying the measurement process. We then presented a new method for deriving UV-through-IR extinction curves, based on the use of stellar atmosphere models to provide estimates of the intrinsic (unreddened) stellar SEDs rather than unreddened (or lightly reddened) โstandardโ stars. We have shown that this โextinction-without-standardsโ technique greatly increases the accuracy of the derived extinction curves, particularly in the cases of low reddening and cool spectral types (i.e., late-B), and allows a realistic estimation of the uncertainties. A side benefit of the technique is the simultaneous determination of fundamental properties of the reddened stars themselves ($`T_{\mathrm{eff}}`$, $`\mathrm{log}g`$, \[m/H\], and $`v_{turb}`$), making the procedure valuable for both stellar and interstellar studies. Given the physical limitations of the ATLAS9 models we currently employ, the technique is limited to near main sequence B stars. However, in principal, the procedure can be adapted to any class of star for which accurate model SEDs are available and for which the signature of interstellar reddening can be distinguished from those of the stellar parameters. Although we developed the procedure based on IUE spectrophotometry in the UV, and photometry in the optical and near-IR (requiring a calibration of synthetic optical and near-IR photometry), the ideal application of the technique would be with spectrophotometric data throughout the UV-through-IR domain, allowing the most detailed examination of the wavelength dependence of the extinction curves.
The specific scientific advantages afforded by the extinction-without-standards technique can be summarized as follows:
Increased Precision: The increased precision in the derived extinction curves allows us to improve our understanding of extinction in a number of ways. First and most simply, we will determine the basic UV-to-IR wavelength dependence of extinction along a wide variety of sightlines more precisely than has been possible in the past. Second โ and given that we know that extinction curve morphology varies widely from sightline-to-sightline (see Figure 7!) โ we will be able to study the form of the variability and search for relationships between various features and wavelength domains using a data set with small and well-defined uncertainties. Such relationships, e.g., the correlation between $`R(V)`$ and the steepness of UV extinction discovered by Cardelli et al. (1989), provide important constraints on the dust grain population causing the extinction. Non-correlations can be equally important. For example, the lack of a correlation between the position and width of the 2175 ร
bump demonstrated by Fitzpatrick & Massa (1986) remains a strong constraint on models for the bump carrier (Draine 2003). In either case, a precise knowledge of the measurement errors is required for transforming an observation into a scientific constraint. Finally, we will be able to place much stronger and more well-defined limits on the relationship between extinction curve morphology and interstellar environment. Curves derived from lines of sight to specific, localized regions, such as towards open clusters, can be particularly useful for relating curve properties to physical processes occurring in the interstellar medium (see, e.g., the study of Cepheus OB3 by Massa & Savage 1984 and Trumpler 37 by Clayton & Fitzpatrick 1987).
Access to lightly-reddened sightlines: The ability to accurately probe low $`E(BV)`$ sightlines (as exemplified in Figures 4 and 9) opens the door to studies of dust in regions that have not been thoroughly explored. These include halo dust, dust in very low density regions, and local dust. Halo dust is especially important since we must contend with its effects every time we look out of the Galaxy. There are indications (Kiszkurno-Koziej & Lequeux 1987) that its properties differ systematically from dust at lower Galactic altitudes, and this result needs to be verified on a star-by-star basis. The nature of low density dust provides insights into the processing which occurs in hostile environments. Clayton et al. (2000) presented observations of dust from low density sightlines, and their results were intriguing. However, they were forced to select sightlines which accumulate fairly substantial color excesses, introducing the possibility of mixed grain populations. Furthermore, the results for their least-reddened sightlines were, as they acknowledge, poorly determined, forcing them to base their conclusions on a global average of properties of their sample. Finally, measuring the properties of local dust is important because it allows us to search for isolated, relatively homogeneous environments with uniform extinction curve shapes. These may signal physically and kinematically isolated regions which would be ideal for follow up interstellar line studies. In addition, Hipparcos parallax data exist for nearby stars and provide an opportunity to study the 3-dimensional structure of local extinction.
Access to the mid-to-late B stars: These stars are especially important because their space density is higher than that of the early-B stars usually used in extinction studies and thus their inclusion increases the number of stars available to create curves for nearby sightlines. Study of the mid-to-late B stars will enable the examination of the spatial structure of local dust absorption more thoroughly than previously possible and will greatly enhance our understanding of the local interstellar medium. The ability to construct accurate curves for mid-to-late B stars will also allow us to study extinction in open clusters and associations which are dominated by these stars, such as the Pleiades, $`\alpha `$ Per cluster, and IC 4665 (see Figure 4).
Automation: Because our model atmosphere approach does not require human intervention, once the basic data have been assembled, it is possible to process large data sets at one time. Naturally, the automated results must be inspected for outliers and data anomalies. Nevertheless, this approach reduces the work load considerably and a first-attempt yielded the regional anomaly shown in Figure 9. Furthermore, since the results are produced in a uniform manner, it is a relatively simple matter to inspect them for correlations between various curve properties, for anomalous curve shapes, and for spatial trends on the sky.
Stellar properties: In addition to dust parameters, our technique provides a meaningful, quantitative physical properties for the reddened stars. The temperature and surface gravity information will be useful for population studies of B stars in the field and in clusters. However, perhaps the most useful property will be the metallicity. We have demonstrated that several stars in the same cluster, which have a range in temperatures and gravities, all yield the same \[m/H\]. This verifies the sensitivity of our fitting procedure to this important quantity. As a result, we are confident that application of our procedure to large scale surveys of reddened, near main sequence B stars can provide a census of the distribution of metallicity throughout the Galaxy and the local universe.
E.F. acknowledges support from NASA grant NAG5-12137, NAG5-10385, and NNG04GD46G. D.M. acknowledges support from NASA grant NNG04EC01P. Some of the data presented in this paper were obtained from the Multimission Archive at the Space Telescope Science Institute (MAST). STScI is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. Support for MAST for non-HST data is provided by the NASA Office of Space Science via grant NAG5-7584 and by other grants and contracts. This publication also makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. |
warning/0506/math0506105.html | ar5iv | text | # Approximations of Set-Valued Functions by Metric Linear Operators
## 1 Introduction
In this work, we adapt approximating operators for real-valued functions to set-valued functions (multifunctions, SVFs), by replacing an operation between numbers with an operation between sets. The known approximation methods, based on Minkowski sums of sets, fail to approximate, when the images of a multifunction are not convex. In case of Bernstein-type operators and subdivision operators there is a phenomenon of โconvexificationโ ().
In a binary operation between sets, the โmetric averageโ, is introduced and the metric piecewise linear interpolant based on it is shown to approximate continuous SVFs with general images. The use of this operation in the adaptation of known approximation methods to SVFs, requires a representation of the approximation operators by repeated binary averages. Such a representation exists for any operator which reproduces constants, but is not unique . This non-uniqueness leads to different operators and it is not clear what are the appropriate adaptations. Spline subdivision schemes represented by repeated averages and the Schoenberg operators defined in terms of the de Boor algorithm are proved to approximate SVFs with general compact images. Yet, for the adaptation of the Bernstein operators based on the de Casteljau algorithm we could obtain an approximation result only for SFVs with images in $`R`$ all consisting of the same number of disjoint intervals .
In this paper we introduce a set-operation on a finite sequence of compact sets, termed โmetric linear combinationโ, which extends the metric average. We adapt approximation methods for real-valued functions to SVFs, replacing linear combinations of numbers by the metric linear combinations of sets. We prove that this adaptation of any linear operator, approximating continuous real-valued functions, approximates continuous SVFs of bounded variation. In particular for Lipschitz continuous SVFs, sharper error estimates are obtained. Approximation results for set-valued functions which are only continuous, are given for a limited class of operators. It should be noted that our adaptation method is not restricted to positive operators. The approximation results are specialized to the Shoenberg spline operators and the Bernstein polynomial operators. Also the adaptation of polynomial interpolation to SVFs is presented as examples of non-positive operators. This adaptation is illustrated by two metric parabolic interpolants.
An outline of the paper is as follows. The next section contains basic definitions, notation and known results. In Section 3 we introduce the metric linear combination between a finite number of ordered sets and define metric linear operators for multifunctions based on it. In Section 4 properties of the metric piecewise linear intorpolant are considered. In particular a representation of it by a specific set of selections is studied. Similar selections are used in and to prove the existence of a continuous selection for a continuous SVF of bounded variation, and of a representation of a Lipschitz SVF respectively. In Section 5 we derive approximation results for the metric linear approximation operators, based on the results in Section 4. Finally, in Section 6, we specialize these results to some classical approximation operators.
## 2 Preliminaries
First we present some definitions and notation.
$``$ $`K(R^n)`$ is the collection of all compact nonempty subsets of $`R^n`$.
$``$ A linear Minkowski combination of two sets $`A`$ and $`B`$ from $`K(R^n)`$ is
$$\lambda A+\mu B=\{\lambda a+\mu b:aA,bB\},$$
with $`\lambda ,\mu R`$.
$``$ The Euclidean distance from a point $`aR^n`$ to a set $`BK(R^n)`$ is defined as
$$\text{dist}(a,B)=\underset{bB}{inf}|ab|,$$
where $`||`$ is the Euclidean norm in $`R^n`$.
$``$The Hausdorff distance between two sets $`A,BK(R^n)`$ is defined by
$$\text{haus}(A,B)=\mathrm{max}\{\underset{aA}{sup}\text{dist}(a,B),\underset{bB}{sup}\text{dist}(b,A)\}.$$
$``$ The set of all projections of $`aR^n`$ into a set $`BK(R^n)`$ is
$$\mathrm{\Pi }_B(a)=\{bB:|ab|=\text{dist}(a,B)\}.$$
$``$ For $`A,BK(R^n)`$ and $`0t1`$, the t-weighted metric average of $`A`$ and $`B`$ is
$$A_tB=\{ta+(1t)b:(a,b)\mathrm{\Pi }(A,B)\}$$
(1)
with $`\mathrm{\Pi }(A,B)=\{(a,b)A\times B:a\mathrm{\Pi }_A(b)\text{or}b\mathrm{\Pi }_B(a)\}`$.
The metric average has the metric property
$`\text{haus}(A_tB,A_sB)`$ $`=`$ $`|ts|\text{haus}(A,B),`$
$`\text{haus}(A_tB,A)`$ $`=`$ $`(1t)\text{haus}(A,B),`$ (2)
$`\text{haus}(A_tB,B)`$ $`=`$ $`t\text{haus}(A,B).`$
$``$ The modulus of continuity of $`f:[a,b]X`$ with images in a metric space $`(X,\rho )`$ is
$$\omega _{[a,b]}(f,\delta )=sup\{\rho (f(x),f(y)):|xy|\delta ,x,y[a,b]\},\delta >0.$$
(3)
In this paper $`X`$ is either $`R^n`$ or $`K(R^n)`$, and $`\rho `$ is either the Euclidean distance or the Hausdorff distance respectively.
The property of the modulus that we use is
$$\omega _{[a,b]}(f,\lambda \delta )(1+\lambda )\omega _{[a,b]}(f,\delta ).$$
(4)
$``$ By $`Lip([a,b],)`$ we denote the set of all Lipschitz functions $`f:[a,b]X`$ satisfying
$$\rho (f(x),f(y))|xy|,x,y[a,b],$$
where $``$ is a constant independent of $`x`$ and $`y.`$
$``$ A variation of $`f:[a,b]X`$ on a partition $`\chi =\{x_0<\mathrm{}<x_N:x_i[a,b],i=0,\mathrm{},N\}`$ is defined by
$$V(f,\chi )=\underset{i=1}{\overset{N}{}}\rho (f(x_i),f(x_{i1})),$$
The total variation of $`f`$ on $`[a,b]`$ is
$$V_a^b(f)=\underset{\chi }{sup}V(f,\chi ).$$
We say that $`f`$ is of bounded variation if $`V_a^b(f)<\mathrm{}`$ and define in this case
$$v_f(x)=V_a^x(f),x[a,b].$$
(5)
It is obvious that $`v_f`$ is nondecreasing. If $`f`$ is also continuous then $`v_f`$ is continuous as well. For completeness we prove it.
###### Proposition 2.1.
A function $`f:[a,b]X`$ is continuous and of bounded variation on $`[a,b]`$ if and only if $`v_f`$ is a continuous function on \[a,b\].
###### Proof.
The sufficiency follows from
$$\rho (f(x),f(y))V_x^y(f)=v_f(y)v_f(x),\mathrm{for}x<y.$$
(6)
To prove the other direction, fix $`x[a,b]`$ and $`\epsilon >0`$. By the uniform continuity of $`f`$ on $`[a,b]`$, $`\rho (f(z),f(y))<\epsilon /2`$ if $`|zy|<\delta `$ for some $`\delta >0`$. First we show that $`v_f`$ is continuous from the left. We can always choose $`\chi =\{a=x_0<x_1<\mathrm{}<x_N=x\}`$ such that
$$V_a^x(f)<V(f,\chi )+\epsilon /2=\underset{i=1}{\overset{N}{}}\rho (f(x_i),f(x_{i1}))+\epsilon /2,$$
and $`xx_{N1}<\delta `$. Thus
$$V_a^x(f)<\underset{i=1}{\overset{N1}{}}\rho (f(x_i),f(x_{i1}))+\epsilon ,$$
implying that $`v_f(x)v_f(x_{N1})<\epsilon `$. By the monotonicity of $`v_f`$we get for every $`x_{N1}<y<x`$
$$v_f(x)v_f(y)<\epsilon .$$
Similarly one can show the continuity of $`v_f`$ from the right. Thus we obtain that $`v_f`$ is continuous at $`x`$ and consequently it is continuous on $`[a,b]`$. โ
From (6) we conclude that
$$\omega _{[a,b]}(f,\delta )\omega _{[a,b]}(v_f,\delta ).$$
(7)
$``$ By CBV we denote the set of all functions $`f:[a,b]X`$ which are continuous and of bounded variation.
$``$ For a set-valued function $`F:[a,b]K(R^n)`$, any single-valued function $`f:[a,b]R^n`$ with $`f(x)F(x)`$, $`x[a,b]`$ is called a selection of $`F`$.
###### Definition 2.2.
A set of selections of $`F`$, $`\{f^\alpha :\alpha ๐\}`$, is termed a representation of $`F`$ if
$$F(x)=\{f^\alpha (x):\alpha ๐\},x[a,b].$$
We denote this shortly by $`F=\{f^\alpha :\alpha ๐\}.`$
## 3 Linear operators on SVFs based on a metric linear combination of ordered sets
In this section we introduce a new operation on a finite number of ordered sets. Using this operation we present a new adaptation of linear operators to multifunctions.
###### Definition 3.1.
Let $`\{A_0,A_1,\mathrm{},A_N\}`$ be a finite sequence of compact sets. A vector $`(a_0,a_1,\mathrm{},a_N)`$ with $`a_iA_i`$, $`i=0,\mathrm{},N`$, for which there exists $`j`$, $`0jN`$ such that
$$a_{i1}\mathrm{\Pi }_{A_{i1}}(a_i),\mathrm{\hspace{0.17em}1}ij\text{and}a_{i+1}\mathrm{\Pi }_{A_{i+1}}(a_i),jiN1$$
is called a metric chain of $`\{A_0,\mathrm{},A_N\}`$.
An illustration of such a metric chain is given in Figure 3.2.
$$a_0\mathrm{\Pi }_{A_0}\left(a_1\right)a_{j1}\mathrm{\Pi }_{A_{j1}}\left(a_j\right)a_jA_ja_{j+1}\mathrm{\Pi }_{A_{j+1}}\left(a_j\right)a_N\mathrm{\Pi }_{A_N}\left(a_{N1}\right)$$
###### Figure 3.2.
Thus each element of each set $`A_i`$, $`i=0,\mathrm{},N`$ generates at least one metric chain. We denote by $`CH(A_0,\mathrm{},A_N)`$ the collection of all metric chains of $`\{A_0,\mathrm{},A_N\}`$. The set $`CH(A_0,\mathrm{},A_N)`$ depends on the order of the sets $`A_i`$, $`i=0,\mathrm{},N`$.
With this notion of metric chains we can introduce a new operation between sets.
###### Definition 3.3.
A metric linear combination of a sequence of sets $`A_0,\mathrm{},A_N`$ with coefficients $`\lambda _0,\mathrm{},\lambda _NR`$, is
$$\underset{i=0}{\overset{N}{}}\lambda _iA_i=\{\underset{i=0}{\overset{N}{}}\lambda _ia_i:(a_0,\mathrm{},a_N)CH(A_0,\mathrm{},A_N)\}.$$
(8)
Since for two sets $`CH(A,B)=\mathrm{\Pi }(A,B)`$, in the special case $`N=1`$ and $`\lambda _0,\lambda _1[0,1]`$, $`\lambda _0+\lambda _1=1`$, the metric linear combination is the metric average. The following are two important properties of the metric linear combination which can be easily seen from the definition.
$`\begin{array}{cc}(\mathrm{i})\hfill & \underset{i=0}{\overset{N}{}}\lambda _iA_i=\underset{i=0}{\overset{N}{}}\lambda _{Ni}A_{Ni},\hfill \\ (\mathrm{ii})\hfill & \mathrm{For}\lambda _0,\mathrm{},\lambda _N\mathrm{such}\mathrm{that}\underset{i=0}{\overset{N}{}}\lambda _i=1,\underset{i=0}{\overset{N}{}}\lambda _iA=A.\hfill \end{array}`$
With this operation, a large class of linear operators can be adapted to SVFs.
Let $`A_\chi `$, $`\chi =\{x_0,\mathrm{},x_N\}`$ be a linear operator of the form
$$A_\chi (f,x)=\underset{i=0}{\overset{N}{}}c_i(x)f(x_i),$$
(9)
defined on real-valued functions, with domain containing $`\chi `$.
###### Definition 3.4.
Let $`F:[a,b]K(R^n)`$, $`\chi [a,b]`$ and let $`\{F(x_i),i=0,\mathrm{},N\}`$ be samples of $`F`$ at $`\chi `$. For $`A_\chi `$ of the form (9), we define a metric linear operator $`A_\chi ^M`$ on $`F`$ by
$$A_\chi ^MF(x)=A_\chi ^M(F,x)=\underset{i=0}{\overset{N}{}}c_i(x)F(x_i).$$
(10)
We term this operator the metric analogue of (9).
Note that due to property (ii), the metric analogue of a linear operator which preserves constants, preserves constant multifunctions. The analogue of property (ii) does not hold for Minkowski linear combinations with some negative coefficients, even for convex sets. This is one reason why only positive operators, based on Minkowski sum, were applied to set-valued functions. As is shown in the sequel, Definition 3.4 allows to define also non-positive operators.
The analysis of the approximation properties of $`A_\chi ^MF`$ is based on properties of the metric piecewise linear approximation operator.
## 4 Metric piecewise linear approximations of SVFs
ยฟFrom now on $`F:[a,b]K(R^n)`$, $`\{F_i=F(x_i)\}_{i=0}^N`$, where $`a=x_0<\mathrm{}<x_N=b`$ and $`\chi =(x_0,\mathrm{},x_N)`$ denotes a partition of $`[a,b]`$. We use the notation $`CH=CH(F_0,\mathrm{},F_N)`$, and $`\delta _{max}=\mathrm{max}\{\delta _i:\mathrm{\hspace{0.33em}\hspace{0.17em}0}iN1\}`$, $`\delta _{min}=\mathrm{min}\{\delta _i:\mathrm{\hspace{0.33em}\hspace{0.17em}0}iN1\}`$ with values $`\delta _i`$ defined as $`\delta _i=(x_{i+1}x_i)`$, $`i=0,\mathrm{},N1`$. In case of a uniform partition, we have $`\delta _{max}=\delta _{min}=h=(ba)/N`$ and denote such a partition by $`\chi _N`$.
###### Definition 4.1.
The metric piecewise linear approximation to $`F`$ is
$$S_\chi ^M(F,x)=\{\lambda _i(x)f_i+(1\lambda _i(x))f_{i+1}:(f_0,\mathrm{},f_N)CH\},x[x_i,x_{i+1}],$$
where
$$\lambda _i(x)=(x_{i+1}x)/(x_{i+1}x_i).$$
(11)
By construction, the set-valued function $`S_\chi ^MF`$ has a representation by selections
$$S_\chi ^MF=\{s(\chi ,\phi ):\phi CH(F_0,\mathrm{},F_N)\},$$
(12)
where $`s(\chi ,\phi )`$ is a piecewise linear single-valued function interpolating the data $`(x_i,f_i),`$ $`i=0,\mathrm{},N,`$ with $`\phi =(f_0,\mathrm{},f_N)`$.
Recall the piecewise linear interpolant based on the metric average, introduced in :
$$S_\chi ^{MA}(F,x)=F_i_{\lambda _i(x)}F_{i+1},x[x_i,x_{i+1}]$$
with $`\lambda _i(x)`$ defined by (11).
It is easy to see by the triangle inequality for the Hausdorff metric and by (2) that for a continuous set-valued function $`F`$
$$\mathrm{haus}(F(x),S_\chi ^{MA}(F,x))2\omega _{[a,b]}(F,\delta _{max}),x[a,b].$$
(13)
###### Remark 4.2.
It is not unexpected that $`S_\chi ^{MA}FS_\chi ^MF`$.
Indeed, for a fixed $`x[x_i,x_{i+1}]`$ and for any $`yS_\chi ^{MA}(F,x)`$,
$$y=\lambda _i(x)f_i+(1\lambda _i(x))f_{i+1}$$
with $`(f_i,f_{i+1})\mathrm{\Pi }(F_i,F_{i+1})`$. Then there exists a metric chain $`\phi =(f_0,\mathrm{},f_i,f_{i+1},\mathrm{},f_N)`$, $`\phi CH`$, such that $`y=s(\chi ,\phi )(x)`$. Also it is obvious that for any $`x[a,b]`$ and any $`\phi CH`$, $`s(\chi ,\phi )(x)S_\chi ^{MA}(F,x)`$.
In the following we show that $`S_\chi ^MF`$, and its piecewise linear selections (12) โinheritโ some continuity properties of a continuous multifunction $`F`$. The following lemma and corollary consider Lipschitz continuous SVFs.
###### Lemma 4.3.
Let $`FLip([a,b],)`$, and let $`\chi `$ be a partition of $`[a,b]`$. Then
$$S_\chi ^MFLip([a,b],).$$
###### Proof.
For $`x,y[x_j,x_{j+1}]`$ the claim of the lemma follows from the metric property (2). Now, let $`x[x_j,x_{j+1}]`$ and $`y[x_k,x_{k+1}],`$ where $`0jkN1.`$ Using the triangle inequality, (2) and the Lipschitz continuity of $`F`$, we get
$`\text{haus}(S_\chi ^M(F,x),S_\chi ^M(F,y))`$
$`{\displaystyle \frac{x_{j+1}x}{x_{j+1}x_j}}\mathrm{haus}(F_j,F_{j+1})+\mathrm{haus}(F_{j+1},F_k)+{\displaystyle \frac{yx_k}{x_{k+1}x_k}}\mathrm{haus}(F_k,F_{k+1})`$
$`(x_{j+1}x+x_kx_{j+1}+yx_k)|yx|.`$
###### Corollary 4.4.
Under the conditions of Lemma 4.3, for any $`s(\chi ,\phi )`$ in 12
$$s(\chi ,\phi )Lip([a,b],).$$
The proof of this corollary is similar to the proof of the previous lemma and uses the observation that for $`kj`$
$`|s(\chi ,\phi )(x_{j+1})s(\chi ,\phi )(x_k)|{\displaystyle \underset{l=j+1}{\overset{k1}{}}}|s(\chi ,\phi )(x_l)s(\chi ,\phi )(x_{l+1})|`$
$`{\displaystyle \underset{l=j+1}{\overset{k1}{}}}\mathrm{haus}(S_\chi ^M(F,x_l),S_\chi ^M(F,x_{l+1})){\displaystyle \underset{l=j+1}{\overset{k1}{}}}(x_{l+1}x_l)=|x_kx_{j+1}|.`$
Now we consider the case when $`F`$ is a general continuous function.
###### Lemma 4.5.
Let $`F:[a,b]K(R^n)`$ be a continuous set-valued function. Then for any partition $`\chi `$ of $`[a,b]`$
$$\omega _{[a,b]}(S_\chi ^MF,\delta )5\omega _{[a,b]}(F,\delta ).$$
###### Proof.
By definition, for any $`\delta >0`$
$$\omega _{[a,b]}(S_\chi ^MF,\delta )=sup\{\mathrm{haus}(S_\chi ^M(F,x),S_\chi ^M(F,y)):|xy|\delta ,x,y[a,b]\}.$$
In case $`x,y[x_j,x_{j+1}]`$, $`|xy|\delta `$, the claim of the lemma is obtained using (11), the metric property (2) and (4),
$$\begin{array}{cc}\hfill \mathrm{haus}(S_\chi ^M(F,x),S_\chi ^M(F,y))& =\frac{|xy|}{\delta _j}\mathrm{haus}(F_j,F_{j+1})\frac{\delta }{\delta _j}\omega _{[a,b]}(F,\delta _j)\hfill \\ & \frac{\delta }{\delta _j}\left(1+\frac{\delta _j}{\delta }\right)\omega _{[a,b]}(F,\delta )2\omega _{[a,b]}(F,\delta ).\hfill \end{array}$$
(14)
Now, let $`x[x_j,x_{j+1}]`$, $`y[x_k,x_{k+1}],`$ $`0j<kN1`$ and $`|xy|\delta `$. By the triangle inequality
$$\begin{array}{cc}\hfill \mathrm{haus}(S_\chi ^M(F,x),S_\chi ^M(F,y))& \mathrm{haus}(S_\chi ^M(F,x),S_\chi ^M(F,x_{j+1}))\hfill \\ & +\mathrm{haus}(S_\chi ^M(F,x_{j+1}),S_\chi ^M(F,x_k))\hfill \\ & +\mathrm{haus}(S_\chi ^M(F,x_k),S_\chi ^M(F,y)),\hfill \end{array}$$
(15)
while by the interpolation property of $`S_\chi ^MF`$ and since $`|x_kx_{j+1}|\delta `$, we have
$$\mathrm{haus}(S_\chi ^M(F,x_{j+1}),S_\chi ^M(F,x_k))\omega _{[a,b]}(F,\delta ).$$
(16)
Applying (14) and (16) to (15) we obtain the claim of the lemma. โ
###### Corollary 4.6.
For any $`s(\chi ,\phi )`$ in (12) and any $`x,y[x_j,x_{j+1}]`$, $`0jN1`$
$$|s(\chi ,\phi )(x)s(\chi ,\phi )(y)|2\omega _{[a,b]}(F,|xy|)$$
(17)
Also,
$$\omega _{[a,b]}(s(\chi ,\phi ),\delta )4\omega _{[a,b]}(F,\delta ),\delta \delta _{min}.$$
(18)
The proof of this corollary is similar to the proof of assertion (14).
We cannot generalize (18) for arbitrary $`0<\delta ba`$ if $`F`$ is only continuous. Yet we can get an estimate for $`\omega _{[a,b]}(s(\chi ,\phi ),\delta )`$ if $`F`$ is continuous and of bounded variation.
###### Lemma 4.7.
Let $`FCBV([a,b])`$. Then for any $`s(\chi ,\phi )`$ in (12),
$$\omega _{[a,b]}(s(\chi ,\phi ),\delta )4\omega _{[a,b]}(F,\delta )+\omega _{[a,b]}(v_F,\delta )5\omega _{[a,b]}(v_F,\delta ).$$
###### Proof.
Denote $`s=s(\chi ,\phi )`$. For a given $`\delta >0`$, let $`x[x_j,x_{j+1}]`$, $`y[x_k,x_{k+1}]`$, $`0jkN1`$, such that $`|xy|\delta `$. Using the definition of $`s(\chi ,\phi )`$ and of $`S_\chi ^MF`$ we get
$`|s(x)s(y)||s(x)s(x_{j+1})|+{\displaystyle \underset{l=j+1}{\overset{k1}{}}}|s(x_{l+1})s(x_l)|+|s(y)s(x_k)|`$
$`{\displaystyle \frac{x_{j+1}x}{\delta _j}}|s(x_{j+1})s(x_j)|+{\displaystyle \underset{l=j+1}{\overset{k1}{}}}\mathrm{haus}(F(x_{l+1}),F(x_l))+{\displaystyle \frac{yx_k}{\delta _k}}|s(x_{k+1})s(x_k)|.`$
Now, (17) yields
$$|s(x)s(y)|4\omega _{[a,b]}(F,\delta )+V_{x_{j+1}}^{x_k}(F)4\omega _{[a,b]}(F,\delta )+\omega _{[a,b]}(v_F,\delta ).$$
Taking the supremum over $`|xy|\delta `$ and using (7), we complete the proof. โ
## 5 Approximation by metric linear operators
We use the metric piecewise linear approximation to obtain error estimates for metric linear operators.
Let $`A_\chi ^MF`$ be defined by (10) and $`S_\chi ^MF`$ be a metric piecewise linear multifunction as defined in Section 4. By Definition 3.4
$$A_\chi ^MFA_\chi ^M(S_\chi ^MF).$$
(19)
Moreover by (9), (10) and (12)
$$A_\chi ^M(S_\chi ^MF)=\{A_\chi s(\chi ,\phi ):\phi CH(F_0,\mathrm{},F_N)\}.$$
(20)
###### Remark 5.1.
In contrast to our previous definition of positive operators for SVFs based on the metric average , the metric analogues (10) of two linear operators of the form (9), which are identical on single-valued functions, are identical on SVFs. For example, in spline subdivision schemes are not identical to the Schoenberg spline operators for SVFs.
The metric analogues of linear operators of the form (9), which approximate real-valued functions, are approximating SVFs. By (19), (20) the approximation results depend on the way $`A_\chi `$ approximates piecewise linear real-valued functions.
In what follows $`\varphi :[a,b]\times R_+R_+`$ is a continuous real-valued function, non-decreasing in its second argument, satisfying $`\varphi (x,0)=0`$, and $`๐ฎ_\chi `$ denotes the set of piecewise linear continuous single-valued functions, with values in $`R^n`$ and knots at $`\chi `$.
###### Theorem 5.2.
Let $`A_\chi `$ be of the form (9), such that for any $`s๐ฎ_\chi Lip([a,b],)`$
$$|A_\chi (s,x)s(x)|C\varphi (x,\delta _{max}).$$
(21)
Then if $`FLip([a,b],)`$,
$$\mathrm{haus}(A_\chi ^M(F,x),F(x))=2\delta _{max}+C\varphi (x,\delta _{max}).$$
###### Proof.
By (19)
$$\begin{array}{c}\hfill \mathrm{haus}(A_\chi ^M(F,x),F(x))\mathrm{haus}(A_\chi ^M(S_\chi ^MF,x),S_\chi ^M(F,x))+\mathrm{haus}(S_\chi ^M(F,x),F(x)),\end{array}$$
(22)
while by (20)
$$\mathrm{haus}(A_\chi ^M(S_\chi ^MF,x),S_\chi ^M(F,x))\underset{\phi CH}{sup}|A_\chi (s(\chi ,\phi ),x)s(\chi ,\phi )(x)|.$$
In view of Corollary 4.4 and (21)
$$\underset{\phi CH}{sup}|A_\chi (s(\chi ,\phi ),x)s(\chi ,\phi )(x)|C\varphi (x,\delta _{max})$$
(23)
The proof is completed by substituting (23) and (13) in (22). โ
For general continuous SVFs we cannot prove an analogous approximation result. Yet for continuous multifunctions of bounded variation we get a weaker approximation result, by applying Lemma 4.7 instead of Corollary 4.4 in the proof of Theorem 5.2.
###### Theorem 5.3.
Let $`FCBV([a,b])`$, and let $`A_\chi `$ be of the form (9), satisfying
$$|A_\chi (s,x)s(x)|C\omega _{[a,b]}(s,\varphi (x,\delta _{max})),s๐ฎ_\chi .$$
(24)
Then
$$\mathrm{haus}(A_\chi ^M(F,x),F(x))=2\omega _{[a,b]}(F,\delta _{max})+5C\omega _{[a,b]}(v_F,\varphi (x,\delta _{max})).$$
For continuous SVFs which are not of bounded variation we can prove an approximation result only for uniform partitions and for a limited class of linear operators.
###### Theorem 5.4.
Let $`A_N`$ be a linear operator of the form (9), defined on a uniform partition $`\chi _N`$ and let $`h=(ba)/N`$. If
$$|A_N(s,x)s(x)|C\varphi (x,\omega _{[a,b]}(s,h)),s๐ฎ_\chi ,$$
(25)
then for a continuous $`F`$
$$\mathrm{haus}(A_N^M(F,x),F(x))=2\omega _{[a,b]}(F,h)+C\varphi (x,4\omega _{[a,b]}(F,h)).$$
The proof of this result repeats the proof of Theorem 5.2, but replaces Corollary 4.4 by (18) of Corollary 4.6.
## 6 Examples
In this section we present metric analogues for SVFs of the Schoenberg spline operators and the Bernstein polynomial operators and give approximation results. We conclude by two examples demonstrating the operation of metric analogues of parabolic interpolants. To our knowledge so far only positive operators were applied to SVFs. The two examples we present assert that such interpolation between sets is reasonable.
### 6.1 Metric Bernstein operators
The Benstein operator $`B_N(f,x)`$ for $`fC[0,1]`$ is
$$B_N(f,x)=\underset{i=0}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{i}\right)x^i(1x)^{Ni}f\left(\frac{i}{N}\right).$$
(26)
It is known (see , Chapter 10) that there exists a constant $`C`$ independent of $`f`$ such that
$$|f(x)B_N(f,x)|C\omega _{[0,1]}(f,\sqrt{x(1x)/N}).$$
(27)
The classical Bernstein operator for $`F:[0,1]K(R^n)`$ with sums of numbers replaced by Minkowski sums of sets is
$$B_N^{Mn}(F,x)=\underset{i=0}{\overset{N}{}}\left(\genfrac{}{}{0pt}{}{N}{i}\right)x^i(1x)^{Ni}F\left(\frac{i}{N}\right).$$
(28)
It was shown in that for $`x(0,1)`$ the limit of $`B_N^{Mn}(F,x)`$ when $`N\mathrm{},`$ is the convex hull of $`F(x)`$, therefore these operators cannot approximate SVFs with general images.
In Bernstein operators for set-valued functions are defined procedurally in terms of the de Casteljau algorithm, with the metric average as a basic binary operation,
$`F_i^0=F(i/N),i=0,\mathrm{},N,`$
$`F_i^k=F_i^{k1}_{\mathrm{\hspace{0.17em}1}x}F_{i+1}^{k1},i=0,1,\mathrm{},Nk,k=1,\mathrm{},N,`$ (29)
$`B_N^{MA}(F,x)=F_0^N.`$
We do not know whether these operators approximate multifunctions with general compact images in $`R^n`$, yet they approximate multifunctions with compact images in $`R`$ all consisting of the same number of disjoint intervals .
Here we investigate the metric analogue of the Bernstein operators for SVFs.
###### Definition 6.1.
For $`F:[0,1]K(R^n)`$ the metric Bernstein operator is
$`B_N^M(F,x)`$ $`={\displaystyle \underset{i=0}{\overset{N}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{N}{i}}\right)x^i(1x)^{Ni}F\left({\displaystyle \frac{i}{N}}\right)`$
$`=\{{\displaystyle \underset{i=0}{\overset{N}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{N}{i}}\right)x^i(1x)^{Ni}f_i:(f_0,\mathrm{},f_N)CH\},`$
where $`CH=CH(F(0),F(1/N)\mathrm{},F(1))`$.
By Theorem 5.2 and by (27) we conclude that
###### Corollary 6.2.
Let $`FLip([0,1],)`$, then
$$\mathrm{haus}(B_N^M(F,x),F(x))2/N+C\sqrt{x(1x)/N}.$$
Moreover by Theorem 5.3 and by (27)
###### Corollary 6.3.
Let $`FCBV([0,1])`$, then
$$\mathrm{haus}(B_N^M(F,x),F(x))2\omega _{[0,1]}(F,1/N)+5C\omega _{[0,1]}(v_F,\sqrt{x(1x)/N}).$$
Since (25) does not hold for these operators, Theorem 5.4 cannot be applied.
### 6.2 Metric Schoenberg operators
For a uniform partition $`\chi _N`$, the โclassicalโ set-valued analogues of the Schoenberg spline operators for $`F:[0,1]K(R^n)`$ is
$$S_{m,N}^{Mn}(F,x)=\underset{i=0}{\overset{N}{}}F(i/N)b_m\left(Nxi\right),$$
(30)
where $`b_m\left(x\right)`$ is the B-spline of order $`m`$ (degree $`m1`$) with integer knots and support $`[0,m]`$, and where the linear combination is in the Minkowski sense. An example, given in , shows that the operators in (30) with $`m=2`$ and $`N\mathrm{}`$ cannot approximate $`F`$ with general compact images, in any point of $`[0,1]\chi _N`$.
A Shoenberg operator based on the metric average is introduced in , by a procedural definition in terms of repeated binary averages according to the de Boor algorithm. It is proved that for Hรถlder continuous set-valued functions, the approximation rate is the Hรถlder exponent.
Here we consider the metric analogue of the Schoenberg operators.
###### Definition 6.4.
The metric Shoenberg operator of order $`m`$ for a set-valued function $`F:[0,1]K(R^n)`$ and a uniform partition $`\chi _N`$ is defined by
$`S_{m,N}^M(F,x)={\displaystyle \underset{i=0}{\overset{N}{}}}b_m\left(Nxi\right)F\left({\displaystyle \frac{i}{N}}\right)=\{{\displaystyle \underset{i=0}{\overset{N}{}}}b_m\left(Nxi\right)f_i:(f_0,\mathrm{},f_N)CH\},`$
where $`CH=CH(F(0),F(1/N)\mathrm{},F(1))`$.
By Theorem 5.4 and the known approximation result in case of single-valued functions (see , Chapter XII), we obtain
###### Corollary 6.5.
Let $`F`$ be a continuous SFV defined on $`[0,1]`$. Then
$$\mathrm{haus}(S_{m,N}^M(F,x),F(x))=2\left(1+2\frac{m+1}{2}\right)\omega _{[0,1]}(F,1/N),x[\frac{m1}{N},1]$$
with $`t`$ the maximal integer not greater than $`t`$.
The approximation result in the specific case of Lipschitz continuous SVFs, can be further improved by applying Theorem 5.2.
###### Corollary 6.6.
For $`FLip([0,1],)`$,
$$\mathrm{haus}(S_{m,N}^M(F,x),F(x))=\left(2+\frac{m+1}{2}\right)\frac{}{N}.$$
### 6.3 Metric polynomial interpolants
###### Definition 6.7.
Let $`F:[a,b]K(R^n)`$, and let $`\chi `$ be a partition of $`[a,b]`$. The metric polynomial interpolation operator is given by
$`P_\chi ^M(F,x)={\displaystyle \underset{i=0}{\overset{N}{}}}l_i(x)F(x_i)=\{{\displaystyle \underset{i=0}{\overset{N}{}}}l_i(x)f_i:(f_0,\mathrm{},f_N)CH(F(x_0),\mathrm{},F(x_N))\},`$
with $`l_i(x)`$ the $`i`$-th Lagrange polynomial,
$$l_i(x)=\underset{j=0,ji}{\overset{N}{}}\frac{xx_j}{x_ix_j}.$$
To illustrate our method we apply the metric parabolic interpolation operator to three sets in $`R`$. We consider two different examples.
The first example: $`x_0=0`$, $`x_1=2`$, $`x_2=6`$;
$$F(x_0)=[2,8],F(x_1)=\{5\},F(x_2)=\{5\}.$$
The second example: $`x_0=0`$, $`x_1=4`$, $`x_2=8`$;
$$F(x_0)=[2,4][6,8],F(x_1)=[4.5,5.5],F(x_2)=[2,4][6,8].$$
The two set-valued interpolants are illustrated in Figure 6.8 and Figure 6.9 respectively.
###### Figure 6.8.
Parabolic interpolation - first example.
###### Figure 6.9.
Parabolic interpolation - second example.
In the above figures the sets in black are $`F(x_0)`$, $`F(x_1)`$, $`F(x_2)`$ and the gray curves are the parabolic interpolants to the selections in (12). |
warning/0506/cond-mat0506352.html | ar5iv | text | # Saturated Ferromagnetism from Statistical Transmutation in Two Dimensions
## Abstract
The total spin of the ground state is calculated in the $`U\mathrm{}`$ Hubbard model with uniform magnetic flux perpendicular to a square lattice, in the absence of Zeeman coupling. It is found that the saturated ferromagnetism emerges in a rather wide region in the space of the flux density $`\varphi `$ and the electron density $`n_\mathrm{e}`$. In particular, the saturated ferromagnetism at $`\varphi =n_\mathrm{e}`$ is induced by the formation of a spin-1/2 boson, which is a composite of an electron and the unit flux quantum.
Ferromagnetism remains a challenging problem in spite of being among the best known phenomena in condensed matter physics. In particular, the origin of ferromagnetism in electron systems is fundamentally quantum-mechanical and non-perturbative TasakiReview . There are rather few established mechanisms of ferromagnetism, especially those of the saturated (complete) ferromagnetism. Saturated ferromagnetism is defined as the ground state of the many-electron system with spin-independent interaction having the maximum possible total spin.
Nagaokaโs theorem is one of few rigorous results on the saturated ferromagnetism Nagaoka . It guarantees that the saturated-ferromagnetic state is the unique ground state when a single hole is inserted in the half-filled Hubbard model with infinite on-site repulsion $`U`$. Unfortunately, this theorem is limited to the single-hole case. Numerical studies suggest that Nagaoka ferromagnetism is unstable in the thermodynamic limit at finite hole densities Putikka-etal .
The flat-band ferromagnetism is another rigorous result Mielke ; Tasaki92 . Namely, the saturated ferromagnetism is proved rigorously under certain conditions, in electron systems with a (nearly) flat dispersion in the lowest band for a single electron. The ferromagnetism in a system with the low electron density and singular density of states near the Fermi level Hlubina-etal may also be related to the flat-band mechanism. However, as these results still have limited applicability, it is worth pursuing other mechanisms of (saturated) ferromagnetism.
The difficulty in realizing the ferromagnetism in many-electron systems may be attributed to the Pauli principle for electrons. In the absence of interaction, the ground state of the system is generally paramagnetic, because the lower energy bands are filled with up and down spins. If we consider a system of bosons rather than fermions, the intrinsic tendency to favor paramagnetism may be absent. In fact, it was proved that in a continuous system with spinful bosons, one of the ground states is always fully polarized if explicit spin-dependent interactions are absent Suto ; Eisenberg-Lieb . This statement holds also in a lattice model Eisenberg-Lieb ; Fledderjohann-etal . In particular, for the infinite-$`U`$ Hubbard model with spin-1/2 bosons, the total spin of the ground state is shown to be maximal for all hole densities, unlike in the electronic Hubbard model Fledderjohann-etal .
Thus, the saturated ferromagnetism could emerge in an electron system if the statistics of the electron is transmuted to bosonic. In fact, the statistical transmutation is indeed possible in two-dimensional (2D) systems Semenoff ; Fradkin , and it has been applied to fractional quantum Hall effect Girvin-MacDonald .
Combining these ideas, we can expect that the 2D electron system in the presence of an external gauge (magnetic) field exhibits the (saturated) ferromagnetism thanks to the formation of the composite boson. Namely, when the applied magnetic field amounts to unit flux quantum per electron, the magnetic-flux quantum may be assigned to an electron. The composite particle consisting of an electron and the attached flux is then expected to have spin $`1/2`$ and to obey the Bose statistics with hard-core constraint. In this way, in the mean-field level, the original system can be mapped into a spin-1/2 boson system without magnetic field, which exhibits the saturated ferromagnetism Suto ; Eisenberg-Lieb ; Fledderjohann-etal . However, as the โflux attachmentโ argument is not rigorous, whether this mechanism actually leads to the ferromagnetism in an electron system has to be checked.
In quantum Hall systems in the continuum, the fully spin-polarized ground state is favored for the filling factor $`\nu =1`$ (and in general, for $`\nu =1/m`$ with $`m`$ odd), without the Zeeman energy Zhang-Chakraborty ; Rezayi ; Maksym . This is referred to as quantum Hall ferromagnets MacDonald-etal . While this ferromagnetism is usually associated with the antisymmetric nature of the orbital part of the wavefunction, it could also be regarded as a consequence of the formation of a spin-1/2 boson which is composed of an electron and $`m`$ flux quanta MacDonald-etal . On the other hand, the dispersion in quantum Hall systems in the continuous space is completely flat (Landau levels). Thus the saturated ferromagnetism may also be understood as a special case of the flat-band ferromagnetism comment2 . It is not clear whether the statistical transmutation is essential to realize the saturated ferromagnetism in quantum Hall systems.
In order to clarify this issue, in this Letter we study an electron system on a lattice. We demonstrate an example of saturated ferromagnetism which is due entirely to statistical transmutation and is distinct from the flat-band variety.
Let us introduce the $`U\mathrm{}`$ Hubbard model on a square lattice, with the gauge (magnetic) flux $`\varphi `$ per plaquette Schofield-etal :
$``$ $`=`$ $`{\displaystyle \underset{ij\sigma }{}}[t_{ij}(\varphi _{ij})c_{i\sigma }^{}c_{j\sigma }+\mathrm{H}.\mathrm{c}.]+U{\displaystyle \underset{i}{}}n_in_i,`$
$`t_{ij}(\varphi _{ij})t\mathrm{exp}(\mathrm{i2}\pi \varphi _{ij}/\varphi _0),`$
$`\varphi ={\displaystyle \underset{\mathrm{oriented}\mathrm{plaquette}}{}}\varphi _{ij},\varphi _0h/e1,`$
where $`ij`$ refers to the nearest-neighbor pairs. Periodic boundary conditions are imposed in both directions, unless explicitly mentioned otherwise. In the $`U\mathrm{}`$ limit, we have
$$=\underset{ij\sigma }{}[t_{ij}(\varphi _{ij})\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+\mathrm{H}.\mathrm{c}.],$$
(2)
where $`\stackrel{~}{c}_{i\sigma }=c_{i\sigma }(1n_{i,\sigma })`$, which means that double occupancy at each site is excluded. We stress that, in our model (2), we have not included the Zeeman coupling of spins to the magnetic field. The model is therefore completely isotropic in the spin space, and the total spin is a conserved quantum number.
Exact numerical diagonalization for $`4\times 4,\sqrt{18}\times \sqrt{18}`$, and $`\sqrt{20}\times \sqrt{20}`$ clusters is employed in our study. We study the system with various values of the electron density $`n_\mathrm{e}`$ and the flux per plaquette $`\varphi `$. In particular, we need to investigate the case $`\varphi =n_\mathrm{e}`$ where the statistical transmutation to boson would occur. Under periodic boundary conditions, the total flux of the system is quantized to an integer. Thus $`\varphi `$ can take only integral multiples of $`1/N`$, where $`N`$ is the number of sites (plaquettes). In order to study all the possible values of $`\varphi `$, we need to use the string gauge Hatsugai-etal : Choosing a plaquette $`S`$ as a starting one, we draw $`N1`$ outgoing arrows (strings) from the plaquette $`S`$, so that each plaquette other than $`S`$ is the endpoint of a string. Then we set $`\varphi _{ij}`$ on a link $`ij`$ to $`\varphi ๐ฉ_{ij}`$ taking account of the orientation, where $`๐ฉ_{ij}`$ is the number of strings cutting the link $`ij`$. We have checked that the single-electron spectrum for a small system size ($`N20`$) in the string gauge approximately reproduces the Hofstadter butterfly in the thermodynamic limit Hofstadter .
The total spin $`S_{\mathrm{tot}}`$ at zero temperature can be evaluated from the expectation value of $`(๐บ_{\mathrm{tot}})^2=(_{\mathrm{}}๐บ_{\mathrm{}})^2`$ in the ground state, where $`๐บ_{\mathrm{}}`$ is the spin operator at site $`\mathrm{}`$. In Fig. 1 we show the scaled total spin $`S_{\mathrm{tot}}/S_{\mathrm{max}}`$ in the $`\varphi `$$`n_\mathrm{e}`$ plane. Here $`S_{\mathrm{max}}`$ is given by $`N_\mathrm{e}/2`$ with $`N_\mathrm{e}`$ being the number of electrons. Red regions correspond to saturated-ferromagnetic states. In addition to Nagaoka ferromagnetism in the single-hole case with $`\varphi =0`$, we find two common features irrespective of the system size. (i) Saturated ferromagnetism appears along a straight line with $`\varphi =n_\mathrm{e}`$ (or $`\varphi =1n_\mathrm{e}`$) except $`0.6<n_\mathrm{e}<\mathrm{\hspace{0.33em}0.7}`$. (ii) Saturated ferromagnetism appears in the triangular region surrounded by three straight lines: $`n_\mathrm{e}=0`$, $`\varphi =n_\mathrm{e}`$, and $`\varphi =1n_\mathrm{e}`$.
The result (i) confirms the expectation based on the statistical transmutation. Moreover, we find that, in most cases of $`\varphi =n_\mathrm{e}`$, the saturated ferromagnetism is robust against twisting the boundary condition. This is in contrast to Nagaoka ferromagnetism, where the total spin of the ground state is changed from maximum to zero as the boundary is twisted Kusakabe-Aoki . Besides, we have checked for $`N=16`$ and $`18`$ that the saturated-ferromagnetic ground state at $`\varphi =n_\mathrm{e}`$ is nondegenerate except for the trivial $`2S_{\mathrm{max}}+1`$-fold degeneracy. This is consistent with the ferromagnetism for spin-1/2 bosons Eisenberg-Lieb ; Fledderjohann-etal .
In order to distinguish the ferromagnetism due to the statistical transmutation from possible โflat-bandโ varieties, we define the spectral functions
$`D^{}(\omega )={\displaystyle \frac{1}{N}}{\displaystyle \underset{\mathrm{},n}{}}|\mathrm{\Psi }_n(N_{}1,N_{};\varphi )|c_{\mathrm{}}|\mathrm{\Psi }_0(N_{},N_{};\varphi )|^2`$
$`\times \delta [\omega +E_n(N_{}1,N_{};\varphi )E_0(N_{},N_{};\varphi )+\mu ],`$ (3)
$`D^+(\omega )={\displaystyle \frac{1}{N}}{\displaystyle \underset{\mathrm{},n}{}}|\mathrm{\Psi }_n(N_{}+1,N_{};\varphi )|c_{\mathrm{}}^{}|\mathrm{\Psi }_0(N_{},N_{};\varphi )|^2`$
$`\times \delta [\omega E_n(N_{}+1,N_{};\varphi )+E_0(N_{},N_{};\varphi )+\mu ].`$ (4)
Here $`\mu `$ is the chemical potential, and $`|\mathrm{\Psi }_n(N_{},N_{};\varphi )`$ denotes an eigenstate with energy $`E_n(N_{},N_{};\varphi )`$ in the system with $`N_{}`$ up-spins, $`N_{}`$ down-spins, and the flux $`\varphi `$. We define the index $`n`$ so that $`n=0`$ corresponds to the ground state with the given $`N_{}`$ and $`N_{}`$. In the following, we set $`N_{}=N_{}=N_\mathrm{e}/2`$ for $`N_\mathrm{e}`$ even. $`D^\pm `$ can be estimated numerically by the continued-fraction method Gagliano-Balseiro . Below, we show the results for two values of $`n_\mathrm{e}`$; $`n_\mathrm{e}=4/20`$ and $`n_\mathrm{e}=18/20`$ as representatives of the โlow electron densityโ and โhigh electron densityโ regimes, respectively.
First let us focus on the โlow electron densityโ case, $`n_\mathrm{e}=4/20`$. Figure 2(a) shows the evolution of $`D^\pm (\omega )`$ with varying $`\varphi `$. The sum of $`D^+`$ and $`D^{}`$ corresponds to the density of states (either occupied or unoccupied). Apparently it is always spread over a similar range of energy, representing the โbandwidthโ which is about $`8t`$ although there is some $`\varphi `$-dependence.
On the other hand, at this density, we find a crucial difference in the spectral function $`D^{}`$ corresponding to magnetism. Namely, $`D^{}`$ is concentrated in a narrow range of energy when the system exhibits the saturated ferromagnetism for $`\varphi =4/20,5/20,\mathrm{},9/20`$ (and $`1\varphi `$). In contrast, when the saturated ferromagnetism is absent, $`D^{}`$ is spread over a region of energy. Intuitively, the spectral function $`D^{}`$ corresponds to the density of states occupied by electrons. The narrow distribution of $`D^{}`$ compared to the โbandwidthโ indicates a variant of the narrow or nearly flat-band ferromagnetism Hlubina-etal .
In fact, there is more difference in $`D^{}`$ between the cases with and without saturated ferromagnetism, than what is visible in Fig. 2. $`D^{}`$ vanishes completely below a certain threshold ($`\omega /t=0.550,0.229,0.474`$ respectively for $`\varphi =4/20,6/20`$ and $`9/20`$) when the system exhibits the saturated ferromagnetism. In contrast, there is a continuous โshoulderโ of low intensity (invisible in Fig. 2) down to much lower energy $`\omega /t20`$ when saturated ferromagnetism is absent. This seems to be consistent again with our interpretation.
In particular, for $`\varphi =4/20(=n_\mathrm{e})`$, the spectral function $`D^{}`$ is localized within a narrow energy band below an apparent gap around the Fermi level. This appears similar to the quantum Hall ferromagnet at $`\nu =1`$ in the continuum MacDonald-etal . In this case, the statistical transmutation mechanism and the flat-band mechanism appear indistinguishable.
Now let us discuss the saturated ferromagnetism observed in the โhigh electron densityโ regime, at $`\varphi =n_\mathrm{e}=18/20`$. Figure 2(b) shows the spectral functions in this case. Clearly, the spectral function $`D^{}`$ spreads over almost the entire bandwidth even though the system does exhibit the saturated ferromagnetism. Thus the ferromagnetism at $`\varphi =n_\mathrm{e}=18/20`$ is difficult to be understood in terms of the flat-band mechanism, and appears to be exclusively due to the statistical transmutation mechanism. In fact, changing the value of $`\varphi `$ destroys the saturated ferromagnetism at this electron density, as expected from the statistical transmutation scenario.
In order to further confirm the statistical transmutation scenario at $`\varphi =n_\mathrm{e}`$, we define the following operators: $`b_\mathrm{}\sigma =\mathrm{e}^{\mathrm{i}๐ฅ_{\mathrm{}}}c_\mathrm{}\sigma `$ and $`b_\mathrm{}\sigma ^{}=c_\mathrm{}\sigma ^{}\mathrm{e}^{\mathrm{i}๐ฅ_{\mathrm{}}}`$, where $`๐ฅ_{\mathrm{}}=m_{i(\mathrm{})}\theta _\mathrm{}in_i`$ with $`n_i=_{\sigma =,}c_{i\sigma }^{}c_{i\sigma }`$. In general, $`m`$ denotes the number of magnetic-flux quanta, and we set $`m=1`$ in the present case. $`\theta _\mathrm{}i`$ is the argument of the vector drawn from site $`i`$ to site $`\mathrm{}`$. Note that the relation $`\theta _\mathrm{}i\theta _i\mathrm{}=\pm \pi `$ holds for $`i\mathrm{}`$. Then we can prove that these operators satisfy the commutation relations of bosons for two different sites. There is a large freedom in determining explicit values of $`\theta _i\mathrm{}`$. Here we follow a prescription introduced in Ref. Cabra-Rossini , although this prescription unavoidably breaks translation invariance of a periodic cluster. The order parameter for the condensation of the composite bosons may be defined by
$$O_\mathrm{B}=\frac{1}{N}\underset{\mathrm{}}{}|\mathrm{\Psi }_0(N_{}1,N_{};\varphi \frac{1}{N})|b_{\mathrm{}}|\mathrm{\Psi }_0(N_{},N_{};\varphi )|^2.$$
(5)
We again set $`N_{}=N_{}=N_\mathrm{e}/2`$ for $`N_\mathrm{e}`$ even, and $`N_{}=N_{}+1=(N_\mathrm{e}+1)/2`$ for $`N_\mathrm{e}`$ odd. Figure 3(a) shows $`O_\mathrm{B}`$ as a function of $`n_\mathrm{e}(=\varphi )`$. The order parameter has a pronounced enhancement in both the high-density region ($`0.7<n_\mathrm{e}<1`$) and the low-density one ($`0<n_\mathrm{e}<\mathrm{\hspace{0.33em}0.3}`$). Figures 3(b) and 3(c) depict the $`\varphi `$-dependence of $`O_\mathrm{B}`$ for $`n_\mathrm{e}=6/20`$ and $`16/20`$, respectively. We find a salient growth at $`\varphi =n_\mathrm{e}`$ for both densities. We have also observed for $`n_\mathrm{e}=6/20`$ that the order parameter given by Eq. (5) with $`m=1`$ has a peak at $`\varphi =1n_\mathrm{e}`$. These results again support the ferromagnetism based on the statistical transmutation.
Finally, we discuss the region $`0.6<n_\mathrm{e}<\mathrm{\hspace{0.33em}0.7}`$, where the saturated ferromagnetism is absent in spite of $`\varphi =n_\mathrm{e}`$. As a candidate of the competing order, we consider the spin chirality Wen-Wilczek-Zee defined by the order parameter $`\chi _{\mathrm{ch}}=(1/N)_{\mathrm{}}๐บ_{\mathrm{}}๐บ_{\mathrm{}+\widehat{y}}\times ๐บ_{\mathrm{}+\widehat{x}}`$, where $`\widehat{x}`$ ($`\widehat{y}`$) is the unit vector along $`x`$ ($`y`$) direction, and $`\mathrm{}`$ denotes the expectation value in the ground state. In fact, Nagaoka ferromagnetism in the single-hole case is known to be destroyed by development of the spin chirality in the presence of a perpendicular magnetic field Schofield-etal . Along the line $`\varphi =n_\mathrm{e}`$, we have confirmed that the spin chirality vanishes when the system exhibits saturated ferromagnetism. On the other hand, the chiral order is indeed developed when the saturated ferromagnetism is absent (not shown).
In summary, we have calculated the total spin of the ground state in the $`U\mathrm{}`$ Hubbard model with magnetic flux ($`\varphi `$) perpendicular to a square lattice and revealed regions of saturated ferromagnetism. The saturated ferromagnetism at $`\varphi =n_\mathrm{e}`$ is argued to be due to formation of spinful composite bosons. Statistical transmutation may therefore play a key role in ferromagnetism in strongly correlated systems, just as it did in fractional quantum Hall effect.
The present mechanism may be relevant to future experiments on an artificial crystal of a square lattice with quantum dots (i.e., a quantum dot superlattice) Kimura-etal . A large lattice constant of such a crystal would enable us to observe the magnetic-field effect at a modest magnetic field of a few tesla. In this situation, the orbital motion rather than the Zeeman effect could be essential for the emergence of ferromagnetism. Another possibility is to induce the effective gauge field internally without an applied magnetic field Taguchi-etal .
We thank D.S. Hirashima and M. Koshino for valuable discussions, and I. Herbut for critical reading of the manuscript. Y.S. was supported in part by JSPS Research Fellowships for Young Scientists. The numerical calculations were performed partly at the Supercomputer Center of the Institute for Solid State Physics, University of Tokyo. The present work is supported in part by Grant-in-Aid for Scientific Research and 21st Century COE programs โNanometer-Scale Quantum Physicsโ at Tokyo Institute of Technology and โORIUMโ at Nagoya University, all from MEXT of Japan. |
warning/0506/quant-ph0506008.html | ar5iv | text | # Squeezing and amplitude-squared squeezing effects on the dynamics of two nonidentical two-level atoms
Electronic address:\]bash@ssu.samara.ru
## Abstract
Squeezing and amplitude-squared squeezing for two two-level nonidentical atoms in lossless cavity has been investigated assuming the field to be initially in the coherent state. The time-dependent squeezing parameters has been calculated. The influence of the relative differences of two coupling constants on the squeezing parameters has been analyzed.
Squeezing phenomena attract much attention over the last few decades. The squeezed states of light were investigated intensively both from theoretical and experimental point of view Scully and attract considerable attention because their possible practical applications for high-precision optical measurements, optical communications and optical processing Yamamoto . A variety of schemes for producing squeezed states has been proposed.
The possibility of squeezing phenomenon in Jaynes-Cummings model (JCM) was analyzed by several authors starting with the Meystre and Zubairy Meystre . The multiphoton, nondegenerate two-mode and two-atom generalizations of JCM have also been shown to produce squeezing Kien ,Kien1 . The field squeezing in the two-atom JCM with one and multiphoton transitions has been investigated in several papers for initial coherent, squeezed, vacuum and thermal field input Kien1 ,Mir . Last years some interest has arisen in higher-order squeezing Hong . One type of higher-order squeezing, namely, squeezing of the square of the field amplitude or in brief the amplitude-squared squeezing (ASS) has been proposed by Hillery H . The ASS has been shown to exist in one- and multiphoton JCM Yang and two-atom JCMMir ,B .
In recent years, the model with two nonidentical two-level atoms in cavity has attracted a considerable attention in the study of the collective atom-field interaction. The exact solution of this model for lossless cavity and the field which is at resonance with the atomic transitions has been calculated firstly for one-photon transitions by Zubairy et al. Z , for two-photon transitions by Jex J and for $`m`$-photon transitions by Xu et al. X . Based on these solutions both the collapse-revival phenomenon of the atomic coherence for initial coherent Z , binomial S and squeezed field state X and the photon statistics X ,A have been considered. The entanglement of two nonidentical atoms, interacting with the thermal field in the cavity with loss has been studied in Zh . Agarwal and co-authors have investigated the two-photon absorption K and large two-photon vacuum Rabi oscillations P in system of two nonidentical atoms taking into account the detuting. In this paper we consider the squeezing and ASS in the system of the two atoms with different coupling constants which interacts with one mode of coherent field in lossless cavity. We analyse the dependence of the squeezing on the relative difference of two coupling constants.
Let us consider a system of two nonidentical two-level atoms interacting with a single-mode quantized electromagnetic field in a lossless resonant cavity via the one-photon-transition mechanism. The Hamiltonian of the considered system in the rotating wave approximation is
$$H=\mathrm{}\omega a^+a+\underset{i=1}{\overset{2}{}}\mathrm{}\omega _0R_i^z+\underset{i=1}{\overset{2}{}}\mathrm{}g_i(R_i^+a+R_i^{}a^+),$$
(1)
where $`a^+`$ and $`a`$ are the creation and annihilation operators of photons of the cavity field, respectively, $`R_f^+`$ and $`R_f^{}`$ are the raising and the lowering operators for the $`i`$th atom, $`\omega `$ and $`\omega _0`$ are the frequencies of the field mode and the atoms, $`g_i`$ is the coupling constant between the $`i`$th atom and the field. We assume the field to be at one-photon resonance with the atomic transition, i.e. $`\omega _0=\omega `$ .
We denote by $`+`$ and $``$ the excited and ground states of single atom and by $`n`$ the Fock state of the electromagnetic field. The two-atom wave function can be expressed as a combination of state vectors of the form $`v_\mathit{1},v_\mathit{2}=v_\mathit{1}v_\mathit{2}`$, where $`v_\mathit{1},v_\mathit{2}=+,`$. Let the atoms are initially in the ground state $`,`$ and the field is initially in a coherent state $`\alpha `$,
$$\alpha =\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{exp}\left(\frac{\alpha ^2}{2}\right)\frac{\alpha ^n}{\sqrt{n!}},$$
where $`\alpha =\alpha e^{ฤฑ\phi }`$ and $`\overline{n}=\alpha ^2`$ is the initial mean photon number or dimensionless intensity of the cavity field.
The time-dependent wave function of the total system $`\mathrm{\Psi }(t)`$ obeys the $`\mathrm{Schr}\ddot{\mathrm{o}}\mathrm{dinger}`$ equation
$$ฤฑ\mathrm{}\dot{\mathrm{\Psi }}(t)=H\mathrm{\Psi }(t).$$
$`(2)`$
Using the Hamiltonian (1) the wave function is found to be
$$\mathrm{\Psi }(t)=_{n=0}^{\mathrm{}}\mathrm{exp}[ฤฑ(n1)\omega t]\mathrm{exp}(\frac{\alpha ^2}{2})\frac{\alpha ^n}{\sqrt{n!}}\times $$
$$\times [C_1^{(n)}(t)+,+;n2+C_2^{(n)}(t)+,;n1+C_3^{(n)}(t),+;n1+C_4^{(n)}(t),;n].$$
$`(3)`$
With the help of formulas (1)-(3) we can obtain the equations of motion for probability coefficients $`C_i^n(t)`$. These equations must be written separately for $`n=0,n=1`$ and $`n2`$:
$$\dot{C}_i^{(0)}=0(i=1,2,3,4);$$
$`(4)`$
$$\dot{C}_1^{(1)}=0,\dot{C}_2^{(1)}=ฤฑg_1C_4^{(1)},\dot{C}_3^{(1)}=ฤฑg_2C_4^{(1)},\dot{C}_4^{(1)}=ฤฑ(g_1C_2^{(1)}+g_2C_3^{(1)});$$
$`(5)`$
and for $`n2`$
$$\dot{C}_1^{(n)}=ฤฑ(g_2\sqrt{n1}C_2^{(n)}+g_1\sqrt{n1}C_3^{(n)},$$
$$\dot{C}_2^{(n)}=ฤฑ(g_2\sqrt{n1}C_1^{(n)}+g_1\sqrt{n}C_4^{(n)},$$
$$\dot{C}_3^{(n)}=ฤฑ(g_1\sqrt{n1}C_1^{(n)}+g_2\sqrt{n}C_4^{(n)},$$
$$\dot{C}_4^{(n)}=ฤฑ(g_1\sqrt{n}C_2^{(n)}+g_2\sqrt{n}C_3^{(n)}.$$
$`(6)`$
For atoms initially prepared in their ground state we have the initial conditions for probability coefficients
$$C_4^{(n)}(0)=1,C_1^{(n)}(0)=C_2^{n)}(0)=C_3^{(n)}(0)=0(n=0,1,2,\mathrm{}).$$
$`(7)`$
The solutions of Eqs. (4)-(6) with initial conditions (7) are found to be
$$C_1^{(0)}(t)=C_2^{(0)}(t)=C_3^{(0)}(t)=0,C_4^{(0)}(t)=1;$$
$`(8)`$
$$C_1^{(1)}(t)=0,C_2^{(1)}(t)=\frac{ฤฑ\mathrm{sin}(\sqrt{1+R^2}t)}{\sqrt{1+R^2}},$$
$$C_3^{(1)}(t)=\frac{ฤฑR\mathrm{sin}(\sqrt{1+R^2}t)}{\sqrt{1+R^2}},C_4^{(1)}(t)=\mathrm{cos}(\sqrt{1+R^2}t)$$
$`(9)`$
and for $`n2`$
$$C_1^{(n)}(t)=\frac{2R\sqrt{(n1)n}}{\beta }[\mathrm{cos}(\lambda _+t)\mathrm{cos}(\lambda _{}t)],$$
$$C_2^{(n)}(t)=\frac{4ฤฑR^2(n1)\sqrt{n}}{\beta }\left\{\frac{\lambda _+^2+(1R^2)n}{\lambda _+[\beta (1+R^2)]}\mathrm{sin}(\lambda _+t)\frac{\lambda _{}^2+(1R^2)n}{\lambda _{}[\beta +(1+R^2)]}\mathrm{sin}(\lambda _{}t)\right\},$$
$$C_2^{(n)}(t)=\frac{4ฤฑR(n1)\sqrt{n}}{\beta }\left\{\frac{\lambda _+^2(1R^2)n}{\lambda _+[\beta (1+R^2)]}\mathrm{sin}(\lambda _+t)\frac{\lambda _{}^2(1R^2)n}{\lambda _{}[\beta +(1+R^2)]}\mathrm{sin}(\lambda _{}t)\right\},$$
$$C_4^{(n)}(t)=\frac{8R^2(n1)n}{\beta }\left[\frac{\mathrm{cos}(\lambda _+t)}{\beta (1+R^2)}+\frac{\mathrm{cos}(\lambda _{}t)}{\beta +(1+R^2)}\right],$$
$`(10)`$
where
$$\lambda _\pm =\sqrt{(1+R^2)(2n1)\pm \beta }/\sqrt{2},$$
$$\beta =\sqrt{(2n1)^2(1+R^2)^24(n1)n(1R^2)^2},R=g_2/g_1.$$
To investigate the photon squeezing we introduce the two slowly varying quadrature components $`X_1,X_2`$ of field, defined by
$$X_1=\frac{1}{2}(ae^{ฤฑ\omega t}+a^+e^{ฤฑ\omega t}),$$
$$X_2=\frac{1}{2ฤฑ}(ae^{ฤฑ\omega t}a^+e^{ฤฑ\omega t}).$$
Thus $`[X_1,X_2]=ฤฑ/2`$, which implies the uncertainty relation $`(\mathrm{\Delta }X_1)^2(\mathrm{\Delta }X_2)^21/16`$, where $`(\mathrm{\Delta }X_i)^2=X_i^2X_i^2(i=1,2)`$ are variances of quadrature components. Normal squeezing occurs when variances satisfy the relation $`(\mathrm{\Delta }X_i)^2<1/4(i=1\mathrm{or}\mathrm{\hspace{0.17em}2})`$. The condition for squeezing one can write in the form $`S_i<0`$, where squeezing parameters are
$`S_i={\displaystyle \frac{(\mathrm{\Delta }X_i)^2)1/4}{1/4}}=4(\mathrm{\Delta }X_i)^21(i=1,2).`$
The value $`S_i=1`$ corresponds to 100% squeezing in $`i`$th quadrature component. In terms of photon creation and annihilation operators we can rewrite squeezing parameters in the form
$$S_1=2a^+a+2Rea^2e^{2ฤฑ\omega t}4(Reae^{ฤฑ\omega t})^2,$$
$`(11)`$
$$S_2=2a^+a2Rea^2e^{2ฤฑ\omega t}4(Imae^{ฤฑ\omega t})^2.$$
$`(12)`$
Using (3) we can obtain
$$a^+a=\overline{n}\left[2\underset{n=2}{\overset{\mathrm{}}{}}p_nC_1^{(n)}^2+\underset{n=1}{\overset{\mathrm{}}{}}p_n(C_2^{(n)}^2+C_3^{(n)}^2)\right]=A_0,$$
$$e^{ฤฑ\omega t}a=\alpha \{\underset{n=2}{\overset{\mathrm{}}{}}p_n(C_1^{(n)})^{}C_1^{(n+1)}\sqrt{\frac{n1}{n+1}}+\underset{n=1}{\overset{\mathrm{}}{}}p_n[(C_2^{(n)})^{}C_2^{(n+1)}+$$
$$+(C_3^{(n)})^{}C_3^{(n+1)}]\sqrt{\frac{n}{n+1}}+\underset{n=0}{\overset{\mathrm{}}{}}p_n(C_4^{(n)})^{}C_4^{(n+1)}\}=\alpha A_1,$$
$$e^{2ฤฑ\omega t}a^2=\alpha ^2\{\underset{n=2}{\overset{\mathrm{}}{}}p_n(C_1^{(n)})^{}C_1^{(n+2)}\sqrt{\frac{(n1)n}{(n+1)(n+2)}}+\underset{n=1}{\overset{\mathrm{}}{}}p_n[(C_2^{(n)})^{}C_2^{(n+2)}+$$
$$\text{+}(C_3^{(n)})^{}C_3^{(n+2)}]\sqrt{\frac{n}{n+2}}+\underset{n=0}{\overset{\mathrm{}}{}}p_n(C_4^{(n)})^{}C_4^{(n+2)}\}=\alpha ^2A_2.$$
$`(13)`$
The parameter of initial coherent state is $`\alpha =\sqrt{\overline{n}}\mathrm{exp}i\phi .`$ Let below $`\phi =0`$. Then, for squeezing parameters $`S_1`$ and $`S_2`$ one can write
$$S_1=2A_0+2\overline{n}A_24\overline{n}A_1^2,$$
$`(14)`$
$$S_2=2A_02\overline{n}A_2.$$
$`(15)`$
To define the squeezing of the square of the field amplitude or amplitude-squared squeezing (ASS) we can introduce the quantities H
$`Y_1={\displaystyle \frac{1}{2}}(a^2e^{2ฤฑ\omega t}+a^{+2}e^{2ฤฑ\omega t}),`$
$`Y_2={\displaystyle \frac{1}{2ฤฑ}}(a^2e^{2ฤฑ\omega t}a^{+2}e^{2ฤฑ\omega t}).`$
The operators $`Y_1`$ and $`Y_2`$ correspond to the real and imaginary parts, respectively, of the field amplitude squared and obey the commutation relation $`[Y_1,Y_2]=i(2n+1)`$, where $`n=a^+a`$. The uncertainty relation for these two quantities has the form
$`(\mathrm{\Delta }Y_1)^2(\mathrm{\Delta }Y_2)^2n+1/2^2.`$
The ASS state in $`Y_1`$ exists if $`(\mathrm{\Delta }Y_2)^2<n+1/2`$ and similarly for $`Y_2`$. Then, we can introduce the squeezing parameters for ASS in the following form
$$Q_i=\frac{(\mathrm{\Delta }Y_i)^2n+1/2}{n+1/2}=n+1/2^1((\mathrm{\Delta }Y_i)^21.$$
The SSFA is obtained whenever $`Q_i<0`$ for $`i=1`$ or $`i=2`$ and $`Q_i=1`$ will correspond to 100% SSFA. In terms of photon creation and annihilation operators we can rewrite SSFA squeezing parameters in the form B
$$Q_1=\frac{1}{4}n+1/2^1\left[2a^{+2}a^2+2Rea^4e^{4ฤฑ\omega t}4(Rea^2e^{2ฤฑ\omega t})^2\right],$$
$`(16)`$
$$Q_2=\frac{1}{4}n+1/2^1\left[2a^{+2}a^22Rea^4e^{4ฤฑ\omega t}4(Ima^2e^{2ฤฑ\omega t})^2\right].$$
$`(17)`$
From (3) we have
$$a^{+2}a^2=\underset{n=4}{\overset{\mathrm{}}{}}p_n(n2)(n3)C_1^{(n)}^2++\underset{n=3}{\overset{\mathrm{}}{}}p_n(n1)(n2)[C_2^{(n)}^2+C_3^{(n)}^2]+$$
$$+\underset{n=2}{\overset{\mathrm{}}{}}p_nC_4^{(n)}^2=A_3,$$
$$e^{4ฤฑ\omega t}a^2=\alpha ^2\{\underset{n=2}{\overset{\mathrm{}}{}}p_n(C_1^{(n)})^{}C_1^{(n+4)}\sqrt{\frac{(n1)n}{(n+3)(n+4)}}+\underset{n=1}{\overset{\mathrm{}}{}}p_n[(C_2^{(n)})^{}C_2^{(n+4)}+$$
$$\text{+}(C_3^{(n)})^{}C_3^{(n+4)}]\sqrt{\frac{n}{n+4}}+\underset{n=0}{\overset{\mathrm{}}{}}p_n(C_4^{(n)})^{}C_4^{(n+24)}\}=\alpha ^24A_4.$$
$`(18)`$
With taking into account the Eqs. (13),(16)-(18) we can rewrite the ASS parameters $`Q_1`$ and $`Q_2`$ in the form
$$Q_1=\frac{1}{4}n+1/2^1\left[2A_3+2\overline{n}^2A_44\overline{n}^2A_2^2\right],$$
$`(19)`$
$$Q_2=\frac{1}{4}n+1/2^1\left[2A_32\overline{n}^2A_4\right].$$
$`(20)`$
Using the expressions (11)-(20) we have calculated the squeezing parameters $`S_i`$ and $`Q_i`$ for various initial photon numbers $`\overline{n}`$ and relative differences of two coupling constant $`R`$.
Fig. 1 presents the long time behaviour of parameters $`S_1`$ and $`S_2`$ for $`\overline{n}=0.2`$ and $`R=0.5`$. For small field intensities $`\overline{n}`$ as soon as $`t>0`$ we observe negative values of $`S_1`$ (squeezing in the first field quadrature component) and positive values of $`S_2`$. As times goes on, $`S_1`$ and $`S_2`$ start oscillating and reversing sign. The maximum degree of subsequent squeezing may be larger than that for the first squeezing. These features have much in common with that for the case of single or two identical atoms Meystre ,Kien1 . With increasing of $`\overline{n}`$ the degree of squeezing in $`S_1`$ and the number of squeezing intervals decreases.
Figs. 2-5 present the short time behaviour of squeezing parameter $`S_1`$ (the first squeezing) for different small field input intensities $`\overline{n}`$ and values of relative differences of two coupling constants. Obviously, that for case $`R=0`$ we have dealings with a single atom and the case $`R=1`$ corresponds to two identical atoms. For small input intensity $`\overline{n}`$ (letโs say $`0\overline{n}0.3`$) the degree of first squeezing increases with decreasing of $`R`$ (as $`R`$ decreases from 1 to 0 the maximum obtainable degree of squeezing increases from 20% to 27% for $`\overline{n}=0.2`$). For field intensities $`\overline{n}0.3`$ the maximum degree of squeezing is insensitive to choice of $`R`$. But for larger intensity input (lets say $`\overline{n}>0.3`$) the dependence of the degree of squeezing from $`R`$ is reversed. When, for instance, $`\overline{n}=0.4`$ the increasing of $`R`$ from 0 to 1 leads to increasing the degree of squeezing from 18% to 28%. Note that at the beginning of time scale the squeezing parameter $`S_1`$ for model with two nonidentical atoms takes the positive values in contrast to that for single atom or two identical atoms and the first squeezing of $`S_1`$ is reached with some delay time. But this features is distinct only for relative large initial intensities. In Fig.4 we show the short time behaviour of squeezing parameters $`S_1`$ for models with $`\overline{n}=0.8`$ and different $`R`$. For $`\overline{n}>0.8`$ the $`R`$ -dependence of the degree of squeezing has nonmonotone character. Note that for large input intensities, the parameter $`S_1`$ exhibits weak first squeezing and with increasing $`\overline{n}`$ the squeezing is vanished at first for intermediate values of $`R`$ (See Fig. 5).
Fig. 6 presents the long time behaviour of ASS parameters $`Q_1`$ and $`Q_2`$ for $`\overline{n}=0.2`$ and $`R=0.5`$. These parameters for small input intensity parameters are carried out in much the same way as $`S_1`$ and $`S_2`$ but the amount of squeezing for ASS is less than that for second-order squeezing. The maximum degree of squeezing in $`Q_1`$ decreases with increasing of the parameter $`R`$. The dependence $`Q_1`$ and $`Q_2`$ from intensity $`\overline{n}`$ have the more complicated character but for large intensities $`\overline{n}`$ the ASS is weak in both components.
Figs. 7,8 present the short time behaviour of squeezing parameter $`Q_1`$ (the first ASS) for different field intensities $`\overline{n}`$ and different values of relative differences of two coupling constants $`R`$. For small input intensity $`\overline{n}`$ (letโs say $`0\overline{n}0.7`$) the degree of first ASS increases with decreasing of $`R`$ (as $`R`$ decreases from 1 to 0 the maximum obtainable degree of squeezing increases from 5% to 1.5% for $`\overline{n}=0.4`$). For $`\overline{n}>0.7`$ the $`R`$ -dependence of the degree of squeezing has nonmonotone character. In particular for model with $`\overline{n}=0.8`$ the maximum of ASS is equal 6% when $`R=0.5`$. Similarly to ordinary squeezing the first ASS is appeared with some delay time when $`0<R<1`$ and with increasing of the input intensity the ASS is vanished at first for intermediate values of $`R`$.
Thus, we have considered the effects of squeezing and amplitude-squared squeezing of the cavity field mode in the model with two nonidentical atoms. The case in which the field is initially in a coherent state together with the atoms in the ground state has been examined. The long and short time behaviour of the squeezing and ASS parameters have been calculated. The influence of the relative differences of two coupling constants on the squeezing parameters has been analyzed. The investigation of the model with multiphoton transitions and other initial states of field and atoms is the aim of our subsequent papers.
This work was supported by RFBR grant 04-02-16932. |
warning/0506/cond-mat0506178.html | ar5iv | text | # Conformational transformations induced by the charge-curvature interaction
## I Introduction
The functioning of biological macromolecules is determined by their tertiary structure and different types of semiflexible polymers in solutions exhibit a variety of conformational phase transitions geom . Even modest conformational changes modify long-range electronic interactions in oligopeptides wolfgang , they may remove steric hindrances and open the pathways for molecular motions which are not available in rigid proteins feitel . In particular, it has been recently shown viduna that flexibility increases the hydrogen accessibility of DNA fragments and in this way facilitates strand breaks in DNA molecules. There is also a strong belief that the conductivity of DNA is due to thermal motion of small polarons conw ; yoo ; alex .
Research in solitonic properties of the chains with a bending has been initiated in the recent years peyrard ; polyakova ; curv1 ; arch ; george . In particular, it was shown that the bending of the chain could manifest itself as an effective trap for nonlinear excitations peyrard ; curv1 ; george and that the energy of excitations decreases when the curvature of the bending increases curv1 .
Quite recently, there has been a growing interest in studying nonlinear charge and energy transport in soft condensed systems (polymers, membranes) with self-consistent account of coupling between nonlinear excitations and the shape of the systems. It has been demonstrated saxena-1998 that a mismatch of length scales in the presence of magnetic solitons leads to an elastic deformation on a soft magnetic surface. A phenomenological model for describing the conformational dynamics of biopolymers via the nonlinearity-induced buckling and collapse instability was proposed in ming02 . It was shown there that the nonlinear excitations may cause local softening of polymer bonds. That is the effective bending rigidity of a chain may become negative nearby the nonlinear excitation and in this way the nonlinear excitation causes a buckling instability of the chain.
In this paper we study a simple model for electron-curvature interactions on closed molecular aggregates. In particular we show that due to the interaction between electrons and the bending degrees of freedom the circular shape of the aggregate may be become unstable and the aggregate takes the shape of an ellipse or a polygon. It is shown that the interaction between complexes may stabilize the ringlike shape of the The paper is organized as follows. The theory is applied to study the conformational transitions in light-harvesting complexes in purple bacteria. In Sec. II we describe a model. In Sec. III we present an analytical solution to the Euler-Lagrange equation. In Sec.III we compare our analytical results to results obtained directly by numerical simulations. In Sec. IV the interplay between shape of complexes and intercomplex interaction is studied. Sec. V is devoted to application to conformational transformations in light-harvesting complexes. Sec. VI presents some concluding remarks.
## II The model
Let us consider a chain consisting of L units labelled by an index $`n`$, and located at the points $`\stackrel{}{r}_n=\{x_n,y_n,z_n\}`$. We are interested in the case when the chain is closed and so we impose the periodicity condition on the coordinates $`\stackrel{}{r}_n`$
$`\stackrel{}{r}_n=\stackrel{}{r}_{n+L}.`$ (1)
The chain flexibility is accounted for by employing a microstructure consisting of many sequentially joined rigid rods and by incorporating a bend potential at each point of rotation chaubal . Thus the Hamiltonian of such a polymer chain has the form
$`H=U+H_{el},`$ (2)
with the potential energy of inter-unit interactions $`U(\mathrm{},\stackrel{}{r}_n,\stackrel{}{r}_{n+1},..)`$ which we take in the form
$`U=U_S+U_B`$ (3)
where
$`U_s={\displaystyle \frac{\sigma }{2}}{\displaystyle \underset{n}{}}\left(|\stackrel{}{r}_n\stackrel{}{r}_{n+1}|a\right)^2`$ (4)
is the stretching energy in the harmonic approximation. Here $`a`$ is an equilibrium distance between units (in what follows we assume $`a=1`$) and $`\sigma `$ is a dimensionless elastic modulus of the stretching rigidity of the chain, and
$`U_b={\displaystyle \frac{k}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{\kappa _n^2}{1\kappa _n^2/\kappa _{max}^2}}`$ (5)
is the bending energy. Furthermore,
$`\kappa _n|\stackrel{}{\tau }_n\stackrel{}{\tau }_{n1}|=2\mathrm{sin}{\displaystyle \frac{\alpha _n}{2}}`$ (6)
determines the curvature of the chain at the point $`n`$. The vector
$`\stackrel{}{\tau }_n={\displaystyle \frac{\stackrel{}{r}_{n+1}\stackrel{}{r}_n}{|\stackrel{}{r}_{n+1}\stackrel{}{r}_n|}}`$ (7)
is the unit tangent vector, $`\alpha _n`$ is the angle between the tangent vectors $`\stackrel{}{\tau }_n`$ and $`\stackrel{}{\tau }_{n1}`$, $`\kappa _{max}=2\mathrm{sin}\left(\alpha _{max}/2\right)`$ with $`\alpha _{max}`$ being the maximum bending angle, $`k`$ is the elastic modulus of the bending rigidity of the chain.
We assume that there is a small amount of extra electrons ( or holes) on the chain. The Hamiltonian of electrons and their interaction with conformational degrees of freedom is
$`H_{el}=H_{hop}+H_{elconf},`$ (8)
where
$`H_{hop}=J{\displaystyle \underset{n}{}}\left|\psi _n\psi _{n+1}\right|^2,`$ (9)
describes the motion of electrons along the chain, and
$`H_{elconf}={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\chi \left(|\psi _{n+1}|^2+|\psi _{n1}|^2\right)\kappa _n^2`$ (10)
gives the interaction between electrons and bending degrees of freedom. In Eqs. (8)-(10) $`\psi _n(t)`$ is the electron wave function at the site $`n`$, the parameter $`J`$ describes the electron hopping in the chain, and $`\chi `$ is the curvature-electron coupling constant (see Appendix for details).
The quantity
$`\nu {\displaystyle \frac{1}{L}}{\displaystyle \underset{n}{}}|\psi _n|^2`$ (11)
gives the total density of extra electrons which can move along the chain and participate in the formation of the conformational state of the system. We will neglect the interaction between electrons. It is legitimate when the total density of electrons in the chain $`\nu `$ is small. Combining Eqs. (5) and (10), we notice that the effective bending rigidity changes close the points where the electron is localised. For positive values of the coupling constant $`\chi `$ there is a local softening of the chain, while for $`\chi `$ negative there is a local hardening of the chain.
In what follows we assume that the chain is planar ($`z_n=0`$) and inextensible ($`\sigma \mathrm{}`$):
$`|\stackrel{}{r}_n\stackrel{}{r}_{n+1}|=1.`$ (12)
## III Continuum approach
We are interested here in the case when the characteristic size of the excitation is much larger than the lattice spacing. This permits us to replace $`\psi _n(t)`$ by the function $`\psi (s,t)`$ of the arclength $`s`$ which is the continuum analogue of $`n`$. Using the Euler-Mclaurin summation formula abr we get
$`H=U_b+H_{el},`$ (13)
$`U_b={\displaystyle \frac{1}{2}}k{\displaystyle \underset{0}{\overset{L}{}}}\kappa (s)^2๐s`$ (14)
$`H_{el}={\displaystyle \underset{0}{\overset{L}{}}}\left\{J\left|_s\psi \right|^2\chi \kappa (s)^2|\psi |^2\right\}๐s.`$ (15)
Being interested here in small curvature effects: $`\kappa (s)\kappa _{max}`$, we consider the bending energy in the harmonic approximation given by Eq. (14).
### III.1 Ground state of the chain
The electron ground state wave function $`\varphi (s)`$ and the shape of the chain $`\stackrel{}{r}(s)`$ may be obtained by minimizing the functional
$`=H_{el}+U_b`$ (16)
with $`H_{el}`$ and $`U_b`$ given by Eqs (15) and (14) under the constraint
$`\nu ={\displaystyle \frac{1}{L}}{\displaystyle \underset{0}{\overset{L}{}}}\varphi (s)^2๐s,`$ (17)
which is a continuum analog of Eq. (11). The inextensibility constraint (12) reads in the continuum limit
$`|_s\stackrel{}{r}|^2=1.`$ (18)
The inextensibility constraint (18) is automatically taken into account by choosing the parametrization
$`_sx(s)=\mathrm{sin}\theta (s),_sy(s)=\mathrm{cos}\theta (s)`$ (19)
where the angle $`\theta (s)`$ satisfies the conditions
$`\theta (s+L)=2\pi +\theta (s)`$ (20)
and
$`{\displaystyle \underset{0}{\overset{L}{}}}\mathrm{cos}\theta (s)๐s={\displaystyle \underset{0}{\overset{L}{}}}\mathrm{sin}\theta (s)๐s=0`$ (21)
which follow from Eq. (1). Note that, in the continuum limit, the curvature of the chain $`\kappa (s)`$ given by Eq. (6) can be expressed as
$`\kappa (s)=|_s^2\stackrel{}{r}(s)|.`$ (22)
Thus, in the frame of the parametrization (19), $`\kappa (s)=_s\theta `$ and the functional (16) takes the form
$`={\displaystyle \underset{0}{\overset{L}{}}}\left\{J\nu \left(_s\phi \right)^2+\left({\displaystyle \frac{k}{2}}\chi \nu \phi ^2\right)\left(_s\theta \right)^2\right\}๐s,`$ (23)
where the rescaled function $`\phi (s)=\sqrt{\nu }\varphi (s)`$ which satisfies the normalization condition
$`{\displaystyle \frac{1}{L}}{\displaystyle \underset{0}{\overset{L}{}}}\phi ^2(s)=1,`$ (24)
has been introduced.
The Euler-Lagrange equations for the problem of minimizing $``$, given by Eq. (23) under the constraint (24) become
$`_s^2\phi +{\displaystyle \frac{\chi }{J}}\left(_s\theta \right)^2\varphi \lambda \phi =0,`$ (25)
$`_s\left(_s\theta \left(1w\phi ^2\right)\right)=0`$ (26)
where $`\lambda `$ is the Lagrange multiplier and
$`w={\displaystyle \frac{2\chi \nu }{k}}`$ (27)
is a coupling constant which characterizes the strength of the charge-curvature interaction in terms of the bending rigidity of the chain and the charge density. We are interested in solutions of Eq. (25) subject to the periodic boundary conditions
$`\phi (s)=\phi (s+{\displaystyle \frac{L}{n}}).`$ (28)
where $`n`$ is an integer which characterizes the shape of the chain (see below). Integrating Eq. (26), we get
$`_s\theta ={\displaystyle \frac{A}{1w\phi ^2}}`$ (29)
where $`A`$ is an integration constant. Taking into account the condition (20) we obtain that the integration constant $`A`$ is determined by the relation
$`A={\displaystyle \frac{2\pi }{L}}I`$ (30)
where the functional $`I`$ is given by the relation
$`{\displaystyle \frac{1}{I}}={\displaystyle \frac{1}{L}}{\displaystyle \underset{0}{\overset{L}{}}}{\displaystyle \frac{ds}{1w\phi ^2}}.`$ (31)
From Eqs. (19) and (29) we see that the shape of the chain is determined by the equations
$`x(s)={\displaystyle \underset{0}{\overset{s}{}}}sin\theta (s^{})๐s^{},y(s)={\displaystyle \underset{0}{\overset{s}{}}}cos\theta (s^{})๐s^{},`$
$`\theta (s)={\displaystyle \frac{2\pi }{LI}}{\displaystyle \underset{0}{\overset{s}{}}}{\displaystyle \frac{1}{1w\phi ^2(s^{})}}๐s^{}.`$ (32)
### III.2 Solution of the Euler-Lagrange equations
There are two kinds of solutions to Eqs (25) and (29).
* Circular chain.
Charge is uniformly distributed along the chain
$`\phi =1`$ (33)
where the normalization condition (17) has been used, and the curvature of the chain is constant
$`\kappa (s)_s\theta ={\displaystyle \frac{A}{1w}}.`$ (34)
This case corresponds to a circular chain
$$x=R\mathrm{sin}\frac{s}{R},y=R\mathrm{cos}\frac{s}{R}.$$
The radius $`R`$ of the circle can be obtained by puting Eq. (34) into the boundary condition (20). As a result we have
$`R={\displaystyle \frac{L}{2\pi }}.`$ (35)
The energy of the circular chain is thus
$`_{circ}={\displaystyle \frac{2\pi ^2}{L}}k(1w).`$ (36)
* Polygonally deformed chain.
Let us consider now the case of spatially non-uniform distributed electrons. Inserting Eqs. (29)-(31) into Eq. (23) we get
$`=J\nu {\displaystyle \underset{0}{\overset{L}{}}}\left(_s\phi \right)^2๐s+{\displaystyle \frac{2\pi ^2k}{L}}I.`$ (37)
We restrict our analytical consideration to the case when the charge-curvature coupling is weak and/or the charge density is low: $`w\mathrm{\hspace{0.17em}1}`$. Expanding the functional $`I`$ in terms of the small parameter $`w`$ we obtain from Eq. (37)
$`={\displaystyle \frac{2\pi ^2k}{L(1+w)}}+J\nu {\displaystyle \underset{0}{\overset{L}{}}}\left\{\left(_s\phi \right)^2{\displaystyle \frac{G\nu }{R^2}}\phi ^4w{\displaystyle \frac{G\nu }{R^2}}\phi ^6\right\}๐s`$ (38)
where
$`G={\displaystyle \frac{2\chi ^2}{Jk(1+w)^2}}`$ (39)
is an effective nonlinear parameter. For small $`w`$ one can neglect the last term in Eq. (38) and the Euler-Lagrange equation for the functional (38) then takes the form
$`_s^2\phi +2{\displaystyle \frac{G\nu }{R^2}}\phi ^3\lambda \phi =0.`$ (40)
Straightforward calculations show that Eq. (40) has a solution of the form
$`\phi =R\sqrt{{\displaystyle \frac{\lambda }{(2m)G\nu }}}dn\left(\sqrt{{\displaystyle \frac{\lambda }{(2m)}}}s|m\right)`$ (41)
where $`dn(u|m)`$ is the Jacobi elliptic function with the modulus $`m`$ abr . Inserting Eq. (41) into the boundary condition (28) and the normalization condition (24), we find that the Lagrange multiplier $`\lambda `$ and the modulus $`m`$ are determined by the equations
$`\sqrt{{\displaystyle \frac{\lambda }{(2m)}}}{\displaystyle \frac{L}{n}}=2K(m)`$ (42)
$`G\nu ={\displaystyle \frac{n^2}{\pi ^2}}K(m)E(m)`$ (43)
and the charge distribution along the chain is given by
$`\phi =\sqrt{{\displaystyle \frac{K}{E}}}dn\left({\displaystyle \frac{2nK}{L}}s|m\right)`$ (44)
where $`K(m)`$ and $`E(m)`$ are the complete elliptic integrals of the first kind and the second kind, respectively abr .
The curvature of the chain is given by the equation
$`\kappa (s)_s\theta {\displaystyle \frac{1}{R}}I\left(1+w{\displaystyle \frac{K}{E}}dn^2\left({\displaystyle \frac{2nsK}{L}}|m\right)\right).`$ (45)
Integrating Eq. (45), we get
$`\theta (s)={\displaystyle \frac{2\pi }{L}}I\left(s+w{\displaystyle \frac{L}{2nE}}E\left(\alpha |m\right)\right)`$ (46)
where $`E\left(\alpha |m\right)`$ is the incomplete elliptic integral of the second kind, and $`\alpha =am\left(\frac{2nK}{L}s|m\right)`$ is the amplitude function abr . By using the Fourier expansion for the amplitude function for small $`m`$ we obtain from Eq. (46)
$`\theta (s){\displaystyle \frac{2\pi }{L}}(1+w)Is+{\displaystyle \frac{w}{4n}}Im\mathrm{sin}\left({\displaystyle \frac{2n\pi }{L}}s\right).`$ (47)
Inserting Eq. (47) into the closure condition (21), we find that it is satisfied for $`n2`$. Eqs. (III.1) and (47) describe a polygon: for $`n=2`$ it is an elliptically deformed chain, while for $`n=3`$ it has a triangular shape (see Fig. 1 ). We see from Eqs. (44) and (45) that the polygon structure is a result of the self-consistent interaction between electrons and bending degrees of freedom: extrema of the curvature and of the charge density correlate: in the case of the softening electron-curvature interaction ($`\chi >0`$ maxima of curvature and charge density coincide, while in the case of the hardening interaction the minima of the curvature coincide with the maxima of the charge density. Eq. (43) shows that, for a given value of the nonlinear parameter $`G`$, the $`n`$-gon structure appears when the charge density exceeds the threshold value $`\nu _n`$:
$`\nu >\nu _n{\displaystyle \frac{n^2}{4G}}`$ (48)
.
The energy difference between the $`n`$-gon structure and the circular chain is given by the expression
$`_n_{circ}={\displaystyle \frac{4\pi ^2}{3L}}{\displaystyle \frac{GJ\nu ^2}{E^2}}\left(3E^2(2m)EK(1m)K^2\right).`$ (49)
The normalized energy difference
$`\mathrm{\Delta }_n={\displaystyle \frac{_n_{circ}}{_{circ}}}`$ (50)
for $`n=2,3`$ versus the charge density is shown in Fig. 2. We note that when the charge density is above the critical value the deformed structure with spatially inhomogeneous charge distribution is energetically more favorable than the circular system with a uniformly distributed charge. The state with elliptically deformed chain $`n=2`$ is the energetically most preferable.
Note also that our analytical approach was based on the assumption that $`w\phi ^2\mathrm{\hspace{0.17em}1}`$. Taking into account Eq. (44), this means that it is legitimate to consider not too sharp distributions which correspond to $`wK(m)1`$ or $`m1\mathrm{exp}\{0.72/w\}.`$
## IV Numerical studies
To check our results we have performed also several numerical studies. To this end we carried out the dynamical simulations of the equations
$`\eta {\displaystyle \frac{d}{dt}}\stackrel{}{r}_n={\displaystyle \frac{H}{\stackrel{}{r}_n}},`$ (51)
$`i{\displaystyle \frac{d}{dt}}\psi _n={\displaystyle \frac{H}{\psi _n^{}}}`$ (52)
with the Hamiltonian $`H`$ being defined by Eqs. (2)-(10). Thus the conformational dynamics is considered in an overdamped regime with the friction coefficient $`\eta `$. Then we took as our starting configurations systems involving the electric charge density of (almost) the same magnitude ($`\psi _n`$) at all points (we broke the symmetry by increasing the density at one point of the chain by $`1\%`$). Initially, all the lattice points were placed at symmmetric points on the circle of an appropriate radius (see Fig.3). We performed such simulations for several values of the charge density. Due to the absorption the energy of the system was decreasing during the evolution and the system was evolving towards a minimum. At the same the points of the chain were moving from their initial to their new positions.
We considered both the cases of the hardening and of the softening electron-curvature interaction.
In the case of a hardening electron-curvature interaction ($`\chi <\mathrm{\hspace{0.17em}0}`$) a typical final distribution of the chain points is shown in Fig.4. The charge density and the curvature distributions are in full agreement with the results of our analytical considerations: the curve is more flat where the density of the electrons is maximal. In fact, the ellipse-like shape is rather robust as it arises for a large range of parameters (of the strength of the hardening electron-curvature interaction and of the anharmonicity coefficient $`\kappa _{max}`$).
Complexes with a softening electron-curvature interaction are much more flexible. Their equilibrium shape depends drastically both on the anharmonicity and on the charge density. Figs 5 and 6 demonstrate how drastically the shape of the complex and the charge distribution along the chain can change as a function of the total charge density $`\nu `$: increasing the total charge by $`5\%`$ can lead to the localisation of almost the whole charge of the system at one place.
In our numerical work we also studied the stability of our โfinalโ field configurations - i.e. the configurations which we thought the system was settling at. This we studied by perturbing the system. Such perturbations were introduced in two stages. First we changed the electric charge of the configuration by multiplying all โfinalโ values of the electric charge by a constant factor $`\mu `$; this had the effect of changing the energy of the system. Then we performed the new minimisation and, when the system appeared to have settled at the new โfinalโ configuration, we changed back its $`\psi _n`$ by a new multiplication by $`1/\mu `$. As $`_n|\psi _n|^2`$ is conserved during the evolution, the final system had the same value of it as the original โunperturbedโ fields. The results of the further minimisation were then compared with the original โfinalโ fields.
When we applied this technique to our field configuration shown in Fig. 6 we found that the system was really unchanged by this perturbation; in fact the perturbation led to an overall rotation of the system by one lattice point, but the sequence of values of the fields was essentially the same thus showing the stability of the found minimum.
## V Effects of intercomplex interaction
The aim of this section is to investigate how the interaction between complexes influences the shape. We will consider the system which is described by the Hamiltonian
$`={\displaystyle \underset{j}{}}_j+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}U_{ij},`$ (53)
where $`_j`$ is the energy of the $`j`$-th aggregate which is given by Eq. (23) with $`\phi `$ replaced by $`\phi _j`$ and $`\theta `$ replaced by $`\theta _j`$, and $`U_{ij}`$ is the interaction energy between particles. The latter we will take in the form of the Gay-Berne potential gayberne which is a generalization of the Lennard-Jones 12-6 potential and is widely used to study translational and orientational ordering in systems of aspherical molecules. We consider small deviations from the ringlike structure of aggregates and so we neglect the difference of the well depths for side-to-side and end-to-end configurations. In this case the Gay-Berne potential between two parallel uniaxial molecules is given by
$`U_{ij}U\left(\stackrel{}{r}_{ij}\right)={\displaystyle \frac{U_0}{\sqrt{1\zeta ^2}}}\left[\left({\displaystyle \frac{\sigma _0}{r_{ij}+\sigma \left(\widehat{\stackrel{}{r}}_{ij}\right)+\sigma _0}}\right)^{12}\left({\displaystyle \frac{\sigma _0}{r_{ij}+\sigma \left(\widehat{\stackrel{}{r}}_{ij}\right)+\sigma _0}}\right)^6\right]`$ (54)
where $`\stackrel{}{r}_{ij}=r_{ij}\widehat{\stackrel{}{r}}_{ij}`$ is the interparticle vector,
$`\sigma \left(\widehat{\stackrel{}{r}}_{ij}\right)=\sigma _0\left[1{\displaystyle \frac{2\zeta }{1+\zeta }}\left(\widehat{\stackrel{}{r}}_{ij}\stackrel{}{e}\right)^2\right]^{1/2}`$ (55)
is the anisotropy parameter where $`\stackrel{}{e}`$ is a unit vector specifying the axes of symmetry. The anisotropy coefficient $`\zeta `$ is determined by the lengths of the major and minor axes $`\sigma _{||}`$ and $`\sigma _{}`$
$`\zeta ={\displaystyle \frac{\sigma _{||}^2\sigma _{}^2}{\sigma _{||}^2+\sigma _{}^2}}`$ (56)
and $`\sigma _0`$ gives a charasteristic length scale while $`U_0`$ determines the intensity of the interaction.
The centers of densely packed circular aggregates of the radius $`R`$ create a two-dimensional triangular lattice $`\stackrel{}{r}_j=j_1\stackrel{}{a}_1+j_2\stackrel{}{a}_2,(j(j_1,j_2),j_1,j_2=0,\pm 1,\pm 2,\mathrm{})`$ with the basic vectors $`\stackrel{}{a}_1=\mathrm{}(1,0)`$ and $`\stackrel{}{a}_2=\mathrm{}(1/2,\sqrt{3}/2)`$ where $`\mathrm{}`$ is the lattice constant (see Fig. 7). Being elliptically deformed in such a way that the major axes of all aggregates are parallel to the $`x`$-axis, the centres of densely packed elliptical aggregates create a lattice $`\stackrel{}{r}_j=j_1\stackrel{}{b}_1+j_2\stackrel{}{b}_2`$ with the basic vectors $`\stackrel{}{b}_1=\mathrm{}(1+u,0)`$ and $`\stackrel{}{b}_2=\mathrm{}(1/2(1+u),\sqrt{3}/2(1u))`$ (see Fig. 8), where the parameter $`u`$ is given by
$`u={\displaystyle \frac{\sigma _{||}\sigma _{}}{\sigma _{||}+\sigma _{}}}.`$ (57)
Now we study the ground state of this system by using a trial function approach. We assume that this state is spatially homogeneous and relying on the results of the previous section, we assume that the electron trial function and the trial curvature can be taken in the form
$`\phi _j=\left(\mathrm{cos}\alpha +\sqrt{2}\mathrm{sin}\alpha \mathrm{cos}\left({\displaystyle \frac{2s}{R}}\right)\right),`$ (58)
$`_s\theta _j={\displaystyle \frac{1}{R}}\left(1+\gamma \mathrm{cos}\left({\displaystyle \frac{2s}{R}}\right)\right)`$ (59)
The functions (58) and (59) can be considered as a truncated Fourier expansion of the solutions (41) and (45) in which the coefficients $`\alpha `$ and $`\gamma `$ are variational parameters and $`R`$ is the radius of the cylindrically symmetric aggregate. The function (58) satisfies both the periodicity condition (28) and the number of particles constraint (17). In the limit $`\gamma <1`$ the shape of the curve with the curvature given by Eq. (59) is parametrically determined by the expressions
$`x(s)=R\left((1+{\displaystyle \frac{\gamma }{4}})\mathrm{cos}\left({\displaystyle \frac{s}{R}}\right)+{\displaystyle \frac{\gamma }{12}}\mathrm{cos}\left({\displaystyle \frac{3s}{R}}\right)\right),`$
$`y(s)=R\left((1{\displaystyle \frac{\gamma }{4}})\mathrm{sin}\left({\displaystyle \frac{s}{R}}\right)+{\displaystyle \frac{\gamma }{12}}\mathrm{sin}\left({\displaystyle \frac{3s}{R}}\right)\right).`$ (60)
Thus the lengths of the major and minor axes of the curve (V) are given by
$`\sigma _{||}=R(1+{\displaystyle \frac{\gamma }{3}}),\sigma _{}=R(1{\displaystyle \frac{\gamma }{3}})`$ (61)
and comparing Eqs (57) and (61), we see that $`u=\gamma /3.`$
Inserting Eqs. (58) and (59) into Eqs. (53), (23) and (54) for an energy per aggregate we get
$`{\displaystyle \frac{}{N_a}}=_{tr}+U`$ (62)
where
$`_{tr}={\displaystyle \frac{\pi k}{R}}\left\{1+{\displaystyle \frac{\gamma ^2}{2}}{\displaystyle \frac{1}{8}}\xi \nu \left(8+5\gamma ^2\gamma ^2\mathrm{cos}(2\alpha )+8\sqrt{2}\mathrm{sin}(2\alpha )\right)+{\displaystyle \frac{4J}{k}}\nu \mathrm{sin}^2\alpha \right\}`$ (63)
is the energy of an isolated aggregate, and
$`U={\displaystyle \underset{j=0}{\overset{5}{}}}U\left(\stackrel{}{\mathrm{\Delta }}_j\right)`$ (64)
is the energy due to the interaction between aggregates in the lattice. In Eq. (64) the function $`U\left(\stackrel{}{\mathrm{\Delta }}_j\right)`$ is given by Eqs. (54), (55) with $`\stackrel{}{e}=(1,0)`$ and vectors $`\stackrel{}{\mathrm{\Delta }}_j=\sigma _0((1+\frac{\gamma }{3})\mathrm{cos}\left(\frac{\pi j}{3}\right),(1\frac{\gamma }{3})\mathrm{sin}\left(\frac{\pi j}{3}\right))`$ connect nearest and next-nearest neighbours of the lattice. The interaction energy (64) has a minimum at $`\mathrm{}=2^{1/6}\sigma _0,\gamma =0`$ which corresponds to a system of densely packed circular aggregates. Expanding the function (64) in the vicinity of this point, in powers of the variational parameter $`\gamma `$, we get
$`U={\displaystyle \frac{3}{2}}U_0+cU_0\gamma ^2+\mathrm{}`$ (65)
where the numerical coefficient $`c2.05.`$ According to the variational principle we should satisfy the equations
$`_\alpha _{tr}=0,`$ (66)
$`_\gamma \left(_{tr}+U\right)=0.`$ (67)
From Eq. (66) we get
$`\mathrm{tan}(2\alpha )={\displaystyle \frac{8\sqrt{2}k\xi }{32Jk\xi \gamma ^2}}\gamma .`$ (68)
Inserting Eq. (68) into Eq. (63) and expanding it in terms of $`\gamma `$ we obtain that
$`_{tr}=_{circ}+{\displaystyle \frac{\pi k}{2R}}\left(1{\displaystyle \frac{\nu }{\nu _{cr}}}\right)\gamma ^2+B\gamma ^4+\mathrm{}`$ (69)
where
$`\nu _{cr}={\displaystyle \frac{k}{2\chi }}{\displaystyle \frac{J}{J+\chi }}`$ (70)
is the critical charge density and the notation
$$B=\frac{\pi }{16R}\frac{\chi ^3\nu }{J^3}(2\chi J)$$
is introduced. Thus in the framework of the variational approach the energy of a single complex has a single minimum at $`\gamma =0`$ when $`\nu <\nu _{cr}`$ and in this case the aggregate has a ring-like shape. When $`\nu >\nu _{cr}`$ and $`2\chi >J`$ the energy (69) possesses two equivalent minima with $`\pm \gamma _0(\gamma _00`$) ( see Fig. 9). As it is seen from Eqs. (V) the finite value of $`\gamma _0`$ implies that the aggregate is elliptically deformed either along the $`x`$-axis (when $`\gamma >0`$) or along the $`y`$-axis (when $`\gamma <0`$). Note that in the limit $`\chi J`$, $`\nu _{cr}`$ coincides with $`\nu _2`$ given by Eq. (48). Combining (65) and (69) we see that the interaction between aggregates modifies the condition for appearance of the low-symmetry form. Indeed, even in the case when $`\nu >\nu _{cr}`$ (an isolated aggregate has an ellipse-like shape) in the condensed phase of aggregates for strong enough inter-aggregate interaction we may have an inequality
$`\nu <\nu _{cr}\left(1+{\displaystyle \frac{2cU_0}{\pi k}}R\right)`$ (71)
which means that interacting aggregates are ring-like and create a densely packed crystallic structure with the group symmetry $`D_{6h}`$ (see Fig. 7).When
$`\nu >\nu _{cr}\left(1+{\displaystyle \frac{2cU_0}{\pi k}}R\right)`$ (72)
the ellipse-like shape survives in the condensed phase which has a rectangular unit cell and transforms in accordance with the group symmetry $`D_{2h}`$ ( see Fig.8). There are two equivalent types of arrangement in the low-symmetry phase when the aggregates are elliptically deformed either along the $`x`$-axis (when $`\gamma >0`$) or along the $`y`$-axis (when $`\gamma <0`$). In accordance with this the condensed phase of the aggregates must have a domain structure, i.e. it must consist of various regions in which the direction of long axes are different.
## VI Application to bacteriochlorophyll $`a`$ complexes
The X-ray crystallography shows dermott that bacteriochlorophyll a (Bchl a) molecules in the LH complex are organized in two concentric rings: the B800 and B850 rings. The former consists of nine well-separated Bchl a molecules with an absorption band at 800 nm and the latter consists of 18 Bchl molecules with an absorption band at 860 nm (see Fig. 1 in Ref. kete ). Measurements of the anisotropic properties of the absorption of isolated LH2 complexes from bound to mica surfaces bopp , the fluorescence-excitation spectra of individual LH2 complexes from bacteria Rhodopseudomonas acidophila oijen ; kete ; matsu showed that the complexes are generally not cylindrically symmetric but reveal a deformation of the circular complex into a shape with $`C_2`$ symmetry. Quite recently the shape of the LH2 complex from bacteria Rhodobacter spheroides 2.4.1 in detergent solution has been determined from synchrotron small-angle X-ray scattering data hong . It was shown that, in contrast to the cylindrical crystal structure with a diameter of $`6.8nm`$, the shape of an isolated LH2 complex is an oblate plate with an eccentricity $`ฯต=0.59`$. It was conjectured in Ref. oijen that the extremely dense packing of LH2 in the crystals causes the cylindrical symmetry of the complexes. A variety in shapes and conformations for light-harvesting LH1 complexes and different types of their packing in two-dimensional crystals revealed by atomic force microscopy were reported recently in Ref. bahatyrova .
A qualitative explanation of this phenomena follows from the theory developed in previous sections. It has shown that in the presence of sufficiently strong charge-curvature interaction an isolated complex has an ellipse-like shape while the interaction between complexes in the form of an anisotropic Gay-Berne potential stabilizes the ringlike shapes of the complexes. However if the intensity of the intercomplex interaction is less than some threshold value (see Eqs. (65) and (69)) the non-circular shape of the complex is preserved in the condensed phase.
The model is too simplified for quantitative predictions for light-harvesting complexes but we believe that it contains interesting physics which should be important for further studies of such systems.
## VII Conclusions
In this paper, we have investigated the role of the electron-curvature interaction on the formation of the ground state of closed semiflexible molecular chains. We have found that the coupling between electrons and the bending degrees of freedom of the chain can induce a local softening or hardening of chain bonds, i.e. the effective bending rigidity of the semiflexible chain changes as the density of charge changes along the chain. When the charge density and/or the strength of the electron-curvature coupling exceed a threshold value, the spatially uniform distribution of the charge along the chain and the circular, cylindrically symmetric shape of the chain become unstable. In this case the ground state of the system is characterized by a spatially non-uniform distribution of electrons along the chain and the chain takes on an ellipse-like (or in general polygon-like) form.
## Acknowledgments
Yu. B. G. would like to thank L. Valkunas for attracting his attention to the problem of structural deformations of light-harvesting complexes. Yu.B.G. and W.J.Z. thank L. Brizhik, A. Eremko and B. Piette for their interest and helpful discussions. Yu.B.G. thanks for a Guest Professorship funded by Civilingeniรธr Frederik Christiansens Almennyttige Fond and MIDIT, SNF grant $`\mathrm{\#}`$ 21-02-0500 as well as support from the research center of quantum medicine โVidgukโ. He also thanks Department of Physics, the Technical University of Denmark for hospitality.
## Appendix A
The electron-conformation interaction can be obtained from the Coulomb interaction between an electron and charged groups of the molecule
$`H_{int}={\displaystyle \underset{n,n^{}}{}}V\left(|\stackrel{}{r}_n\stackrel{}{r}_n^{}|\right)|\psi _n|^2`$ (73)
where the matrix element $`V\left(|\stackrel{}{r}_n\stackrel{}{r}_n^{}|\right)`$ describes the Coulomb interaction between an electron which occupies the site $`n`$ in the chain and the charged group at the site $`n^{}`$. From the inextensibility condition (12) we see that the nearest neighbour interaction terms ($`|nn^{}|=1`$) do not contribute to the electron-conformational interaction. Taking into account that
$`\left(\stackrel{}{r}_{n+2}\stackrel{}{r}_n\right)^2=2a^2+2\left(\stackrel{}{r}_{n+1}\stackrel{}{r}_n\right)\left(\stackrel{}{r}_{n+2}\stackrel{}{r}_{n+1}\right)=4a^2a^2\kappa _{n+1}^2`$ (74)
we obtain from Eq. (73) that, in the next-nearest-neighbour approximation, the energy of the electron-conformational interaction takes the form
$`H_{elconfapp}={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\chi \left(\kappa _{n+1}^2+\kappa _{n1}^2\right)|\psi _n|^2`$ (75)
where the notation
$`\chi ={\displaystyle \frac{1}{2}}a{\displaystyle \frac{dV}{dr}}|_{r=2a}`$ (76)
has been used. |
warning/0506/quant-ph0506185.html | ar5iv | text | # Quantum limited velocity readout and quantum feedback cooling of a trapped ion via electromagnetically induced transparency
## I Introduction
Quantum feedback control employs the strategy of acting on a system based on measurement data obtained by continuous observation of the quantum system of interest, thus achieving control of quantum dynamics and preparation of particular quantum states Wiseman and Milburn (1993a); Wiseman (1994); Doherty and Jacobs (1999); Wiseman and Doherty (2005). A prerequisite of developing quantum feedback control is the realization of quantum limited measurements. Quantum optical systems and, more recently, mesoscopic systems have taken a leading role in achieving these requirements towards demonstration of feedback control in the laboratory Fischer et al. (2002); Steck et al. (2004); Mancini et al. (2000); Dunningham et al. (1997); Smith et al. (2002); Geremia et al. (2004); Mancini et al. (1998); Cohadon et al. (1999); Hopkins et al. (2003).
Motivated by the remarkable experimental progress with trapped ions, we will develop in the present paper a theory of quantum feedback cooling of a single ion, based on a continuous readout of the atomic velocity via dispersive interactions with laser light. The idea is to devise an (in principle) quantum limited measurement of the velocity by employing the strong detuning (and thus velocity) dependence of the index of refraction of a driven $`\mathrm{\Lambda }`$-system near the atomic dark state. These dark states are coherent superpositions of two atomic ground states which do not couple to the excited atomic state, which leads to strong suppression of dissipative light scattering. This is the same feature which underlies recent studies of electromagnetically induced transparency, slow light in atomic gases and quantum memory of light in atomic ensembles. The present setup of dispersive readout of the atomic velocity complements and is in contrast to ongoing experiments of quantum feedback cooling of a single two-level ion in front of a mirror Bushev (2004); Bushev et al. (2005), where the position of the ion is continuously monitored by emission of light into the mirror mode, as analyzed theoretically in our recent publication Steixner et al. (2005).
Our discussion of quantum feedback cooling of a single trapped ion in a strongly driven atomic $`\mathrm{\Lambda }`$-system builds on, and connects various well-developed topics in atomic physics and quantum optics, as well as continuous measurement and quantum feedback theory. Thus we find it worthwhile to present both a brief review of the background material and physical key ideas underlying the present work, as well as the main results of the paper in Section II. In Section III we give the technical details of our model and derive the equations for the measured signal and the conditioned evolution of the atomic motion for a weakly excited atom. A Wiseman-Milburn-type Wiseman and Milburn (1993a) master equation for feedback cooling and the resulting temperatures will be discussed in Section IV. Finally, in Section V we make a connection to EIT-laser cooling and describe combination of feedback and laser cooling. The details of the adiabatic elimination of the internal atomic states are given in the Appendixes A and B.
## II Overview and Summary
In this section we present an overview of the concepts and the main results of this paper. Our emphasis will be on explaining the basic physics and strategy behind our quantum feedback scheme, and providing references to later sections where the mathematically inclined reader can find the details of the derivations. In Subsection II.1 we will briefly review EIT and discuss the dependence of the index of refraction on the ion momentum near the dark state resonance. Continuous read out of the ion momentum using homodyne detection will be formulated in Subsection II.2. The main results of the present paper are the equations for feedback cooling in Subsection II.3, and the predictions for the final temperatures and cooling rates in II.4, in particular also in connection to EIT laser cooling of ions Morigi et al. (2000); Morigi (2003); Roos et al. (2000).
### II.1 Electromagnetically Induced Transparency
The phenomenon of EIT is related to a quantum interference effect which in its simplest form can be observed in a three level atom (for a review see: Scully and Zubairy (1997); Lukin (2003) and references therein). To discuss this effect we consider an atom with the internal states $`|g`$, $`|e`$ and $`|r`$ in a $`\mathrm{\Lambda }`$-configuration as shown in Fig. 1. The transition between the states $`|r`$ and $`|e`$ is driven by a strong, resonant laser, while a second light field, the probe field, couples the ground state to the excited state. We denote by $`\mathrm{\Delta }_p=\omega _p\omega _{eg}`$ the detuning of the probe field from the atomic resonance. The atomic Hamiltonian is then given by
$$\begin{array}{c}\hfill H_\mathrm{\Lambda }=\mathrm{}\mathrm{\Delta }_p|gg|+\frac{\mathrm{}}{2}(\mathrm{\Omega }_L|er|+g|eg|+h.c.),\end{array}$$
(1)
where $`\mathrm{\Omega }_L`$ and $`g`$ are the Rabi frequencies of the laser and the probe field respectively. At the two photon resonance, $`\mathrm{\Delta }_p=0`$, the Hamiltonian, $`H_\mathrm{\Lambda }`$ has an adiabatic eigenstate with zero energy, a so-called โdark stateโ,
$$|D\mathrm{\Omega }_L|gg|r.$$
(2)
For an atom in the state $`|D`$, the excitation from the states $`|g`$ and $`|r`$ destructively interfere and the atom decouples from the light. We note that an atom in a dark state involves no excited state population, and is thus immune to decay from the excited state.
The existence of such a dark state leads to remarkable properties of the index of refraction. For weak probe fields, the propagation can be discussed in terms of the linear susceptibility, $`\chi (\omega _p)`$. The real part of $`\chi (\omega _p)`$ is related to the refractive index by $`n=1+\mathrm{Re}(\chi (\omega _p))/2`$, while the imaginary part, $`\mathrm{Im}(\chi (\omega _p))`$, is proportional to the absorption coefficient. For an ensemble of three level atoms where both, $`|g`$ and $`|r`$ are long-lived states, the susceptibility has the characteristic behavior Lukin (2003),
$$\chi (\mathrm{\Delta }_p)\frac{i\mathrm{\Delta }_p}{|\mathrm{\Omega }_L|^2/4+i\mathrm{\Delta }_p(\mathrm{\Gamma }+i\mathrm{\Delta }_p)},$$
(3)
where $`\mathrm{\Gamma }`$ is the decay rate of the excited state. Fig. 1b shows the dependence of $`\chi `$ on the detuning, $`\mathrm{\Delta }_p`$. Around the dark resonance $`\mathrm{\Delta }_p=0`$ we have a steep slope of the refractive index (solid line) which leads to a slow group velocity of the probe field (โslow lightโ) while absorption is strongly suppressed (dotted line), giving rise to โelectromagnetically induced transparencyโ. The width of the transparency window as well as the variation of the refractive index depend on $`\mathrm{\Omega }_L`$ and can be controlled by the laser field. We note the suppression of absorption at the dark state resonance (dotted line in Fig. 1b).
Consider now a single ion in a $`\mathrm{\Lambda }`$-configuration (Fig. 2a) moving in a trapping potential. We will only consider the 1D motion along the propagation direction of both the probe and dressing laser beams. If we adopt for the moment a classical description of the ion motion with $`z(t)`$ the ion trajectory, then the internal dynamics of the ion can again be described by the Hamiltonian (1) with a Doppler-shifted probe detuning
$$\mathrm{\Delta }_p(t)=\mathrm{\Delta }_p+(k_pk_L)v(t),$$
where $`k_p`$ and $`k_L`$ the wave vectors of the running probe field and the dressing laser field, respectively, and $`v(t)\dot{z}(t)`$ is the ion velocity. Thus for a resonant probe field, $`\mathrm{\Delta }_p=0`$, the change in the index of refraction is a linear function of the atomic velocity or momentum with the steep slope given in Eq. (3).
This suggests the strategy to measure the momentum of the ion continuously by monitoring the phase shift due to the varying index of refraction. We note that - while the index of refraction of a single particle is small - EIT will strongly amplify the sensitivity to the velocity. At the same time dissipation due to light scattering is strongly suppressed within the transparency window. These arguments can be easily adapted to a situation where the atomic motion is quantized.
### II.2 Continuous observation of the momentum of the trapped ion
We turn now to a formulation of the continuous read out of the ion momentum as outlined in Fig. 2. The idea is to measure the momentum of the ion via the phase changes of the probe beam as described in the previous subsection. The phase of the probe beam can be determined by homodyning, i.e. by mixing the probe beam with a local oscillator and measuring the homodyne current, $`I_c(t)`$.
The state of the observed system (the moving ion) is described by a conditional density operator $`\mu _c(t)`$, which represents the observerโs knowledge of the current state of the system for a given record of the measured signal, $`I_c(t)`$. We will show in Sec. III that after adiabatic elimination of the excited state of the (weakly driven) ion the measured homodyne current has the form,
$$I_c(t)=2ฯต\mathrm{\Gamma }_0\widehat{p}_c+\sqrt{ฯต\mathrm{\Gamma }_0}\xi (t).$$
(4)
The current is the sum of two terms. The first contribution shows a linear dependence on the conditional expectation value of the momentum operator $`\widehat{p}_c\mathrm{Tr}\{\mu _c\widehat{p}\}`$. Thus by measuring $`I_c(t)`$ we learn the momentum of the moving ion. The second contribution describes a shot noise term with $`\xi (t)`$ a white noise Gaussian process. The signal strength is determined by the rate $`\mathrm{\Gamma }_0`$. It is related to the slope of the refractive index $`\chi `$ at $`\mathrm{\Delta }_p=0`$: an explicit expression is given in Sec. III Eq. (39) below. The parameter $`0<ฯต1`$ takes into account the collection efficiency of the scattered photons (where in the ideal case $`ฯต=1`$). Eq. (4) is derived under the assumption that $`\mathrm{\Omega }_L`$ is large compared to typical Doppler detuning, $`\mathrm{\Delta }_D`$, and is valid on time scale which is slow compared to $`\mathrm{\Omega }_L^1`$. The signal is maximized for the local oscillator phase, $`\varphi =0`$.
According to continuous measurement theory applied to homodyne detection, the conditional density operator $`\mu _c(t)`$ is updated upon observation of the current $`I_c(t)`$ following the Ito equation,
$`d\mu _c(t)`$ $`=`$ $`i\nu [\widehat{a}^{}\widehat{a},\mu _c(t)]dt`$
$`+_M\mu _c(t)dt+\sqrt{ฯต\mathrm{\Gamma }_0}[\widehat{p}]\mu _c(t)dW(t).`$
This equation will be derived in Sec. III, Eq. (38).
The first term in this equation describes the free evolution in the 1D harmonic trap where $`\nu `$ denotes the trap frequency, and $`\widehat{a}`$ ($`\widehat{a}^{}`$) the destruction (creation) operators, respectively. The effects of the continuous observation appear in the second and third term of Eq. (II.2).
The superoperator $`_M`$ determines the back action of the measurement setup on the atomic motion. In the Lamb-Dicke limit $`\eta =2\pi a_0/\lambda _p1`$, where the extension of the atomic wavepacket (size of the harmonic oscillator ground state $`a_0`$) is much smaller than the wavelength of the light, $`\lambda _p`$, it has the form,
$$_M\mu =\frac{\mathrm{\Gamma }_0}{2}[\widehat{p},[\widehat{p},\mu ]].$$
(6)
The action of $`_M`$ tends to diagonalize the density operator in the eigenbasis of the (measured) operator $`\widehat{p}`$. By comparing the decoherence rate, $`\mathrm{\Gamma }_0`$, with the signal strength, $`ฯต\mathrm{\Gamma }_0`$, we see that for $`ฯต<1`$ the measurement is not quantum limited, i.e., more noise is added than required by quantum mechanics Braginsky and Khalili (1992). Although the measurement does not reach the quantum limit, the back action is still minimal for a given collection efficiency.
In the third term of Eq. (II.2) we introduced the notation,
$$[\widehat{c}]\mu =\widehat{c}\mu +\mu \widehat{c}c+\widehat{c}\mu .$$
This term describes the observerโs knowledge of the current state of the system and therefore depends on the measured signal. The Wiener increment $`dW(t)`$ is formally related to the signal noise by $`dW(t)\xi (t)dt`$.
In summary, Eq. (4) for the homodyne current $`I_c(t)`$ and the evolution equation (II.2) for the conditional density matrix constitute the basic equations of continuous observation of the momentum of the ion via homodyne detection.
### II.3 Quantum Feedback Cooling
The information on the atomic momentum contained in the signal $`I_c(t)`$ (4) can be used to act back on the system. Here we are interested in cooling the atomic motion by using the feedback strategy known as โcold dampingโ Vitali et al. (2002). The idea is to apply a force on the atom which is proportional but opposite to its momentum. This force creates an effective friction for the atomic motion and, therefore, leads to a dissipation of kinetic energy.
In our setup the measured signal is already proportional to the average momentum, $`\widehat{p}_c`$, and can be amplified and fed back directly. Thus we consider a feedback Hamiltonian of the form,
$$H_{fb}(t)=\frac{G}{2ฯต}I_c(t\tau )\widehat{z},$$
(7)
where $`G`$ is the dimensionless gain factor, and $`\tau `$ is the finite delay in the feedback loop. Note that $`\tau >0`$, so that $`H_{fb}`$ acts after the measurement. Eq. (II.2) with the feedback Hamiltonian added provides us with a feedback equation describing the time evolution of the system. The goal is now to average this equation over the Gaussian white noise $`\xi (t)`$.
A general theory for direct quantum feedback has been first discussed in a seminal paper by Wiseman and Milburn Wiseman and Milburn (1993a). In particular, they have shown how to average the quantum feedback equation in the limit $`\tau 0^+`$. In our case this assumption implies that the time delay of the feedback is small on the scale of the (adiabatically eliminated) system evolution, a condition which is realistic in the present context. Adopting this formalism we will derive in Sec. IV a master equation for the unconditioned density operator, $`\mu (t)=E[\mu _c(t)]`$. We obtain (for $`\mathrm{\Delta }_p=\mathrm{\Delta }_L=0`$)
$$\begin{array}{cc}\hfill \dot{\mu }=& i\nu [\widehat{a}^{}\widehat{a},\mu ]+_M\mu \hfill \\ & i\mathrm{\Gamma }_0\frac{G}{2}[\widehat{z},\widehat{p}\mu +\mu \widehat{p}]\mathrm{\Gamma }_0\frac{G^2}{8ฯต}[\widehat{z},[\widehat{z},\mu ]].\hfill \end{array}$$
(8)
The first line of this equation describes the free evolution and the measurement back action, $`_M`$ Eq. (6). The terms in the second line of Eq. (8) include the effects of the feedback loop. While the term proportional to $`G`$ causes the expected damping of the motion, the second term leads to a diffusion of the atomic momentum. This diffusion originates from the noise in the measured current which is also amplified and fed back to the system.
The generic dependence of the steady state energy on the feed back gain calculated from Eq. (8) is plotted in Fig. 3. The curves show the expected minimum as a function of $`G`$, at the point where the noise of the feedback loop starts to dominate over the damping force. More detailed results will be presented in the following subsection.
### II.4 Results: Quantum Feedback vs. EIT Laser Cooling
In general, the interaction of atoms with light always leads to some form of laser cooling or heating. In the resonant case $`\mathrm{\Delta }_p=\mathrm{\Delta }_L=0`$ (discussed above), which is required to measure the atomic momentum, heating and cooling rates are equal and cause the diffusion described by $`_M`$.
By detuning the lasers away from the resonance, $`\mathrm{\Delta }_L0`$, the atomic susceptibility, $`\chi (\mathrm{\Delta }_p)`$, and therefore the absorption properties become quite asymmetric (see Fig. 4). This asymmetry is exploited in EIT laser cooling (ELC), where the absorption on the red sideband (a phonon is removed from the motion) is much more likely than on the blue sideband (a phonon is added to the motion). A detailed discussion of ELC can be found in Ref. Morigi (2003).
In Sec. V we show that in the Lamb-Dicke limit we can derive the extension of Eq. (8) which includes quantum feedback cooling as well as ELC. In rotating frame with respect to the trap frequency, $`\nu `$, it can be written in the form,
$$\begin{array}{c}\hfill \dot{\mu }=(A_{}+A_{}^{fb})๐[\widehat{a}]\mu +(A_++A_+^{fb})๐[\widehat{a}^{}]\mu ,\end{array}$$
(9)
with
$$๐[\widehat{a}]\mu =\widehat{a}\mu \widehat{a}^{}\frac{1}{2}\widehat{a}^{}\widehat{a}\mu \frac{1}{2}\mu \widehat{a}^{}\widehat{a}.$$
The total cooling and the total heating rate are divided into a contribution from the atom-laser interaction, $`A_\pm `$, and a contribution from the feedback force, $`A_\pm ^{fb}`$. They are given by,
$`A_\pm `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{2}}\mathrm{Re}[I(\pm \nu )],`$
$`A_\pm ^{fb}`$ $`=`$ $`\mathrm{\Gamma }_0\left(G{\displaystyle \frac{\nu \mathrm{\Gamma }}{\mathrm{\Omega }^2}}\mathrm{Im}[I^{}(\pm \nu )e^{i\varphi }]+{\displaystyle \frac{G^2}{8ฯต}}\right).`$
The relation between the four rates is determined by the function $`I(\nu )`$ which is defined as
$$I(\nu )=\frac{\mathrm{\Omega }^4}{2\mathrm{\Gamma }\nu ^2}\frac{i\nu }{(\mathrm{\Omega }^24\nu (\nu \mathrm{\Delta }_L))+i2\mathrm{\Gamma }\nu }.$$
(10)
Note that the rates $`A_\pm `$ are proportional to the imaginary part of the atomic susceptibility at the sideband frequencies, $`\chi (\omega _p\nu )`$, (see Fig. 4). For a general detuning $`\mathrm{\Delta }_L`$, the measured signal is no longer proportional to $`\widehat{p}`$, and therefore the phase of the local oscillator, $`\varphi `$, appears in the feedback rates.
In the basis of harmonic oscillator states the master equation (9) has the form of a standard rate equation for the trap occupations ($`p_n=n|\rho |n)`$ familiar from laser cooling,
$$\begin{array}{cc}\hfill \dot{p}_n=& (A_{}+A_{}^{fb})\left[(n+1)p_{n+1}np_n\right]\hfill \\ & +(A_++A_+^{fb})\left[(n1)p_{n1}(n+1)p_n\right],\hfill \end{array}$$
(11)
which predicts in steady state a Bose-Einstein distribution with a mean occupation number
$$\overline{n}=\frac{A_++A_+^{fb}}{(A_{}A_+)+(A_{}^{fb}A_+^{fb})}.$$
(12)
We now turn to the discussion of results for the case of pure feedback cooling, and combined feedback and laser cooling:
*Pure feedback cooling*. We first reproduce the results of Eq. (8) by setting $`\mathrm{\Delta }_p=\mathrm{\Delta }_L=0`$. Then $`A_\pm =\mathrm{\Gamma }_0/2`$ and the cooling of the atom is attributed to the feedback mechanism. The minimal energy is reached for $`G=\sqrt{4ฯต}`$ and $`\varphi =0`$. It is given by
$$E_{\mathrm{min}}=\frac{\mathrm{}\nu }{2}\sqrt{\frac{1}{ฯต}}.$$
(13)
This expression shows that the final temperature is only limited by the collection efficiency $`ฯต`$. In the theoretical limit, $`ฯต1`$, it approaches the ground state energy, $`\mathrm{}\nu /2`$.
*Feedback cooling and ELC*. When we tune away from the resonance, $`\mathrm{\Delta }_p=\mathrm{\Delta }_L0`$, we obtain a difference in the laser cooling rates, $`A_{}A_+`$. In Fig. 5 we compare the final temperatures for the optimal feedback gain with the case of pure ELC. For blue detuning, $`\mathrm{\Delta }_L>0`$ the mechanism of ELC sets in and cools the atom close to the ground state. Although the addition of the feedback loop always leads to even lower temperatures, its effect can be neglected because for the present parameters ELC already provides efficient ground state cooling. For red detuning $`\mathrm{\Delta }_L0`$ the absorption spectrum (see Fig. 4) is reversed and the atom is actively heated by the light absorption. In this case a steady state is only reached when the feedback cooling dominates over the laser induced heating.
## III The model
In this section we present a detailed description of our model for the measurement setup which is shown in Fig. 2. The system of interest is the three level atom which is confined by an external trapping potential. This atom is illuminated by a strong laser field to create the transparency effect for the probe field. The outgoing probe light is mixed with a strong local oscillator to perform a homodyne measurement to detect linear shifts of the field.
To describe the dynamics of the atom as well as the detection of scattered field, we start with the total Hamiltonian,
$$H=H_A+H_{AEM}+H_{EM}.$$
(14)
It is the sum of the Hamiltonian for the external and internal states of the atom, $`H_A`$, the free Hamiltonian for the electromagnetic environment, $`H_{EM}`$, and the coupling between the atom and the electromagnetic field, $`H_{AE}`$.
For the internal level structure we consider a $`\mathrm{\Lambda }`$-configuration as shown in Fig. 1a. A classical laser field with frequency $`\omega _L`$ drives the transition between the excited state, $`|e`$ and the second ground or metastable state, $`|r`$. The Rabi frequency for this coupling is denoted by $`\mathrm{\Omega }_L`$. For the external dynamics of the atom we restrict ourselves to a one dimensional model, i.e., we assume that the atom is strongly confined in the x and y directions. Along the z-axis, which coincides with the propagation direction of the probe beam, the atom is trapped by the external potential, $`V(z)`$. Although it is not essential for the following discussion, we further assume that $`V(z)`$ is harmonic, with a trap frequency, $`\nu `$. This assumption allows us to introduce the dimensionless position and momentum operators, $`\widehat{z}:=(\widehat{a}+\widehat{a}^{})/\sqrt{2}`$ and $`\widehat{p}:=i(\widehat{a}^{}\widehat{a})/\sqrt{2}`$, where $`\widehat{a}`$ and $`\widehat{a}^{}`$ denote the usual annihilation and creation operators. With the definition of the external Hamiltonian, $`H_E=\mathrm{}\nu \widehat{a}^{}\widehat{a}`$, and the notation $`\sigma _{ij}=|ij|`$ the atom Hamiltonian is then given by
$$\begin{array}{cc}\hfill H_A=H_E& +\mathrm{}\omega _e|ee|+\mathrm{}\omega _r|rr|\hfill \\ & +\frac{\mathrm{}\mathrm{\Omega }_L}{2}(e^{i\omega _rt}e^{i\eta _L\widehat{z}}\sigma _{er}+h.c).\hfill \end{array}$$
(15)
Here $`\omega _{e,r}`$ denote the eigenfrequencies of the corresponding states $`|e`$ and $`|r`$. The Lamb-Dicke parameter is defined as, $`\eta _r=k_L\sqrt{\mathrm{}/m\nu }`$, where $`m`$ is the mass of the atom and $`k_L`$ is the wave vector (projected on the z-axis) of the laser field.
The electromagnetic environment consist of a three dimensional set of plane wave modes, labelled by their wave vector, $`\stackrel{}{k}`$ and their polarization, $`\lambda `$. In terms of the corresponding annihilation and creation operators, $`\widehat{b}_\lambda (\stackrel{}{k})`$ and $`\widehat{b}_\lambda ^{}(\stackrel{}{k})`$, the free evolution is determined by the Hamiltonian,
$$H_{EM}=\underset{\lambda =1,2}{}d^3k\mathrm{}\omega _\stackrel{}{k}\widehat{b}_\lambda ^{}(\stackrel{}{k})\widehat{b}_\lambda (\stackrel{}{k}).$$
(16)
The electric field of the environment interacts with the internal states of the atom via a dipole coupling. Under the rotating wave approximation the interaction Hamiltonian is
$$\begin{array}{c}\hfill H_{\mathrm{A}\mathrm{E}}=\stackrel{}{\mu }_{eg}\stackrel{}{E}^+(\widehat{z})\sigma _{eg}\stackrel{}{\mu }_{er}\stackrel{}{E}^+(\widehat{z})\sigma _{er}+h.c.,\end{array}$$
(17)
with the standard expression for the electric field,
$$\stackrel{}{E}^+(\stackrel{}{x})=i\underset{\lambda =1,2}{}d^3k\sqrt{\frac{\mathrm{}\omega _k}{2ฯต_0(2\pi )^3}}\stackrel{}{\epsilon }_\lambda (\stackrel{}{k})\widehat{b}_\lambda (\stackrel{}{k})e^{i\stackrel{}{k}\stackrel{}{x}}.$$
(18)
In an experiment the lens system defines a certain spatial mode function for the probe beam. To describe the homodyne detection of the probe field, we divide the total electric field into the field of this particular mode, $`\stackrel{}{E}_p`$, and a remaining set of modes, orthogonal to $`\stackrel{}{E}_p`$,
$$\stackrel{}{E}^+(\stackrel{}{x})=\stackrel{}{E}_p^+(\stackrel{}{x})+\stackrel{}{E}_{}^+(\stackrel{}{x}).$$
(19)
In our model we approximate $`\stackrel{}{E}_p`$ by the one dimensional field,
$$\stackrel{}{E}_p^+(z)=i\stackrel{}{_p}e^{i(k_pz\omega _pt)}+i\stackrel{}{\epsilon }_p_0^{\mathrm{}}๐ka(k)e^{ikz}\widehat{b}_p(k).$$
(20)
The first part in this expression describes the coherent field of the incoming probe beam. Note that with this definition of the operator, $`\stackrel{}{E}_p`$, the initial state of the electromagnetic environment is the vacuum state. The function, $`a(k)`$, determines the coupling of the atom to the dense set of modes, $`\widehat{b}_p(k)`$. It must be adjusted to reproduce the correct results of the real mode function.
The interaction of the atom with the coherent part of the probe field leads to transitions between $`|g`$ and $`|e`$, characterized by the Rabi frequency, $`g=2|\stackrel{}{\mu }_{eg}\stackrel{}{_p}|/\mathrm{}`$. In the following include this term in the atomic evolution and define a new system Hamiltonian, $`H_S=H_A+H_{AEM}|_{\mathrm{coh}}`$. In a frame rotating with the laser frequencies, $`\omega _p`$ and $`\omega _L`$, this Hamiltonian is then given by
$$\begin{array}{cc}\hfill H_S=H_E& \mathrm{}\mathrm{\Delta }_p|ee|\mathrm{}(\mathrm{\Delta }_p\mathrm{\Delta }_L)|rr|\hfill \\ & +\frac{\mathrm{}\mathrm{\Omega }_L}{2}\left(e^{i\eta _r\widehat{z}}\sigma _{er}+e^{i\eta _r\widehat{z}}\sigma _{re}\right)\hfill \\ & +\frac{\mathrm{}g}{2}\left(e^{i\eta _g\widehat{z}}\sigma _{eg}+e^{i\eta _g\widehat{z}}\sigma _{ge}\right).\hfill \end{array}$$
(21)
Here we introduced the detunings $`\mathrm{\Delta }_p=\omega _p\omega _e`$ and $`\mathrm{\Delta }_L=\omega _L(\omega _e\omega _r)`$, and the Lamb-Dicke parameter of the probe field, $`\eta _g=k_p\sqrt{\mathrm{}/m\nu }`$.
The coherent evolution of the driven atom is now determined by the system Hamiltonian $`H_S`$. The coupling between the atom and the dense set of modes of $`\stackrel{}{E}_{}`$ and the non-classical part of $`\stackrel{}{E}_p`$ has two effects. First, it leads to an incoherent dynamic of the atomic state. This includes the decay of the excited state population as well as a diffusion of the atomic momentum due to the random recoil kicks of the emitted photons. Second, the coupling changes the state of the electromagnetic environment which is (partially) detectable by the observer.
### III.1 Master equation
We first look at the incoherent dynamics of the atom and ignore the evolution of the electromagnetic bath. By applying the standard Born-Markov approximation and tracing over the bath degrees of freedom we obtain a master equation for the system density matrix, $`\rho `$. It can be written in the standard form, $`\dot{\rho }=\rho `$, with a Liouville operator Stenholm (1986); Cirac et al. (1992); Morigi (2003),
$$\rho =\frac{i}{\mathrm{}}[H_{\mathrm{eff}}\rho \rho H_{\mathrm{eff}}^{}]+๐ฅ_g(\sigma _{ge}\rho \sigma _{eg})+๐ฅ_r(\sigma _{re}\rho \sigma _{er}).$$
(22)
The effective Hamiltonian, $`H_{\mathrm{eff}}=H_Si\mathrm{}\mathrm{\Gamma }/2|ee|`$, includes the unitary evolution of the atom as well as the decay of the excited state population with a total rate, $`\mathrm{\Gamma }`$. The two โrecyclingโ terms are defined by
$$๐ฅ_{j=r,g}(\rho )=\mathrm{\Gamma }_j_1^1๐uN(u)e^{i\eta _j\widehat{z}u}\rho e^{i\eta _j\widehat{z}u}.$$
(23)
The rates $`\mathrm{\Gamma }_g`$ and $`\mathrm{\Gamma }_r`$ denote the decay rates into the corresponding states, $`|g`$ and $`|r`$. The dipole distribution, $`N(u)=\frac{3}{8}(1+u^2)`$, with $`u=\mathrm{cos}(\phi )`$, determines the probability for emitting a photon under a certain angle, $`\phi `$, with respect to the z-axis.
Master equation (22) describes the full dynamics of a driven atom in a $`\mathrm{\Lambda }`$-configuration. The external and internal degrees of freedom are coupled via the position dependent interaction with the electromagnetic field. In general, the recoil kicks of the emitted photons described by $`๐ฅ_{g,r}`$ lead to momentum diffusion and a heating of the atomic motion. For an appropriate choice of laser detunings the heating can be compensated by photon absorptions (laser cooling). In this paper, we follow a different approach where cooling is provided by an external feedback force.
### III.2 Continuous homodyne detection
In a next step we describe the homodyne detection of the probe field, $`\stackrel{}{E}_p`$. Here we follow the standard theory of homodyne detection (see, e.g. Ref. Gardiner and Zoller (2000); Carmicheal (1993); Wiseman and Milburn (1993b)) to derive the relevant equations for the model specified above.
After the interaction with the atom, the outgoing probe field is desribed by the Heisenberg operator $`\widehat{b}_{p,\mathrm{out}}(t)`$. Using the input-output formalism Gardiner and Zoller (2000) it is related to the incoming field, $`\widehat{b}_{p,in}(t)`$, by
$$\widehat{b}_{p,\mathrm{out}}(t)=\widehat{b}_{p,in}(t)+\sqrt{\gamma }\widehat{c}_p(t),$$
(24)
where $`\widehat{c}_p=e^{i\eta _p\widehat{z}}\sigma _{ge}`$ denotes the atomic โjump operatorโ which couples to the probe field. To relate our model to the real experimental setup, we set $`\gamma =ฯต\mathrm{\Gamma }_g`$, where the collection efficiency $`ฯต`$ determines the fraction of photons which are scattered into the mode of the probe field.
In homodyne detection the outgoing field (24) is mixed with a strong coherent field of the local oscillator. When the transmittance of the beam splitter is close to one the field operator at the position of the detector is
$$\widehat{b}_d(t)=\sqrt{\gamma }\beta e^{i\varphi }e^{i\omega _pt}+\widehat{b}_{p,in}(t)+\sqrt{\gamma }\widehat{c}_p(t).$$
(25)
Here, $`\beta `$ and $`\varphi `$ denote the real amplitude and the phase of the reflected part of the local oscillator. Note that the total probe field defined in Eq. (20) is the sum of a classical and a quantized contribution. The classical part of $`\widehat{b}_{p,in}(t)`$ can simply be absorbed in a redefinition of $`\beta e^{i\varphi }`$. The operator for the homodyne current is,
$$\widehat{I}_h(t)=\underset{\beta \mathrm{}}{lim}\left(\widehat{b}_d^{}(t)\widehat{b}_d(t)\gamma \beta ^2\right)/\beta .$$
(26)
The measured signal $`I_c(t)`$ is then defined as the outcome of the continuous measurement of the current operator, $`\widehat{I}_h(t)`$. Using the results form the theory of homodyne detection Gardiner and Zoller (2000), we obtain
$$I_c(t)=ฯต\mathrm{\Gamma }_g\widehat{c}_pe^{i\varphi }+\widehat{c}_p^{}e^{i\varphi }_c(t)+\sqrt{ฯต\mathrm{\Gamma }_g}\xi (t).$$
(27)
The unconditioned evolution given by master equation (22) and the measurement record, $`I_c(t)`$, determine the evolution of the conditioned density operator, $`\rho _c(t)`$. Following Ref. Wiseman and Milburn (1993b) we obtain,
$$d\rho _c(t)=\rho _c(t)dt+\sqrt{ฯต\mathrm{\Gamma }_g}[\widehat{c}_pe^{i\varphi }]\rho _c(t)dW(t).$$
(28)
Eqs. (27) and (28) represent a full description of the conditioned dynamics of the atom under continuous observation and serve as the starting point for the following discussion.
### III.3 Adiabatic Elimination
The current $`I_c(t)`$ as given in Eq. (27) is still a function of the coupled external and internal states of the atom. In the following we show that for a weakly excited atom we can eliminate the dynamics of the internal states, and the measured signal becomes a linear function of the atomic momentum as given in Eq. (4). In addition we derive the resulting back action of the measurement on the motional state of the atom.
As already noted in Section II, the phase shift of the probe light is a linear function of the atomic momentum as long as the typical Doppler detuning, $`\mathrm{\Delta }_D=\nu \eta _p\widehat{p}`$, is small compared to the width of the transparency window, $`\mathrm{\Omega }_L`$. Therefore, we can apply perturbation theory in the parameter $`\mathrm{\Delta }_D/\mathrm{\Omega }_L`$ to derive an effective equation for the external state. As a first step in our calculation we perform a unitary transformation,
$$U=e^{i\eta _p\widehat{z}|ee|}e^{i\overline{\eta }\widehat{z}|rr|},$$
(29)
where we set $`\overline{\eta }:=\eta _p\eta _L`$. In the new basis the system Hamiltonian $`\stackrel{~}{H}_S=U^{}H_SU`$ is given by
$$\begin{array}{cc}\hfill \stackrel{~}{H}_S=H_E& \mathrm{}\left(\mathrm{\Delta }_p\nu \eta _p\widehat{p}\nu \eta _p^2/2\right)|ee|\hfill \\ \hfill & \mathrm{}\left(\mathrm{\Delta }_p\mathrm{\Delta }_L+\nu \overline{\eta }\widehat{p}\nu \overline{\eta }^2/2\right)|rr|\hfill \\ \hfill +& \frac{\mathrm{}\mathrm{\Omega }_L}{2}\left(\sigma _{er}+\sigma _{re}\right)+\frac{\mathrm{}g}{2}\left(\sigma _{eg}+\sigma _{ge}\right).\hfill \end{array}$$
(30)
As in the classical case (see Section II.1) the position dependence of the atom-laser coupling is transformed into a frequency shift for the states $`|e`$ and $`|r`$. In addition to the Doppler detunings, $`\nu \eta _p\widehat{p}`$ and $`\nu \overline{\eta }\widehat{p}`$, the internal states are also shifted by the appropriate recoil frequencies. They account for the fact, that each absorption or emission of a photon also transfers kinetic energy to the atom.
For the particular choice of laser detunings, $`\mathrm{\Delta }_p=\mathrm{\Delta }_L+\nu \overline{\eta }^2/2`$, and in the absence of a trapping potential the Hamiltonian $`\stackrel{~}{H}_S`$ has a dark eigenstate, $`|\psi _D=|p=0,D`$. Here $`|p=0`$ is the zero momentum eigenstate and $`|D`$ denotes the internal dark state,
$$|D=\left(g|r\mathrm{\Omega }_L|g\right)/\mathrm{\Omega },$$
(31)
with $`\mathrm{\Omega }=\sqrt{\mathrm{\Omega }_L^2+g^2}`$. In the following we assume that this relation between the detunings is fulfilled. The system Hamiltonian can then be written as
$$\stackrel{~}{H}_S=H_E+H_I+H_{\mathrm{int}},$$
(32)
such that $`H_E`$ and $`H_I`$ act on the external or the internal states only, while $`H_{\mathrm{int}}`$ describes the coupling between them,
$$H_{\mathrm{int}}=\mathrm{}\nu \eta _p\widehat{p}|ee|\mathrm{}\nu \overline{\eta }\widehat{p}|rr|.$$
(33)
With the definition $`\mathrm{\Delta }=\mathrm{\Delta }_p\nu \eta _p^2/2`$ the Hamiltonian of the internal states, $`H_I`$, reduces to the one of a driven $`\mathrm{\Lambda }`$-system at the two photon resonance,
$$H_I=\mathrm{}\mathrm{\Delta }|ee|+\frac{\mathrm{}\mathrm{\Omega }_L}{2}\left(\sigma _{er}+\sigma _{re}\right)+\frac{\mathrm{}g}{2}\left(\sigma _{eg}+\sigma _{ge}\right).$$
(34)
In the limit where the external and internal degrees of freedom decouple, $`H_{\mathrm{int}}0`$, the conditioned dynamics determined by Eq. (28) leads to a relaxation of the atom into the state
$$\stackrel{~}{\rho }_c(t\mathrm{})=\stackrel{~}{\mu }_0|DD|,$$
(35)
where $`\stackrel{~}{\mu }_0`$ is an undetermined state of the external degrees of freedom. The interaction with the atomic motion, $`H_{\mathrm{int}}`$, or to be more precise the term $`\mathrm{}\nu \overline{\eta }\widehat{p}|rr|`$ couples the state $`|D`$ to the bright (internal) eigenstates of $`H_I`$, $`|+`$ and $`|`$. They are given by
$`|+`$ $`=`$ $`\mathrm{cos}(\theta )|e+\mathrm{sin}(\theta )\left(g|g+\mathrm{\Omega }_L|r\right)/\mathrm{\Omega },`$
$`|`$ $`=`$ $`\mathrm{sin}(\theta )|e\mathrm{cos}(\theta )\left(g|g+\mathrm{\Omega }_L|r\right)/\mathrm{\Omega },`$
where the mixing angle $`\theta `$ is defined by the relation $`\mathrm{tan}(\theta )=\mathrm{\Omega }/(\sqrt{\mathrm{\Omega }^2+\mathrm{\Delta }^2}\mathrm{\Delta })`$. Theses two states are separated from the dark state by the energies
$$\mathrm{}\mathrm{\Omega }_\pm =\mathrm{}(\mathrm{\Delta }\sqrt{\mathrm{\Omega }^2+\mathrm{\Delta }^2})/2.$$
(36)
Therefore, for small Doppler shifts the population in the bright states is of the order of $`\mathrm{\Delta }_D^2/\mathrm{\Omega }_\pm ^2`$.
In the following we consider the limit where the eigen frequencies of the bright states, $`\mathrm{\Omega }_\pm `$, are much larger than the typical Doppler detuning, $`\mathrm{\Delta }_D`$, as well as the frequency of the trap, $`\nu `$. The first assumption says that the atom is only weakly excited and the total density operator is well approximated by,
$$\stackrel{~}{\rho }_c(t)\stackrel{~}{\mu }_0(t)|DD|.$$
(37)
The second condition, $`\mathrm{\Omega }_L\nu `$, ensures that the internal state of the atom adiabatically follows the evolution of the atomic momentum. If both conditions are satisfied we can adiabatically eliminate the population in the the bright states and derive an effective equation for the evolution of the motional state, $`\stackrel{~}{\mu }_0(t)`$. The details of this calculation are summarized in Appendix A. Finally, we revert the unitary transformation, $`U`$ (29), and trace over the internal states. For the resulting conditioned density operator of the motional state, $`\mu _c:=\mathrm{Tr}_I\{U\stackrel{~}{\rho }_cU^{}\}`$, we obtain the stochastic master equation,
$$\begin{array}{cc}\hfill d\mu _c(t)=& i[\nu \widehat{a}^{}\widehat{a}\lambda ^2\mathrm{\Delta }\widehat{p}^2,\mu _c(t)]dt+_M\mu _c(t)dt\hfill \\ & +\sqrt{ฯต\lambda ^2\mathrm{\Gamma }_g}[\widehat{p}e^{i\varphi }]\mu _c(t)dW(t),\hfill \end{array}$$
(38)
and the homodyne current
$$I_c(t)=2ฯต\lambda ^2\mathrm{\Gamma }_g\mathrm{cos}(\varphi )\widehat{p}_c+\sqrt{ฯต\lambda ^2\mathrm{\Gamma }_g}\xi (t).$$
(39)
In these two equations we defined the parameter, $`\lambda =2g\overline{\eta }\nu /\mathrm{\Omega }^2`$, and the average is take with respect to the conditioned motional state, $`_c=\mathrm{Tr}\{\mu _c\}`$. The measurement back action has the form
$$_M\mu =\frac{\mathrm{\Gamma }_0}{2}\left[\frac{2\stackrel{~}{๐ฅ}_g}{1\stackrel{~}{๐ฅ}_r}(\widehat{p}\mu \widehat{p})\widehat{p}^2\mu \mu \widehat{p}^2\right],$$
(40)
with $`\mathrm{\Gamma }_0=\lambda ^2\mathrm{\Gamma }`$ and
$$\stackrel{~}{๐ฅ}_{j=g,r}(\mu )=\frac{\mathrm{\Gamma }_j}{\mathrm{\Gamma }}_1^1๐uN(u)e^{i\eta _j\widehat{z}(u1)}\mu e^{i\eta _j\widehat{z}(u1)}.$$
(41)
*Discussion*. The results given in the Eqs. (38) and (39) are valid in the limit of a weak probe field, $`g\mathrm{\Omega }_L`$. In that case the dark state, $`|D`$, almost coincides with the ground state, $`|g`$, and the signal strength is maximized for a given strength of the measurement back action, $`\mathrm{\Gamma }_0`$, (see Appendix A). The signal, $`I_c(t)`$, can further be optimized setting $`\varphi =0`$ and by choosing atomic states with a small branching ratio $`\mathrm{\Gamma }_r/\mathrm{\Gamma }_g`$. The remaining difference between $`\lambda ^2\mathrm{\Gamma }_g`$ and $`\mathrm{\Gamma }_0`$ can be absorbed into the definition of $`ฯต`$. For $`\mathrm{\Delta }=0`$ we then end up with Eq. (4) and Eq. (II.2) as given in Section II.2.
## IV Feedback Cooling
As already mentioned in Section II the goal of the continuous momentum observation is to use the information in the signal to manipulate the motion of the atom, e.g., to cool it. In this section we discuss the implementation of the โcold dampingโ feedback strategy. By applying the theory of direct quantum feedback Wiseman and Milburn (1993a) we derive a master equation for the unconditioned state, $`\mu (t)=E[\mu _c(t)]`$.
For the feedback cooling we consider the measurement setup as described in the previous section. In addition, we apply a force on the atom which is proportional but opposite to the measured current. For a single trapped ion, such a feedback loop can be realized by converting the homodyne current into a voltage difference between two trap electrodes. The effect of the feedback loop on the system evolution can be written in the general form,
$$\dot{\mu }_c(t)|_{fb}=I_c(t\tau )๐ฆ\mu _c(t).$$
(42)
The time delay of the feedback loop, $`\tau `$, can usually be neglected compared to the timescale of the atomic motion, $`\nu ^1`$. Nevertheless, for the derivation of the final master equation a finite value of $`\tau `$ is important to obtain the correct operator ordering Wiseman and Milburn (1993a). To implement the idea of โcold dampingโ, we consider the feedback superoperator,
$$๐ฆ\mu =i\frac{G}{2ฯต}[\widehat{z},\mu ],$$
(43)
where $`G`$ denotes the dimensionless gain factor. Note that with this definition of $`๐ฆ`$, the frequency scale of the feedback contribution is again of order $`\mathrm{\Gamma }_0`$.
It has been shown in Ref. Wiseman and Milburn (1993a) that Eq. (42) must be interpreted in the Stratonovich sense. To be compatible with Eq. (II.2) we convert it into the Ito-type equation,
$$\begin{array}{cc}\hfill d\mu _c(t)|_{fb}=& \mathrm{\Gamma }_0\left(2ฯต\widehat{p}_c(t\tau )๐ฆ+\frac{ฯต}{2}๐ฆ^2\right)\mu _c(t)dt\hfill \\ & +\sqrt{ฯต\mathrm{\Gamma }_0}๐ฆ\mu _c(t)dW(t\tau ).\hfill \end{array}$$
(44)
In this form we already see that the noise in the current $`I_c(t)`$ leads to the diffusion term $`๐ฆ^2`$. We can now add Eq. (44) to the conditioned evolution given in Eq. (II.2) and obtain the full stochastic dynamics of the atom under the action of the feedback loop.
To derive a master equation which is independent of the measurement outcome, we perform an ensemble average over the stochastic process, $`\xi (t)`$ and obtain the evolution of the unconditioned density operator, $`\mu =E[\mu _c]`$. When taking the average we must keep in mind that although $`E[dW(t)]=0`$, the Ito increment, $`dW(t\tau )`$, is not independent of $`\mu _c(t)`$, e.g., $`E[dW(t\tau )\mu _c(t)]0`$. The way to perform the average in the limit, $`\tau 0^+`$, can be found in Ref. Wiseman and Milburn (1994). By following this procedure we end up with the master equation (8) given in Section II.
### IV.1 Feedback cooling in the Lamb-Dicke limit
We first look at the solution of Eq. (8) in the Lamb-Dicke limit, $`\eta _p,\eta _L1`$. In this limit the recoil kicks of the emitted photons can be neglected and the backaction of the measurement (40) simplifies to
$$_M\mu =\frac{\mathrm{\Gamma }_0}{2}[\widehat{p},[\widehat{p},\mu ]].$$
(45)
For a harmonic trapping potential, $`H_E=\mathrm{}\nu \widehat{a}^{}\widehat{a}`$, the feedback master equation (8) is then quadratic in the position and the momentum operators. Therefore, the final state is Gaussian and we obtain analytic expressions for the variances of $`\widehat{z}`$ and $`\widehat{p}`$. The resulting steady state energy is given by
$$E=\frac{\mathrm{}\nu }{2}\left(\frac{G\mathrm{\Gamma }_0^2}{2\nu ^2}+\frac{1}{G}+\frac{G}{4ฯต}\right).$$
(46)
The first contribution in the brackets originates from the enhanced uncertainty of the position coordinate as a result of the measurement of the momentum operator. If the measurement strength, $`\mathrm{\Gamma }_0`$, is much smaller than the trap frequency, $`\nu `$, position and momentum coordinates are mixed sufficiently fast and this contribution disappears. In this limit the optimal feedback gain is given by $`G=\sqrt{4ฯต}`$ and we obtain the minimal energy given in Eq. (13).
In Section V we extend the discussion of the feedback cooling in the Lamb-Dicke limit to arbitrary detunings, $`\mathrm{\Delta }`$. Then laser cooling effects play an important role for the final temperatures.
### IV.2 Feedback cooling beyond the Lamb-Dicke limit
When the trapping potential is weak, the extension of the atomic wavepacket can be of the order of the wavelength of the emitted photons. In that case, the energy spacing in the trap, $`\mathrm{}\nu `$, is comparable to the recoil energy, $`E_R`$, and recoil kicks from the emitted photons lead to an additional diffusion of the atomic momentum. Therefore, we must take into account the full expression for the back action term $`_M`$. By expanding Eq. (40) in the parameter $`\mathrm{\Gamma }_r/\mathrm{\Gamma }`$ we can write it as
$$_M\mu =\frac{\mathrm{\Gamma }_0}{2}\left(2\stackrel{~}{๐ฅ}_g\left(\underset{n=0}{\overset{\mathrm{}}{}}\stackrel{~}{๐ฅ}_r^n\right)(\widehat{p}\mu \widehat{p})\widehat{p}^2\mu \mu \widehat{p}^2\right).$$
(47)
The zeroth order term in the sum corresponds to the physical picture where the atom is excited and simply decays back to the ground state, $`|g|D`$. Processes where the atom first decays into the state $`|r`$, is then reexcited again are taken into account by including higher order terms in this sum.
In general, the full expression of $`_M`$ leads to a hierarchy of coupled equations for the moments of $`\widehat{p}`$ which does not break off as in the Lamb-Dicke limit. In the following we restrict our discussion to a finite trapping potential and consider the limit, $`\mathrm{\Gamma }_0\nu `$. As mentioned above, this ensures a mixing of position and momentum coordinates and therefore, an equal reduction of both variances. In this regime non-energy conserving terms can be neglected and we obtain an equation for the mean occupation number
$$\dot{\widehat{n}}=\mathrm{\Gamma }_0[GD]\widehat{n}+\frac{\mathrm{\Gamma }_0}{2}\left[\frac{G^2}{4ฯต}G+1+D\right].$$
(48)
In this expression, the parameter, $`D`$, describes the heating induced by the recoil kicks from the emitted photons. It is given by
$$D=\eta _g^2\stackrel{~}{\alpha }+\frac{\mathrm{\Gamma }_r}{\mathrm{\Gamma }}\left[\frac{\mathrm{\Gamma }_g}{\mathrm{\Gamma }_g\mathrm{\Gamma }_r}(\eta _r^2\stackrel{~}{\alpha }+\eta _r\eta _g)+\frac{\mathrm{\Gamma }_r}{\mathrm{\Gamma }_g\mathrm{\Gamma }_r}\eta _r^2\right],$$
(49)
where $`\stackrel{~}{\alpha }=\frac{1}{2}N(u)(u1)^2๐u=7/10`$. By choosing an appropriate atomic level configuration, $`\mathrm{\Gamma }_g\mathrm{\Gamma }_r`$ and/or $`\eta _g^2\eta _r^2`$, its value is only limited by $`D\eta _g^2\stackrel{~}{\alpha }`$. In this case, and for optimized gain the minimal steady state energy is
$$E_{\mathrm{min}}=\mathrm{}\nu \left(\frac{\eta _g^2\stackrel{~}{\alpha }+\sqrt{4ฯต+\eta _g^4\stackrel{~}{\alpha }^2}}{4ฯต}\right).$$
(50)
This expression shows that the minimal energy changes from the Lamb-Dicke to the non-Lamb-Dicke regime at the parameter values, $`4ฯต\eta _g^2\stackrel{~}{\alpha }`$. In the non-Lamb-Dicke regime, i.e., for weak trapping potential the minimal energy approaches the value,
$$E_{\mathrm{min}}\frac{\stackrel{~}{\alpha }}{ฯต}E_R.$$
(51)
Therefore, the limit for feedback cooling is set by the recoil energy, $`E_R`$, divided by $`ฯต`$, and temperatures well below the Doppler limit, $`k_BT_D=\mathrm{}\mathrm{\Gamma }/2`$, can be reached.
## V Feedback vs. EIT Laser Cooling
In the previous section we focused on the laser detunings $`\mathrm{\Delta }_p\mathrm{\Delta }_L0`$ with the goal to measure the atomic momentum to achieve quantum feedback cooling. As already discussed in Section II.4, in a $`\mathrm{\Lambda }`$-systems EIT laser cooling provides an effective mechanism to cool atoms essentially to the ground state without any further external manipulation. In this section we derive a master equation which describes both effects, feedback cooling and ELC, and discuss the cross-over from pure feedback cooling to ELC.
For the adiabatic elimination of the excited states in Section III.3 and therefore for the validity of the feedback master equation (8) we required, $`\nu \mathrm{\Omega }_\pm `$. This assumption excludes the parameter regime, where ELC achieves the lowest temperatures, $`\nu \mathrm{\Omega }_+`$ Morigi (2003). In this Section we restrict the discussion to the Lamb-Dicke limit, $`\eta _g,\eta _r1`$. This allows us to derive a master equation for the motional state for arbitrary choice of the parameters $`\mathrm{\Delta }`$, $`\mathrm{\Omega }_L`$ and $`\nu `$.
We start with the full model for the three level atom coupled to the radiation field as introduced in Section III. To optimize the feedback cooling effect and to simplify the following discussion we make the assumptions, $`\mathrm{\Gamma }_g=\mathrm{\Gamma }`$, and as in the previous sections, $`g\mathrm{\Omega }_L`$. Under the two photon resonance condition, $`\mathrm{\Delta }_p=\mathrm{\Delta }_L\mathrm{\Delta }`$, the system Hamiltonian, $`H_S`$, given in Eq. (21) can be written as
$$H_S=H_E+H_I+H_\eta .$$
(52)
The Hamiltonian $`H_\eta `$ describes the coupling between external and internal degrees of freedom. Up to first order in the Lamb-Dicke parameters it is given by
$$H_\eta i\mathrm{}\widehat{z}\left[\eta _r\frac{\mathrm{\Omega }_L}{2}(\sigma _{er}\sigma _{re})+\eta _g\frac{g}{2}(\sigma _{eg}\sigma _{ge})\right].$$
(53)
The conditioned dynamics of the full atomic density operator, $`\rho _c(t)`$, is determined by the stochastic master equation (28). As in Section III.3 the goal is to eliminate the internal states and to derive an effective equation for conditioned motional density operator, $`\mu _c(t)`$. The principal strategy is the same: For vanishing Lamb-Dicke parameters the decay of the bright states relaxes the atom into the state, $`\rho _c(t)=\mu _c(t)|DD|`$. The dynamics of $`\mu _c(t)`$ can be derived by including the coupling Hamiltonian $`H_\eta `$ in second order perturbation theory. In contrast to Section III.3 we impose no restrictions on the energies of the bright states, $`\mathrm{\Omega }_\pm `$, which leads to resonant transitions for $`|\mathrm{\Omega }_\pm |\nu `$. Therefore, to guarantee the validity of the perturbation theory we require that $`H_\eta \overline{\eta }g`$ is much smaller than the decay rates of the bright states, $`\mathrm{\Gamma }_+\mathrm{\Gamma }\mathrm{cos}^2(\theta )`$ and $`\mathrm{\Gamma }_{}\mathrm{\Gamma }\mathrm{sin}^2(\theta )`$.
In Appendix B we use the stochastic Schrรถdinger equation formalism for the adiabatic elimination of the internal states. As a result we obtain the conditioned master equation,
$$\begin{array}{cc}\hfill d\mu _c=& i(\nu +\delta )[\widehat{a}^{}\widehat{a},\mu _c]dt\hfill \\ & +A_{}๐[\widehat{a}]\mu _cdt+A_+๐[a^{}]\mu _cdt\hfill \\ & +\sqrt{ฯต\mathrm{\Gamma }_0}[\widehat{C}e^{i\varphi }]\mu _cdW(t).\hfill \end{array}$$
(54)
and the expression for the homodyne current,
$$I_c(t)=ฯต\mathrm{\Gamma }_0\widehat{C}e^{i\varphi }+\widehat{C}^{}e^{i\varphi }_c(t)+\sqrt{ฯต\mathrm{\Gamma }_0}\xi (t).$$
(55)
Here we set $`\delta =\mathrm{\Gamma }_0\mathrm{Im}[I(\nu )+I(\nu )]/4`$ and defined the atomic โjump operatorโ,
$$\widehat{C}=\frac{\sqrt{2}\nu \mathrm{\Gamma }}{\mathrm{\Omega }^2}\left[I(\nu )\widehat{a}+I(+\nu )\widehat{a}^{}\right],$$
(56)
This operator as well as the laser heating and cooling rates, $`A_\pm =\mathrm{\Gamma }_0\mathrm{Re}[I(\pm \nu )]/2`$, depend on the function $`I(\pm \nu )`$, which is defined in Section II, Eq. (10).
As in the previous section we consider a feedback force which is proportional to the measured signal, $`I_c(t)`$. Note that depending on values of $`I(\pm \nu )`$, and the local oscillator phase, $`\varphi `$, the force is proportional to a linear combination of $`\widehat{p}`$ and $`\widehat{z}`$. For the derivation of the feedback master equation we follow the outline given in Section IV and obtain
$$\begin{array}{cc}\hfill \dot{\mu }=& i(\nu +\delta )[\widehat{a}^{}\widehat{a},\mu ]+A_{}๐[\widehat{a}]\mu +A_+๐[a^{}]\mu \hfill \\ & i\mathrm{\Gamma }_0\frac{G}{2}[\widehat{z},\widehat{C}e^{i\varphi }\mu +\mu \widehat{C}^{}e^{i\varphi }]\mathrm{\Gamma }_0\frac{G^2}{8ฯต}[\widehat{z},[\widehat{z},\mu ]].\hfill \end{array}$$
(57)
Under the rotating wave approximation, which is valid for $`\mathrm{\Gamma }_0\nu `$, and by neglecting small shifts of the trap frequency, we end up with master equation (9) given in Section II.
*Discussion:* Fig. 6 shows the dependence of the four different rates $`A_\pm `$, $`A_\pm ^{fb}`$ as a function of the detuning $`\mathrm{\Delta }`$. The cooling and heating rates which originate from the laser interaction, $`A_\pm `$, correspond to the rates for ELC derived in Ref. Morigi (2003). For the parameter regime $`\mathrm{\Omega }>2\nu `$ and for blue detuning, $`\mathrm{\Delta }>0`$, they lead to a minimal temperature for $`\mathrm{\Omega }^24\nu (\nu \mathrm{\Delta })`$. For red detuning, $`\mathrm{\Delta }0`$, the heating rate is larger than the cooling rate and without feedback the system does not reach a steady state. By adjusting the phase $`\varphi `$ the feedback loop always provides additional damping, $`W=A_{}^{fb}A_+^{fb}>0`$, which for $`\mathrm{\Delta }=0`$ and $`\nu \mathrm{\Omega }`$ is given by $`W=G\mathrm{\Gamma }_0`$. The noise added by the feedback loop $`\mathrm{\Gamma }_0G^2/8ฯต`$, imposes a restriction on $`G`$ if one is interested in low steady state energies. The combined effect of feedback and EIT laser cooling lead to a final temperature which is plotted in Fig. 3 in Section II.4.
## VI Conclusion
In this paper we have shown that a continuous readout of the momentum of a single atom can be achieved by employing the high velocity sensitivity of the index of refraction in a driven $`\mathrm{\Lambda }`$-system. The transparency effect for an atom at rest and the linear dependence of the index of refraction on the Doppler shift lead to a homodyne current linear in $`\widehat{p}`$ and a minimal back action on the atomic motion, approaching the quantum limit for $`ฯต1`$.
By applying a force which is proportional to the measured signal, feedback cooling for single ions can be realized. The cooling scheme is applicable in and outside the Lamb-Dicke regime with steady state temperatures well below the Doppler limit. From a fundamental point of view we want to point out, that in the proposed feedback scheme the measured current is fed back *directly* on the trap electrodes. Therefore, its implementation allows for a test of the theory of direct quantum feedback Wiseman and Milburn (1993a) on an individual quantum system close to the ground state.
###### Acknowledgements.
The authors thank G. Morigi, J. Eschner and the experimental group around R. Blatt for useful discussions. Work at Innsbruck was supported in part by the Austrian Science Foundation FWF, European Networks and the Institute for Quantum Information.
## Appendix A Adiabatic elimination
In this appendix we derive the stochastic master equation (38) for the external density operator, $`\mu _c`$, and the expression for the signal $`I_c(t)`$ (39). We start with the conditioned master equation given in Eq. (28) and apply the unitary transformation as defined in Eq. (29). In the new basis the stochastic master equation can be written in the form
$$\begin{array}{c}\hfill d\stackrel{~}{\rho }_c=\left(_I+_\nu +\stackrel{~}{๐ฅ}\right)\stackrel{~}{\rho }_cdt+\sqrt{ฯต\mathrm{\Gamma }_g}[\sigma _{ge}e^{i\varphi }]\stackrel{~}{\rho }_cdW(t),\end{array}$$
(58)
where we divided the total Liouville operator into the three contributions,
$`_I(\rho )`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H_I,\rho ]{\displaystyle \frac{\mathrm{\Gamma }}{2}}(|ee|\rho +\rho |ee|),`$
$`_\nu (\rho )`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}[H_E+H_{\mathrm{int}},\rho ],`$
$`\stackrel{~}{๐ฅ}(\rho )`$ $`=`$ $`\mathrm{\Gamma }\stackrel{~}{๐ฅ_g}(\sigma _{ge}\rho \sigma _{eg})+\mathrm{\Gamma }\stackrel{~}{๐ฅ_r}(\sigma _{re}\rho \sigma _{er}),.`$
The action of the recycling operators, $`\stackrel{~}{๐ฅ}_{g,r}`$, is defined in Section III Eq. (41).
In the following we write the total density operator in terms of the eigenbasis of $`H_I`$ as, $`\stackrel{~}{\rho }_c=_{i,j}\stackrel{~}{\mu }_{ij}|ij|`$ with $`i,j\{+,,D\}`$. Reinserting this decomposition into Eq. (58) we obtain a set of coupled equations for the external operators $`\stackrel{~}{\mu }_{ij}`$. By grouping the 9 elements $`\stackrel{~}{\mu }_{ij}`$ into a single vector $`\stackrel{}{\mu }`$, the resulting set of equations can be written in the form
$$d\stackrel{}{\mu }=(๐_I+๐_\nu +๐)\stackrel{}{\mu }dt+\sqrt{ฯต\mathrm{\Gamma }_g}\left(๐\mathrm{Tr}\{\stackrel{}{V}\stackrel{}{\mu }\}\right)\stackrel{}{\mu }dW(t).$$
(59)
The entries of the matrices $`๐_I`$, $`๐_\nu `$, $`๐`$, $`๐`$ and the vector, $`\stackrel{}{V}`$, can be derived in a straight forward (but lengthy) way by writing the operators in Eq. (58) in terms of the states, $`|+`$, $`|`$ and $`|D`$. Note that the entries of $`๐_\nu `$ and $`๐`$ still contain superoperators acting on the external operators $`\stackrel{~}{\mu }_{ij}`$.
The goal is to derive an effective equation for the population in the dark state, $`\stackrel{~}{\mu }_0:=\stackrel{~}{\mu }_{DD}`$, for the parameter regime $`\nu ,\mathrm{\Delta }_D\mathrm{\Omega }_\pm `$. Formally we can combine both conditions and make a series expansion in the trap frequency, $`\nu `$. According to the structure of $`๐_I`$ we group the external operators $`\stackrel{~}{\mu }_{ij}`$ into the two vectors $`\stackrel{}{\mu }_1:=(\stackrel{~}{\mu }_{D+},\stackrel{~}{\mu }_{+D},\stackrel{~}{\mu }_D,\stackrel{~}{\mu }_D)^T`$ and $`\stackrel{}{\mu }_2:=(\stackrel{~}{\mu }_{++},\stackrel{~}{\mu }_{},\stackrel{~}{\mu }_+,\stackrel{~}{\mu }_+)^T`$. By ordering the entries of $`\stackrel{}{\mu }`$ such that, $`\stackrel{}{\mu }=(\stackrel{}{\mu }_2^T,\stackrel{}{\mu }_1^T,\mu _0)^T`$, the matrices $`๐_I`$, $`๐`$ and $`๐_\nu `$ have the block form
$$๐_I+๐=\left(\begin{array}{ccc}L_I^2+J^2& 0& 0\\ J^{12}& L_I^1& 0\\ J^{02}& 0& 0\end{array}\right),๐_\nu =\left(\begin{array}{ccc}L_\nu ^2& L_\nu ^{21}& 0\\ L_\nu ^{12}& L_\nu ^1& L_\nu ^{10}\\ 0& L_\nu ^{01}& L_\nu ^0\end{array}\right).$$
From the structure of $`๐_I`$ and $`๐_\nu `$ we see that for $`\nu 0`$ we have $`\stackrel{}{\mu }_1๐ช(\nu )`$ and $`\stackrel{}{\mu }_2๐ช(\nu ^2)`$. Therefore, up to second order in $`\nu `$ the equation for $`\stackrel{~}{\mu }_0`$ is
$$\begin{array}{cc}\hfill d\stackrel{~}{\mu }_0=& \left(L_\nu ^0\stackrel{~}{\mu }_0+L_\nu ^{01}\stackrel{}{\mu }_1+J^{02}\stackrel{}{\mu }_2\right)dt\hfill \\ & +\sqrt{ฯต\mathrm{\Gamma }_g}\left(\stackrel{}{v}\stackrel{}{\mu }_1\mathrm{Tr}\{\stackrel{}{v}\stackrel{}{\mu }_1\}\mu _0\right)dW(t),\hfill \end{array}$$
(60)
with
$$\stackrel{}{v}=\frac{\mathrm{\Omega }_L}{\mathrm{\Omega }}(e^{i\varphi }\mathrm{cos}(\theta ),e^{i\varphi }\mathrm{cos}(\theta ),e^{i\varphi }\mathrm{sin}(\theta ),e^{i\varphi }\mathrm{sin}(\theta )).$$
(61)
The equations for $`\stackrel{}{\mu }_1`$ and $`\stackrel{}{\mu }_2`$ are given by
$`\dot{\stackrel{}{\mu }}_1`$ $`=`$ $`L_I^1\stackrel{}{\mu }_1+L_\nu ^{10}\stackrel{~}{\mu }_0+๐ช(\nu ^2),`$
$`\dot{\stackrel{}{\mu }}_2`$ $`=`$ $`\left(L_I^2+J^2\right)\stackrel{}{\mu }_2+L_\nu ^{21}\stackrel{}{\mu }_1+๐ช(\nu ^3).`$
They can be integrated and up to the relevant orders of $`\nu `$ we obtain the formal solution
$`\stackrel{}{\mu }_1`$ $`=`$ $`(L_I^1)^1L_\nu ^{10}\stackrel{~}{\mu }_0+๐ช(\nu ^2),`$
$`\stackrel{}{\mu }_2`$ $`=`$ $`\left(L_I^2+J^2\right)^1L_\nu ^{21}(L_I^1)^1L_\nu ^{10}\stackrel{~}{\mu }_0+๐ช(\nu ^3).`$
Resubstituting these expressions into Eq. (60) the resulting equation can be written in the form
$$\begin{array}{cc}\hfill d\stackrel{~}{\mu }_0=& i[\widehat{h}_{\mathrm{eff}}\stackrel{~}{\mu }_0\stackrel{~}{\mu }_0\widehat{h}_{\mathrm{eff}}^{}]dt+\lambda ^2\mathrm{\Gamma }(\widehat{p}\stackrel{~}{\mu }_0\widehat{p})dt\hfill \\ & +\sqrt{\frac{ฯต\mathrm{\Gamma }_g\lambda ^2\mathrm{\Omega }_L^2}{\mathrm{\Omega }^2}}[\widehat{p}e^{i\varphi }]\stackrel{~}{\mu }_0dW(t),\hfill \end{array}$$
(62)
with $`\lambda =2\overline{\eta }\nu g\mathrm{\Omega }_L/\mathrm{\Omega }^3`$, the non-hermitian operator
$$\widehat{h}_{\mathrm{eff}}=\nu a^{}a\frac{g\lambda }{2}\widehat{p}\lambda ^2\mathrm{\Delta }\widehat{p}^2i\frac{\lambda ^2\mathrm{\Gamma }}{2}\widehat{p}^2.$$
(63)
and the โrecyclingโ term,
$$=\frac{\frac{\mathrm{\Omega }_L^2}{\mathrm{\Omega }^2}\stackrel{~}{๐ฅ}_g+\frac{g^2}{\mathrm{\Omega }^2}\stackrel{~}{๐ฅ}_r}{1\frac{g^2}{\mathrm{\Omega }^2}\stackrel{~}{๐ฅ}_g\frac{\mathrm{\Omega }_L^2}{\mathrm{\Omega }^2}\stackrel{~}{๐ฅ}_r}.$$
(64)
In the formal expression for $``$ the inversion of operator is justified in the limit of a weak probe field, $`g\mathrm{\Omega }_L`$ and $`\mathrm{\Gamma }_r\mathrm{\Gamma }`$. From the stochastic term in Eq. (62) we see that these conditions also maximize the signal strength for a given decoherence rate, $`\lambda ^2\mathrm{\Gamma }`$.
In the original basis the evolution for $`\mu _c`$ is given by the relation,
$$d\mu _c=\mathrm{Tr}_I\{Ud\stackrel{~}{\mu }_0|DD|U^{}\}.$$
(65)
Due to the overlap $`|r|D|^2=g^2/\mathrm{\Omega }^2`$ the action of $`U`$ on the external states reduces to the action of the operator $`\mathrm{exp}(i\overline{\eta }\widehat{z}g^2/\mathrm{\Omega }^2)`$. Therefore, the only effect of the basis transformation is the cancellation of the term $`\lambda g\widehat{p}/2`$ in the effective Hamiltonian, $`\widehat{h}_{\mathrm{eff}}`$.
For the expression of the measured signal $`I_c(t)`$ (39) we can simply repeat the calculations from above. Using the same notation as in Eq. (59) it can be written as
$$I_c(t)=ฯต\mathrm{\Gamma }_g\mathrm{Tr}\{\stackrel{}{V}\stackrel{}{\mu }\}+\sqrt{ฯต\mathrm{\Gamma }_g}\xi (t).$$
(66)
We see that the first term already appeared in the stochastic master equation (59) and can be evaluated along the same lines. By multiplying the resulting expression by a factor $`\lambda `$, we obtain the current given in Eq. (39).
## Appendix B Adiabatic elimination in the Lamb-Dicke limit
In this appendix we derive the conditioned evolution of the external atomic state in the Lamb-Dicke limit, $`\eta _p,\eta _L1`$. We start with the stochastic Schrรถdinger equation for the total wavefunction, $`|\mathrm{\Psi }`$, which includes the state of the atom as well as the state of the electromagnetic environment. For the system Hamiltonian $`H_S`$ (52) and the atom-field interaction described in Section III it is given by
$$\begin{array}{cc}\hfill d|\mathrm{\Psi }=& \left(\frac{i}{\mathrm{}}H_S\frac{\mathrm{\Gamma }}{2}|ee|\right)|\mathrm{\Psi }dt\hfill \\ \hfill +& \sqrt{\mathrm{\Gamma }}_1^1๐u\sqrt{N(u)}e^{i\eta _p\widehat{z}u}\sigma _{ge}๐B_u^{}(t)|\mathrm{\Psi }.\hfill \end{array}$$
(67)
The noise increment operators, $`dB_u^{}(t)`$, fulfill the Ito rules $`dB_u(t)dB_u^{}^{}(t)=\delta (uu^{})dt`$ and correspond to the emission of photons under an angle $`\alpha =\mathrm{arccos}(u)`$ with respect to the z-axis.
We decompose the total wave function in terms of the eigenstates of $`H_I`$ as $`|\mathrm{\Psi }=_i|\mathrm{\Psi }_i|i`$, with $`i=+,,D`$. For vanishing Lamb-Dicke parameters, $`\eta _j0`$, and after some transient deviations the system evolves into the state, $`|\mathrm{\Psi }=|\mathrm{\Psi }_D|D`$. The coupling between the external and internal degrees of freedom, $`H_{\mathrm{int}}`$, leads to finite contributions from the bright states which are of the order of the Lamb-Dicke parameter, $`|\mathrm{\Psi }_\pm ๐ช(\overline{\eta })`$. In the following we treat the coupling to the excited states in perturbation theory to derive an effective equation for $`|\mathrm{\Psi }_D`$ which is valid up to second order in $`\overline{\eta }`$.
In the interaction picture with respect to the external Hamiltonian, $`H_E=\mathrm{}\nu \widehat{a}^{}\widehat{a}`$, the equation for the dark state wave function is
$$d|\mathrm{\Psi }_D\frac{\overline{\eta }g}{2}\widehat{z}(t)|\mathrm{\Psi }_edt+\sqrt{\mathrm{\Gamma }}_1^1๐u\sqrt{N(u)}๐B_u^{}(t)|\mathrm{\Psi }_e,$$
(68)
with $`|\mathrm{\Psi }_e=\mathrm{cos}(\theta )|\mathrm{\Psi }_++\mathrm{sin}(\theta )|\mathrm{\Psi }_{}`$. Note that in this equation we already used the assumption of a weak probe field, and set $`\mathrm{\Omega }=\sqrt{\mathrm{\Omega }_L^2+g^2}\mathrm{\Omega }_L`$. Under the same assumption the equations for the bright states are given by
$$\frac{d}{dt}\left(\begin{array}{c}|\mathrm{\Psi }_+\\ |\mathrm{\Psi }_{}\end{array}\right)=๐\left(\begin{array}{c}|\mathrm{\Psi }_+\\ |\mathrm{\Psi }_{}\end{array}\right)+\frac{\overline{\eta }g}{2}\widehat{z}(t)\left(\begin{array}{c}\mathrm{cos}(\theta )\\ \mathrm{sin}(\theta )\end{array}\right)|\mathrm{\Psi }_D,$$
(69)
where we defined the matrix
$$๐=\left(\begin{array}{cc}i\mathrm{\Omega }_++\frac{\mathrm{\Gamma }}{2}\mathrm{cos}^2(\theta )& \frac{\mathrm{\Gamma }}{4}\mathrm{sin}(2\theta )\\ \frac{\mathrm{\Gamma }}{4}\mathrm{sin}(2\theta )& i\mathrm{\Omega }_{}+\frac{\mathrm{\Gamma }}{2}\mathrm{sin}^2(\theta )\end{array}\right).$$
(70)
Up to first order in $`\overline{\eta }`$, the solution for the excited states is
$$\left(\begin{array}{c}|\mathrm{\Psi }_+\\ |\mathrm{\Psi }_{}\end{array}\right)=\frac{\overline{\eta }g}{2}\left[_{\mathrm{}}^te^{๐(ts)}\widehat{z}(s)๐s\right]\left(\begin{array}{c}\mathrm{cos}(\theta )\\ \mathrm{sin}(\theta )\end{array}\right)|\mathrm{\Psi }_D.$$
(71)
By inserting the time dependence of the position operator, $`\widehat{z}(t)=(\widehat{a}e^{i\nu t}+\widehat{a}^{}e^{i\nu t})/\sqrt{2}`$, we can evaluate this integral and obtain the evolution of the wavefunctions, $`|\mathrm{\Psi }_\pm `$. This solution is then reinserted into Eq. (68) to get the equation of the dark state wave function up to order $`\overline{\eta }^2`$. As long as the final dynamics in the interaction picture is slow compared to the trap frequency, $`\nu `$, we can use the rotating wave approximation and neglect terms proportional to $`e^{\pm i2\nu t}`$. The resulting equation is then given by
$$\begin{array}{cc}\hfill d|\mathrm{\Psi }_D=& \frac{\overline{\eta }^2g^2}{8}\left[\stackrel{~}{I}(\nu )\widehat{a}\widehat{a}^{}+\stackrel{~}{I}(\nu )\widehat{a}^{}\widehat{a}\right]|\mathrm{\Psi }_Ddt\hfill \\ \hfill +& \sqrt{\frac{\overline{\eta }^2g^2}{8}\mathrm{\Gamma }}_1^1๐u\sqrt{N(u)}\stackrel{~}{C}(t)๐B_u^{}(t)|\mathrm{\Psi }_D,\hfill \end{array}$$
(72)
where we set $`\stackrel{~}{C}(t)=(\stackrel{~}{I}(\nu )\widehat{a}(t)+\stackrel{~}{I}(\nu )\widehat{a}^{}(t))`$, and defined the function $`\stackrel{~}{I}(\nu )`$ by
$$\stackrel{~}{I}(\nu ):=(\mathrm{cos}(\theta ),\mathrm{sin}(\theta ))\left[_0^{\mathrm{}}e^{๐\tau }e^{i\nu \tau }๐\tau \right]\left(\begin{array}{c}\mathrm{cos}(\theta )\\ \mathrm{sin}(\theta )\end{array}\right).$$
(73)
Apart from the motional state of the atom, the wavefunction $`|\mathrm{\Psi }_D`$ still incudes the full state of the electromagnetic environment. To obtain the conditioned dynamics for the external density operator, $`\mu _c`$, we first decompose the set of noise increment operators into the two contributions,
$`dB_p^{}(t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{ฯต}}}{\displaystyle _{1ฯต}^1}๐u\sqrt{N(u)}๐B_u^{}(t),`$
$`dB_{}^{}(t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1ฯต}}}{\displaystyle _1^{1ฯต}}๐u\sqrt{N(u)}๐B_u^{}(t).`$
As in Section III the parameter, $`ฯต`$, determines the fraction of the photons which are scattered into the mode of the probe beam. The increment operator $`dB_p^{}(t)`$ obeys the Ito rule, $`dB_p(t)dB_p^{}(t)=dt`$, and corresponds to the emission of photons, which are focused on the detector.
A rigorous, but rather technical way to convert the stochastic Schrรถdinger equation into a stochastic master equation can be found in Ref. Gardiner and Zoller (2000); Goetsch and Graham (1994). A convenient shortcut is, to first derive the unconditioned master equation for $`\mu =\mathrm{Tr}_{EM}\{|\mathrm{\Psi }_D\mathrm{\Psi }_D|\}`$ and then perform a partial โunravellingโ of this master equation as discussed in Ref. Wiseman and Milburn (1993b). For a phase $`\varphi `$ of the local oscillator we finally get the equation
$$\begin{array}{cc}\hfill d\mu _c=& i\frac{\overline{\eta }^2g^2}{8}\mathrm{Im}[\stackrel{~}{I}(\nu )+\stackrel{~}{I}(+\nu )][\widehat{a}^{}\widehat{a},\mu _c]dt\hfill \\ & +\frac{\overline{\eta }^2g^2}{4}\left(\mathrm{Re}[\stackrel{~}{I}(\nu )]๐[\widehat{a}]+\mathrm{Re}[\stackrel{~}{I}(+\nu )]๐[\widehat{a}^{}]\right)\mu _cdt\hfill \\ & +\sqrt{ฯต\mathrm{\Gamma }\frac{\overline{\eta }^2g^2}{8}}[\stackrel{~}{C}(t)e^{i\varphi }]\mu _cdW(t),\hfill \end{array}$$
(74)
The measured current has the form
$$I_c(t)=ฯต\mathrm{\Gamma }\frac{\overline{\eta }^2g^2}{8}\stackrel{~}{C}(t)e^{i\varphi }+\stackrel{~}{C}^{}(t)e^{i\varphi }_c+\sqrt{ฯต\mathrm{\Gamma }\frac{\overline{\eta }^2g^2}{8}}\xi (t).$$
(75)
To make a comparison to results of Section III and Section IV we introduce the rescaled function, $`I(\nu )`$, by setting $`I(\nu )=\stackrel{~}{I}(\nu )\mathrm{\Omega }^4/(8\nu ^2\mathrm{\Gamma })`$. With this definition and a rescaling of the current we finally obtain Eqs. (54) and (55). |
warning/0506/hep-ex0506080.html | ar5iv | text | # Measurements of ๐ต Decays to Two Kaons
The Belle Collaboration
## Abstract
We report measurements of $`B`$ meson decays to two kaons using 253 fb<sup>-1</sup> of data collected with the Belle detector at the KEKB energy-asymmetric $`e^+e^{}`$ collider. We find evidence for signals in $`B^+\overline{K}^0K^+`$ and $`B^0K^0\overline{K}^0`$ with significances of $`3.0\sigma `$ and $`3.5\sigma `$, respectively. (Charge-conjugate modes included) The corresponding branching fractions are measured to be $`(B^+\overline{K}^0K^+)=(1.0\pm 0.4\pm 0.1)\times 10^6`$ and $`(B^0K^0\overline{K}^0)=(0.8\pm 0.3\pm 0.1)\times 10^6`$. These decay modes are examples of hadronic $`bd`$ transitions. No signal is observed in the decay $`B^0K^+K^{}`$ and we set an upper limit of $`3.7\times 10^7`$ at 90% confidence level.
Recent precise measurements of the branching fractions br and partial rate asymmetries acp from the decays $`BK\pi ,\pi \pi `$ provide essential information to understand the $`B`$ decay mechanism and to probe possible contributions from new physics. The rates for these decays constrain the hadronic $`bs`$ and $`bu`$ amplitudes. Here we report results on $`B^0K^0\overline{K}^0`$ and $`B^+\overline{K}^0K^+`$ decays, which are examples of $`bd`$ hadronic transitions. We also discuss a search for $`B^0K^+K^{}`$, which is sensitive to effects of final-state interactions (FSI) fsi . The results are based on a sample of 275 million $`B\overline{B}`$ pairs collected with the Belle detector at the KEKB $`e^+e^{}`$ asymmetric-energy (3.5 on 8 GeV) collider KEKB operating at the $`\mathrm{{\rm Y}}(4S)`$ resonance.
The Belle detector is a large-solid-angle magnetic spectrometer that consists of a silicon vertex detector (SVD), a 50-layer central drift chamber (CDC), an array of aerogel threshold Cherenkov counters (ACC), a barrel-like arrangement of time-of-flight scintillation counters (TOF), and an electromagnetic calorimeter (ECL) comprised of CsI(Tl) crystals located inside a superconducting solenoid coil that provides a 1.5 T magnetic field. An iron flux-return located outside the coil is instrumented to detect $`K_L^0`$ mesons and to identify muons (KLM). The detector is described in detail elsewhere Belle . Two different inner detector configurations were used. For the first sample of 152 million $`B\overline{B}`$ pairs (Set I), a 2.0 cm radius beampipe and a 3-layer silicon vertex detector were used; for the latter 123 million $`B\overline{B}`$ pairs (Set II), a 1.5 cm radius beampipe, a 4-layer silicon detector and a small-cell inner drift chamber were usedUshiroda .
Charged kaons are required to have a distance of closest approach to the interaction point (IP) in the beam direction ($`z`$) of less than 4 cm and less than 0.1 cm in the transverse plane. Charged kaons and pions are identified using $`dE/dx`$ information and Cherenkov light yields in the ACC. The $`dE/dx`$ and ACC information are combined to form a $`K`$-$`\pi `$ likelihood ratio, $`(K/\pi )=_K/(_K+_\pi )`$, where $`_K`$ $`(_\pi )`$ is the likelihood that the track is a kaon (pion). Charged tracks with $`(K/\pi )>0.6`$ are regarded as kaons. Furthermore, charged tracks that are positively identified as electrons or muons are rejected. The electron identification uses information composed of $`E/p`$ and $`dE/dx`$, shower shape, track matching $`\chi ^2`$, and ACC light yields, while information from the KLM, $`dE/dx`$ and ACC are combined to identify muons. The kaon identification efficiency and misidentification rate are determined from a sample of kinematically identified $`D^+D^0\pi ^+,D^0K^{}\pi ^+`$ decays, where the kaons from the $`D`$ decay are selected in the same kinematic region as in $`BK\overline{K}`$ decays. The kaon efficiency is measured to be $`(84.24\pm 0.13)\%`$ for Set I and $`(82.84\pm 0.14)\%`$ for Set II, while the pion-fake-kaon rates are $`(5.40\pm 0.08)\%`$ and $`(6.86\pm 0.11)\%`$, respectively.
Candidate $`K^0`$ mesons are reconstructed through the $`K_S^0\pi ^+\pi ^{}`$ decay. We pair oppositely-charged tracks assuming the pion hypothesis and require the invariant mass of the pair to be within 18 MeV/$`c^2`$ of the nominal $`K_S^0`$ mass. Furthermore, the intersection point of the $`\pi ^+\pi ^{}`$ pair must be displaced from the IP.
Two variables are used to identify $`B`$ candidates: the beam-constrained mass, $`M_{\mathrm{bc}}\sqrt{E_{\text{beam}}^2p_B^2}`$, and the energy difference, $`\mathrm{\Delta }EE_B^{}E_{\text{beam}}^{}`$, where $`E_{\text{beam}}^{}`$ is the run dependent beam energy and $`E_B^{}`$ and $`p_B^{}`$ are the reconstructed energy and momentum of the $`B`$ candidates in the center-of-mass (CM) frame, respectively. Events with $`M_{\mathrm{bc}}>5.20`$ GeV/$`c^2`$ and $`|\mathrm{\Delta }E|<0.3\mathrm{GeV}`$ are selected for analysis.
The dominant background is from $`e^+e^{}q\overline{q}(q=u,d,s,c)`$ continuum events. Event topology and $`B`$ flavor tagging information are used to distinguish between the spherically distributed $`B\overline{B}`$ events and the jet-like continuum backgrounds. We combine a set of modified Fox-Wolfram moments pi0pi0 into a Fisher discriminant. A signal/background likelihood is formed, based on a GEANT-based geant Monte Carlo (MC) simulation, from the product of the probability density function (PDF) for the Fisher discriminant and that for the cosine of the angle between the $`B`$ flight direction and the positron beam. The continuum suppression is achieved by applying a requirement on a likelihood ratio $`=_s/(_s+_{q\overline{q}})`$, where $`_{s(q\overline{q})}`$ is the signal ($`q\overline{q}`$) likelihood. Additional background discrimination is provided by $`B`$ flavor tagging. For each event, the standard Belle flavor tagging algorithm tagging provides a discrete variable indicating the probable flavor of the tagging $`B`$ meson, and a quality $`r`$, a continuous variable ranging from zero for no flavor tagging information to unity for unambiguous flavor assignment. An event with a high value of $`r`$ (typically containing a high-momentum lepton) is more likely to be a $`B\overline{B}`$ event, and a looser $``$ requirement can be applied. We divide the data into $`r>0.5`$ and $`r0.5`$ regions. A selection requirement on $``$ for events in each $`r`$ region of Set I and Set II is applied according to a figure-of-merit defined as $`N_s^{\mathrm{exp}}/\sqrt{N_s^{\mathrm{exp}}+N_{q\overline{q}}^{\mathrm{exp}}}`$, where $`N_s^{\mathrm{exp}}`$ denotes the expected signal yields based on MC simulation and the assumed branching fractions, $`1.0\times 10^6`$, and $`N_{q\overline{q}}^{\mathrm{exp}}`$ denotes the expected $`q\overline{q}`$ yields from sideband data ($`M_{\mathrm{bc}}<5.26`$ GeV/$`c^2`$).
Background contributions from $`\mathrm{{\rm Y}}(4S)B\overline{B}`$ events are investigated using a large MC sample, which includes events from $`bc`$ transitions and charmless decays. After all the selection requirements, no $`B\overline{B}`$ background is found for the $`B^0K^0\overline{K}^0`$ mode. Owing to $`K`$-$`\pi `$ misidentification, large $`B^0K^+\pi ^{}`$ and $`B^+K^0\pi ^+`$ feed-across backgrounds appear in the $`B^0K^+K^{}`$ and $`B^+\overline{K}^0K^+`$ modes, respectively. A small charmless three-body contribution is found at low $`\mathrm{\Delta }E`$ values for these two modes.
The signal yields are extracted by performing unbinned two dimensional maximum likelihood (ML) fits to the ($`M_{\mathrm{bc}}`$, $`\mathrm{\Delta }E`$) distributions. The likelihood for each mode is defined as
$``$ $`=`$ $`\mathrm{exp}({\displaystyle \underset{s,k,j}{}}N_{s,k,j}){\displaystyle \underset{i}{}}({\displaystyle \underset{s,k,j}{}}N_{s,k,j}๐ซ_{s,k,j,i})`$
$`๐ซ_{s,k,j,i}`$ $`=`$ $`P_{s,k,j}(M_{\mathrm{bc}i},\mathrm{\Delta }E_i),`$
where $`s`$ indicates Set I or Set II, $`k`$ distinguishes between events in the $`r<0.5`$ and $`r>0.5`$ regions, $`i`$ is the identifier of the $`i`$-th event, $`P(M_{\mathrm{bc}},\mathrm{\Delta }E)`$ is the two-dimensional PDF of $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$, $`N_j`$ is the number of events for the category $`j`$, which corresponds to either signal, $`q\overline{q}`$ continuum, a feed-across due to $`K`$-$`\pi `$ misidentification, or background from other charmless three-body $`B`$ decays.
All the signal PDFs ($`P_{s,k,j=\mathrm{signal}}(M_{\mathrm{bc}},\mathrm{\Delta }E)`$) are parametrized by a product of a single Gaussian for $`M_{\mathrm{bc}}`$ and a double Gaussian for $`\mathrm{\Delta }E`$ using MC simulations based on the Set I and Set II detector configurations. The same signal PDFs are used for events in the two different $`r`$ regions. Since the $`M_{\mathrm{bc}}`$ signal distribution is dominated by the beam energy spread, we use the signal peak positions and resolutions obtained from $`B^+\overline{D}{}_{}{}^{0}\pi _{}^{+}`$ data ($`\overline{D}{}_{}{}^{0}K_S^0\pi ^+\pi ^{}`$ sub-decay is used for the $`K^0\overline{K}^0`$ mode, while $`\overline{D}{}_{}{}^{0}K^+\pi ^{}`$ is used for the other two modes) with small mode dependent correlations obtained from MC. The MC-predicted $`\mathrm{\Delta }E`$ resolutions are verified using the invariant mass distributions of high momentum $`D`$ mesons. The decay mode $`\overline{D}{}_{}{}^{0}K^+\pi ^{}`$ is used for $`B^0K^+K^{}`$, $`D^+K_S^0\pi ^+`$ for $`B^+K^0\pi ^+`$ and $`\overline{D}{}_{}{}^{0}K_S^0\pi ^+\pi ^{}`$ for $`B^0K^0\overline{K}^0`$. The parameters that describe the shapes of the PDFs are fixed in all of the fits.
The continuum background in $`\mathrm{\Delta }E`$ is described by a linear function while the $`M_{\mathrm{bc}}`$ distribution is parameterized by an ARGUS function $`f(x)=x\sqrt{1x^2}\mathrm{exp}[\xi (1x^2)]`$, where $`x`$ is $`M_{\mathrm{bc}}`$ divided by half of the total center of mass energy argus . Therefore, the continuum PDF is the product of this ARGUS function and the linear function, where the overall normalization, $`\xi `$ and the slope of the linear function are free parameters in the fit. These free parameters are $`r`$-dependent and allowed to be different in Set I and Set II. The background PDFs for charmless three-body $`B`$ decays for the $`K^+K^{}`$ and $`\overline{K}^0K^+`$ modes are each modeled by a smoothed two-dimensional histogram, obtained from a large MC sample. The feed-across backgrounds for these two modes from the $`K^+\pi ^{}`$ and $`K^0\pi ^+`$ events have $`M_{\mathrm{bc}}\mathrm{\Delta }E`$ shapes similar to the signals with the $`\mathrm{\Delta }E`$ peak positions shifted by $`45`$ MeV. The methods to model the $`K^+K^{}`$ and $`\overline{K}^0K^+`$ signal PDFs are also applied to describe the feed-across background.
When likelihood fits are performed, the yield for each background component ($`N_{s,k,j}`$ where $`j=q\overline{q}`$, feed-across, charmless) is allowed to float independently for each $`s`$ (Set I or Set II), and $`k`$ bin (low or high $`r`$ region). For the signal component, the same branching fraction is required by constraining the number of signal events in each $`(s,k)`$ bin using the measured efficiency in the corresponding $`(s,k)`$ bin. Table 1 summarizes the fit results for each mode. We observe $`13.3\pm 5.6\pm 0.6`$ $`K^0K^+`$ and $`15.6\pm 5.8_{0.6}^{+1.1}`$ $`K^0\overline{K}^0`$ signal events with significances of $`3.0\sigma `$ and $`3.5\sigma `$, respectively. The second errors in the yields are the systematic errors from fitting, estimated from the deviations after varying each parameter of the signal PDFs by one standard deviation, and from modeling the three-body background, studied by excluding the low $`\mathrm{\Delta }E`$ region ($`<0.15`$ GeV) and repeating the fit. At each step, the yield deviation is added in quadrature to provide the fitting systematic errors and the statistical significance is computed by taking the square root of the difference between the value of $`2\mathrm{ln}`$ for the best fit value and zero signal yield. The smallest value is chosen to be the significance including the systematic uncertainty.
Figure 1 shows the $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$ projections of the fits after requiring events to have $`|\mathrm{\Delta }E|<0.06`$ GeV and $`5.271`$ GeV/$`c^2<M_{\mathrm{bc}}<5.289`$ GeV/$`c^2`$, respectively. The feed-across yields are $`47.1\pm 8.7`$ in the $`K^+K^{}`$ mode and $`16.4\pm 6.1`$ in the $`K^0K^+`$ mode. The amounts of the feed-across background are consistent with the expectations of 49.1 $`K^+\pi ^{}`$ and 18.8 $`K^0\pi ^+`$ events, based on MC simulation and measured branching fractions hfag . The MC modeling of the requirement on the likelihood ratio, $``$ is investigated using the $`B^+\overline{D}{}_{}{}^{0}\pi _{}^{+}(\overline{D}{}_{}{}^{0}K_S^0\pi ^+\pi ^{}`$ for $`K^0K^0`$ and $`\overline{D}{}_{}{}^{0}K^+\pi ^{}`$ for the others) samples. The obtained systematic errors are $`\pm 2.9\%`$ for $`B^0K^0\overline{K}^0`$ and $`\pm 6.8\%`$ for the other two modes. The systematic error on the charged track reconstruction efficiency is estimated to be around $`1`$% per track using partially reconstructed $`D^{}`$ events. The resulting $`K_S^0`$ reconstruction is verified by comparing the ratio of $`D^+K_S^0\pi ^+`$ and $`D^+K^{}\pi ^+\pi ^+`$ yields with the MC expectation. The resulting $`K_S^0`$ detection systematic error is $`\pm 4.5\%`$. The final systematic errors are then obtained by quadratically summing the errors due to the reconstruction efficiency and the fitting systematics.
With 275 million $`B\overline{B}`$ pairs, we find evidence of $`B^+\overline{K}^0K^+`$ and $`B^0K^0\overline{K}^0`$ with branching fractions $`(B^+\overline{K}^0K^+)=(1.0\pm 0.4\pm 0.1)\times 10^6`$ and $`(B^0K^0\overline{K}^0)=(0.8\pm 0.3\pm 0.1)\times 10^6`$. These are examples of hadronic $`bd`$ transitions. Our measurements are consistent with preliminary results reported by the BaBar collaboration and agree with some theoretical predictions bjorken ; pqcd ; pqcd\_kk ; fleischer ; chiang . It has been suggested that the branching fraction and CP asymmetry of the mode $`B^0K^0\overline{K}^0`$, which originates from the flavor-changing neutral current process $`\overline{b}\overline{d}s\overline{s}`$, may be sensitive to physics beyond the Standard Model fleischer . Measurements with larger statistics are needed for this purpose. No signal is observed in $`B^0K^+K^{}`$ and we set the upper limit of $`3.7\times 10^7`$ at the 90% confidence level, using the Feldman-Cousins approach feldman taking into account both the statistical and systematic errors pole .
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the NII for valuable computing and Super-SINET network support. We acknowledge support from MEXT and JSPS (Japan); ARC and DEST (Australia); NSFC (contract No. 10175071, China); DST (India); the BK21 program of MOEHRD and the CHEP SRC program of KOSEF (Korea); KBN (contract No. 2P03B 01324, Poland); MIST (Russia); MHEST (Slovenia); SNSF (Switzerland); NSC and MOE (Taiwan); and DOE (USA). |
warning/0506/math0506436.html | ar5iv | text | # Truncated microsupport and hyperbolic inequalities
## 1 Introduction and statement of the main results
Let $`X`$ be a real manifold and let $`F`$ denote an object of the derived category of abelian sheaves on X. The microsupport of $`F`$, denoted by $`SS(F)`$, was introduced by M. Kashiwara and P. Schapira (; ), as a subset of the cotangent bundle $`\pi :T^{}XX`$ describing the directions of non propagation for $`F`$. The truncated microsupport of a given degree $`k`$ (or $`k`$-truncated microsupport), $`SS_k(F)`$, defined by the same authors, is only concerned by degrees of cohomology up to the order $`k`$ and allows us to consider some phenomenon of propagation in specific degrees. Such notion is particularly useful in the framework of the theory of linear partial differential equations. More precisely, when $`F`$ is the complex of holomorphic solutions of a coherent module $``$ over the sheaf $`๐_X`$ of holomorphic differential operators on a complex manifold $`X`$, $`SS_k(F)`$ is completely determined as a subset of the characteristic variety $`\text{Char}()`$, which itself coincides with $`SS(F)`$. In the characteristic case, interesting propagation results (cf. , , ) may be obtained with the truncated microsupport. The truncated microsupport and its functorial properties were studied in and .
It is now natural to study the behaviour of $`SS_k(F)`$ under specialization along a submanifold. That is the main purpose of this work, having in scope the application to $`๐`$-modules, specially to holomorphic solutions of induced systems and to real analytic solutions.
Let $`๐ค`$ be a field. Let $`D^b(๐ค_X)`$ denote the bounded derived category of complexes of sheaves of $`๐ค`$-vector spaces.
Let $`M`$ be a submanifold of $`X`$. We shall identify $`T_{T_M^{}X}(T^{}X)`$, $`T^{}(T_MX)`$ and $`T^{}(T_M^{}X)`$ thanks to the Hamiltonian isomorphism. Unless otherwise specified, we shall follow the notations in . In particular, for $`FD^b(๐ค_X)`$, $`\nu _M(F)`$ denotes the specialization of $`F`$ along $`M`$, an object of $`D^b(๐ค_{T_MX}`$) and $`C_{T_M^{}X}(SS_k(F))`$ denotes the normal cone to $`T_M^{}X`$ along $`SS_k(F)`$. For a morphism $`f:YX`$ we shall use $`f^\mathrm{\#}`$, a correspondence which associates conic subsets of $`T^{}Y`$ to conic subsets of $`T^{}X`$ as well as the operation $`\widehat{+}`$ which associates to pairs of conic closed subsets of $`T^{}X`$ conic closed subsets of $`T^{}X`$.
The main result of this work is the following:
###### Theorem 1.1.
Let $`M`$ be a closed submanifold of $`X`$ and let $`FD^b(๐ค_X)`$. Then:
$$SS_k(\nu _M(F))C_{T_M^{}X}(SS_k(F)).$$
One difficulty in its proof is that the use of distinguished triangles is not always convenient because of the shift they introduce. We also needed to deduce a number of further functorial properties. Namely, as an essential step of the proof of this theorem, we obtain the following estimate:
###### Theorem 1.2.
Let $`Y`$ and $`X`$ be real manifolds, let $`f:YX`$ be a morphism and let $`FD^b(๐ค_X)`$. Then
$$SS_k(f^1F)f^\mathrm{\#}(SS_k(F)).$$
Let us denote by $`f_d`$ and $`f_\pi `$ the canonical morphisms ($`f_d`$ was noted by $`{}_{}{}^{t}f_{}^{}`$ in ):
$`f_\pi :X\times _YT^{}YT^{}Y`$ and $`f_d:X\times _YT^{}YT^{}X`$.
Regarding $`f`$ as the composition of a smooth map with a closed embedding, the proof of Theorem 1.2 relies in two steps. The first is to apply Proposition 4.4 of which proves the estimate when $`f`$ is smooth. The second is Proposition 6.1, where we obtain the estimate
$$SS_k(F|_M)j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X),$$
when $`j`$ is a closed embedding.
Remark that, in that case, $`j^\mathrm{\#}(SS_k(F))=j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X).`$
In particular, when $`f`$ is non characteristic with respect to $`F`$, we get
$$SS_k(f^1F)f_df_\pi ^1(SS_k(F)).$$
Namely, when $`M`$ is non characteristic with respect to $`F`$, in other words,
$$SS(F)T_M^{}XT_M^{}M,$$
we have $`SS_k(F)\widehat{+}T_M^{}X=SS_k(F)+T_M^{}X`$ and
$$j_dj_\pi ^1(SS_k(F)+T_M^{}X)=j_dj_\pi ^1(SS_k(F)).$$
Let now $`Y`$ be a complex closed smooth hypersurface of a complex analytic manifold $`X`$ and assume that $`Y`$ is defined as the zero locus of a holomorphic function $`f`$. Let $`\psi _Y`$ denote the functor of nearby cycles associated to $`Y`$. Recall that $`Y`$ may be regarded as a submanifold $`Y^{}`$ of $`T_YX`$ by a canonical section $`s`$ given by $`s`$ such that $`\psi _Y(F)s^1\nu _Y(F)`$.
Then, Theorem 1.1 entails:
###### Corollary 1.3.
Let $`FD^b(๐ค_X)`$. Then
$$SS_k(\psi _Y(F))s_ds_\pi ^1(C_{T_Y^{}X}(SS_k(F))\widehat{+}T_Y^{}^{}(T_YX)).$$
Let us point out that one interesting application of Proposition 6.1 is the new estimate for the $`k`$-truncated microsupport of the tensor product (see Proposition 6.7).
We end this paper with the application of our results to the complex $`F=Rom_{๐_X}(,๐ช_X)`$ of holomorphic solutions of a coherent $`๐_X`$-module $``$ on a complex manifold $`X`$(see Section 6.2). Let $`๐_X`$ be the sheaf of linear partial differential operators of finite order and $`๐ช_X`$ the sheaf of holomorphic functions. Let $`Y`$ be a complex submanifold of $`X`$ and $`j`$ be the embedding of $`Y`$ in $`X`$. We shall denote by $`_Y`$ the induced system, an object of the derived category of left $`๐_Y`$-modules. Recall that, when $``$ is regular in the sense of , $`_Y`$ has coherent cohomology. Let $`\tau :T_YXY`$ be the projection. Still under the assumption that $``$ is regular along $`Y`$, one defines a coherent $`๐_{T_YX}`$-module $`\nu _Y()`$, the specialisation of $``$ along $`Y`$, satisfying a natural isomorphism
$$\nu _Y(Rom_{๐_X}(,๐ช_X))Rom_{๐_{T_YX}}(\nu _Y(),๐ช_{T_YX}).$$
Moreover, if $`Y`$ has codimension $`1`$, one defines the nearby-cycle module $`\psi _Y()`$, a coherent $`๐_Y`$-module, satisfying a natural isomorphism
$$\psi _Y(F)Rom_{๐_Y}(\psi _Y(),๐ช_Y).$$
We refer for the details on these isomorphisms.
Set $`V=SS(F)=\text{Char}()`$ and denote by $`V=_\alpha V_\alpha `$ the (local) decomposition of $`V`$ in its irreducible components. The notion of orthogonality between a submanifold $`Y`$ of $`X`$ and an involutive subvariety $`V`$ of $`T^{}X`$ will be recalled at Section 6.2. We recall in Lemma 6.8 that if $`Y,V`$ are orthogonal and $`V`$ is irreducible, then $`V^{}=j_d(j_\pi ^1(V))`$ is irreducible and $`\pi (V)`$ has the same codimension of $`\pi ^{}(V^{})`$. Here, $`\pi ^{}:T^{}YY`$ denotes the projection.
As a consequence of Theorem 1.1 together with the results of we obtain:
###### Theorem 1.4.
Let $``$ be a coherent $`๐_X`$-module. Then:
$$SS_k(Rom_{\tau ^1๐_X}(\tau ^1,\nu _Y(๐ช_X)))C_{T_Y^{}X}(SS_k(F)).$$
If, moreover, $``$ is regular along $`Y`$ in the sense of we have:
$$SS_k(Rom_{๐_{T_YX}}(\nu _Y(),๐ช_{T_YX}))C_{T_Y^{}X}(SS_k(F)).$$
From the preceding theorem, the results of and Corollary 1.3 we obtain:
###### Corollary 1.5.
Assume that $``$ is regular along $`Y`$ in the sense of . Then
$$SS_k(Rom_{๐_Y}(\psi _Y(),๐ช_Y))s_ds_\pi ^1(C_{T_Y^{}X}(SS_k(F))\widehat{+}T_Y^{}^{}(T_YX)).$$
Furthermore, Proposition 6.1 together with the results of and Theorem 6.7 of entails:
###### Theorem 1.6.
Assume that $``$ is regular along $`Y`$ in the sense of . Then:
$$SS_k(Rom_{๐_Y}(_Y,๐ช_Y))j_dj_\pi ^1(SS_k(F)\widehat{+}T_Y^{}X).$$
If, moreover, $`Y`$ is non characteristic for $``$ and $`Y`$ is orthogonal to each $`V_\alpha `$ such that codim $`\pi (V_\alpha )k`$, the preceding inclusion becomes an equality, for every $`ik`$:
$$SS_i(Rom_{๐_Y}(_Y,๐ช_Y))=j_dj_\pi ^1(SS_i(F)).$$
Recall that M. Kashiwara has proven in that, when $`Y`$ is non characteristic for $``$,
$$SS(Rom_{๐_Y}(_Y,๐ช_Y))=j_dj_\pi ^1(SS(F)).$$
The condition of orthogonality is required in Theorem 1.6 in order to have the analogous equality up to a given degree $`k`$.
Let us now assume that the complex manifold $`X`$ is the complexified of a real analytic manifold $`M`$. Denote by $`๐_M`$ the sheaf of real analytic functions on $`M`$ and by $`j`$ the embedding of $`M`$ in $`X`$.
Another important application of Theorem 1.2 is:
###### Proposition 1.7.
Let $``$ be a coherent $`๐_X`$-module. Then we have the estimate:
$$SS_k(Rom_{๐_X}(,๐_M))j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X).$$
Let $`_M`$ denote the sheaf of Satoโs hyperfuntions on $`M`$. As an immediate consequence of Proposition 1.7 together with Theorem 6.7 of we get:
###### Corollary 1.8.
Let $``$ be an coherent $`๐_X`$-module. Assume that
$$SS_k(F)T_M^{}XT_X^{}X.$$
Then,
$$\tau ^k(Rom_{๐_X}(,๐_M))\tau ^k(Rom_{๐_X}(,_M)).$$
In particular
$$SS_k(Rom_{๐_X}(,๐_M))=SS_k(Rom_{๐_X}(,_M))j_dj_\pi ^1(SS_k(F))$$
$$j_dj_\pi ^1((\underset{codimY_\alpha <k}{}V_\alpha )(\underset{codimY_\alpha =k}{}T_{Y_\alpha }^{}X)).$$
We shall illustrate this corollary with an example (see Example 6.10) of a propagation phenomenon for real analytic solutions of a class of non elliptic differential operators, which, as far as we know, is new.
When $``$ is elliptic, in other words,
$$SS()T_M^{}XT_X^{}X$$
, we get the estimate:
$$\text{For any k},SS_k(Rom_{๐_X}(,๐_M))=SS_k(Rom_{๐_X}(,_M))$$
$$j_dj_\pi ^1((\underset{codimY_\alpha <k}{}V_\alpha )(\underset{codimY_\alpha =k}{}T_{Y_\alpha }^{}X)).$$
We thank M. Kashiwara and P. Schapira for the useful discussions through the preparation of this work.
## 2 Notations
We will mainly follow the notations in .
Let $`X`$ be a real manifold. We denote by $`\tau :TXX`$ the tangent bundle to $`X`$ and by $`\pi :T^{}XX`$ the cotangent bundle. We identify $`X`$ with the zero section of $`T^{}X`$. Given a smooth submanifold $`Y`$ of $`X`$, $`T_YX`$ denotes the normal bundle to $`Y`$ and $`T_Y^{}X`$ the conormal bundle. Given a submanifold $`Y`$ of $`X`$ and a subset $`S`$ of $`X`$ we denote by $`C_Y(S)`$ the normal cone to $`S`$ along $`Y`$, a closed conic subset of $`T_YX`$.
Let $`f:XY`$ be a morphism of manifolds. We denote by
$`f_\pi :X\times _YT^{}YT^{}Y`$ and $`f_d:X\times _YT^{}YT^{}X`$
the associated morphisms.
Given a subset $`A`$ of $`T^{}X`$, we denote by $`A^a`$ the image of $`A`$ by the antipodal map
$$a:(x,\xi )(x;\xi ).$$
The closure of $`A`$ is denoted by $`\overline{A}`$. For a cone $`\gamma TX`$, the polar cone $`\gamma ^{}`$ to $`\gamma `$ is the convex cone in $`T^{}X`$ defined by
$`\gamma ^{}=\{(x;\xi )TX;x\pi (\gamma )`$ and $`v,\xi 0`$ for any $`(x;v)\gamma \}.`$
Given conic subsets $`A`$ and $`B`$ of $`T^{}X`$, the operations $`A+B`$ and $`A\widehat{+}B`$ are defined in and will be recalled in section 3.
Given an open subset $`\mathrm{\Omega }`$ of $`X`$, as in , we denote by $`N^{}(\mathrm{\Omega })`$ the conormal cone to $`\mathrm{\Omega }`$.
When $`X`$ is an open subset of a real finite-dimensional vector space $`E`$ and $`\gamma `$ is a closed convex cone (with vertex at 0) in $`E`$, we denote by $`X_\gamma `$ the open set $`X`$ endowed with the induced $`\gamma `$-topology of $`E`$.
Let $`๐ค`$ be a field. We denote by $`D(๐ค_X)`$ the derived category of complexes of sheaves of $`๐ค`$-vector spaces on $`X`$ and by $`D^b(๐ค_X)`$ the full subcategory of $`D(๐ค_X)`$ consisting of complexes with bounded cohomologies.
For $`k`$, we denote by $`D^k(๐ค_X)`$ (resp. $`D^k(๐ค_X)`$) the full additive subcategory of $`D^b(๐ค_X)`$ consisting of objects $`F`$ satisfying $`H^j(F)=0`$ for any $`j<k`$ (resp. $`H^j(F)=0`$ for any $`j>k`$). The category $`D^{k+1}(๐ค_X)`$ is denoted by $`D^{>k}(๐ค_X)`$.
Given an object $`F`$ of $`D^b(๐ค_X)`$ and a submanifold $`M`$ of $`X`$, $`\nu _M(F)`$ denotes the specialization of $`F`$ along $`M`$, an object of $`D^b(๐ค_{T_MX})`$.
Let $`F`$ be an object of $`D^b(๐ค_X)`$; we denote by $`SS(F)`$ its microsupport, a closed $`^+`$-conic involutive subset of $`T^{}X`$. For $`pT^{}X`$, $`D^b(๐ค_X;p)`$ denotes the localization of $`D^b(๐ค_X)`$ by the full triangulated subcategory consisting of objects $`F`$ such that $`pSS(F)`$.
Let $`X`$ be a finite-dimensional complex manifold. We denote by $`๐ช_X`$ the sheaf of holomorphic functions, by $`๐_X`$ the sheaf of linear holomorphic differential operators of finite order and by $`๐_X()`$ the filtration by the order. Given a coherent $`๐_X`$-module $``$, we denote by $`\text{Char}()`$ its characteristic variety.
Let $`Y`$ be a closed submanifold, let $`\tau `$ be the projection of $`T_YX`$ on $`Y`$ and let $`V_Y^{}`$ denote the V-filtration on $`๐_X`$ with respect to $`Y`$. Let $`๐_{[T_YX]}`$ denote the sheaf of differential operators on $`T_YX`$ with polynomial coefficients with respect to the fibers of $`\tau `$. Let $`\theta `$ denote the Euler operator on $`T_YX`$. Recall that $``$ is regular along $`Y`$ if for any local section $`u`$ of $``$ there exists a non trivial polynomial $`b`$ of degree $`m`$ such that
$$b(\theta )u(V_Y^1(๐_X)๐_X(m))u.$$
Following Kashiwara in , given an appropriate good $`V_Y^{}`$-filtration on $``$, the specialized of $``$ along $`Y`$, $`\nu _Y()`$, is the coherent $`๐_{T_YX}`$-module generated by the associated graded module. When $`Y`$ is a hypersurface, one defines a coherent $`๐_Y`$-module, the nearby-cycles module $`\psi _Y()`$, as the degree zero homogeneous term of that graded module.
## 3 Review on normal cones in cotangent bundles
For the readerโs convenience we shall recall here some operations on conic subsets in cotagent bundles defined on .
Let $`X`$ be a real manifold, $`(x)`$ a system of local coordinates on $`X`$ and denote by $`(x;\xi )`$ the associated coordinates on $`T^{}X`$. Given two conic subsets $`A`$ and $`B`$ of $`T^{}X`$, one defines the sum
$`A+B=\{(x;\xi )T^{}X;\xi =\xi _1+\xi _2`$ for some $`(x;\xi _1)A`$ and $`(x;\xi _2)B\}`$.
When $`A`$ and $`B`$ are closed, $`A\widehat{+}B`$ is the closed conic set containing $`A+B`$, described as follows: $`(x_0;\xi _0)`$ belongs to $`A\widehat{+}B`$ if and only if there exists sequences $`\{(x_n;\xi _n)\}_n`$ in $`A`$ and $`\{(y_n;\eta _n)\}_n`$ in $`B`$ such that:
$$\{\begin{array}{cc}x_n,y_n\underset{๐}{\overset{}{}}x_0,\hfill & \\ \xi _n+\eta _n\underset{๐}{\overset{}{}}\xi _0,\hfill & \\ |x_ny_n||\xi _n|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
Let $`M`$ be a submanifold of $`X`$. Let $`(x^{},x^{\prime \prime })`$ be a system of local coordinates on $`X`$ such that $`M=\{(x^{},x^{\prime \prime });x^{}=0\}`$ and let $`(x^{},x^{\prime \prime };\xi ^{},\xi ^{\prime \prime })`$ denote the associated coordinates on $`T^{}X`$. Given a subset $`\mathrm{\Lambda }`$ of $`T^{}X`$ we describe the normal cone to $`\mathrm{\Lambda }`$ along $`T_M^{}X`$, $`C_{T_M^{}X}(\mathrm{\Lambda })`$, as follows: $`(x_0^{},x_0^{\prime \prime };\xi _0^{},\xi _0^{\prime \prime })C_{T_M^{}X}(\mathrm{\Lambda })`$ if and only if there exist sequences of real positive numbers $`\{c_n\}_n`$ and $`\{(x_n^{},x_n^{\prime \prime };\xi _n^{},\xi _n^{\prime \prime })\}_n`$ in $`\mathrm{\Lambda }`$ such that:
$$\{\begin{array}{cc}(x_n^{},x_n^{\prime \prime };\xi _n^{},\xi _n^{\prime \prime })\underset{๐}{\overset{}{}}(0,x_0^{\prime \prime };\xi _0^{},0),\hfill & \\ c_n(x_n^{},\xi _n^{\prime \prime })\underset{๐}{\overset{}{}}(x_0^{},\xi _0^{\prime \prime }).\hfill & \end{array}$$
Thanks to the Hamiltonian isomorphism, one gets an embedding of $`T^{}M`$ into $`T_{T_M^{}X}T^{}X`$, and, for a conic subset $`\mathrm{\Lambda }`$ of $`T^{}X`$, the set $`T^{}MC_{T_M^{}X}(\mathrm{\Lambda })`$ is described as follows: $`(x_0^{},x_0^{\prime \prime };\xi _0^{},\xi _0^{\prime \prime })T^{}MC_{T_M^{}X}(\mathrm{\Lambda })`$ if and only if there exists a sequence $`\{(x_n^{},x_n^{\prime \prime };\xi _n^{},\xi _n^{\prime \prime })\}_n`$ in $`\mathrm{\Lambda }`$ such that:
$$\{\begin{array}{cc}(x_n^{\prime \prime };\xi _n^{\prime \prime })\underset{๐}{\overset{}{}}(x_0^{\prime \prime };\xi _0^{\prime \prime }),\hfill & \\ |x_n^{}|\underset{๐}{\overset{}{}}0\hfill & \\ |x_n^{}||\xi _n^{}|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
Let $`f:YX`$ be a morphism of manifolds. The notion of $`f^\mathrm{\#}`$, a correspondence introduced in associating conic subsets of $`T^{}Y`$ to conic subsets of $`T^{}X`$, is rather complicated and we refer the reader to for the details. We just recall the following results:
###### Proposition 3.1.
(cf. Proposition 6.2.4 of ) Let $`\mathrm{\Lambda }`$ be a conic subset of $`T^{}X`$.
(i) Assume that $`f:MX`$ is a closed embedding. Then,
$$f^\mathrm{\#}(\mathrm{\Lambda })=T^{}MC_{T_M^{}X}(\mathrm{\Lambda }).$$
(ii) Let $`(x)`$ (resp. $`(y)`$) be a system of local coordinates on $`X`$ (resp. $`Y`$) and let $`(x;\xi )`$ (resp.$`(y;\eta )`$) be the associated coordinates on $`T^{}X`$ (resp. $`T^{}Y`$). Then
$`(y_0;\eta _0)f^\mathrm{\#}(\mathrm{\Lambda })`$ if and only if there exist a sequence $`\{(x_n;\xi _n)\}_n`$ in $`\mathrm{\Lambda }`$, a sequence $`\{y_n\}_n`$ in $`Y`$ such that
$$y_n\underset{๐}{\overset{}{}}y_0,x_n\underset{๐}{\overset{}{}}f(y_0),({}_{}{}^{t}f_{}^{}(y_n)\xi _n)\underset{๐}{\overset{}{}}\eta _0,|x_nf(y_n)||\xi _n|\underset{๐}{\overset{}{}}0.$$
We shall also need the following description of $`j^\mathrm{\#}`$ when $`j`$ is an embedding:
###### Lemma 3.2.
Let $`M`$ be a closed submanifold of $`X`$ and let $`j`$ denote the embedding of $`M`$ in $`X`$. Let $`\mathrm{\Lambda }`$ be a closed conic subset of $`T^{}X`$. Then:
$$j_dj_\pi ^1(\mathrm{\Lambda }\widehat{+}T_M^{}X)=T^{}MC_{T_M^{}X}(\mathrm{\Lambda }),$$
where we identify $`T_{T_M^{}X}T^{}X`$ and $`T^{}T_MX`$ by the Hamiltonian isomorphism.
###### Proof.
It is enough to prove that
$$j_dj_\pi ^1(\mathrm{\Lambda }\widehat{+}T_M^{}X)=j^\mathrm{\#}(\mathrm{\Lambda }).$$
Let $`pj_dj_\pi ^1(\mathrm{\Lambda }\widehat{+}T_M^{}X)`$ and let $`(x^{},x^{\prime \prime })`$ be a system of local coordinates on $`X`$ in a neighborhood of $`p`$ such that $`M=\{(x);x^{}=0\}`$. Let $`(x;\xi )`$ denote the associated coordinates on $`T^{}X`$. Suppose $`p=(x_0^{\prime \prime };\xi _0^{\prime \prime })`$.
Then there exists $`\xi _0^{}`$ such that $`(0,x_0^{\prime \prime };\xi _0^{},\xi _0^{\prime \prime })\mathrm{\Lambda }\widehat{+}T_M^{}X`$. By definition of $`\widehat{+},`$ there exist sequences $`\{(x_n^{},x_n^{\prime \prime };\xi _n^{},\xi _n^{\prime \prime })\}_n`$ in $`\mathrm{\Lambda }`$ and $`\{(0,y_n^{\prime \prime };\eta _n^{},0)\}_n`$ in $`T_M^{}X`$ such that
$$\{\begin{array}{cc}(x_n^{},x_n^{\prime \prime }),(0,y_n^{\prime \prime })\underset{๐}{\overset{}{}}(0,x_0^{\prime \prime }),\hfill & \\ \xi _n^{\prime \prime }\underset{๐}{\overset{}{}}\xi _0^{\prime \prime },\hfill & \\ \xi _n^{}+\eta _n^{}\underset{๐}{\overset{}{}}\xi _0^{},\hfill & \\ |(x_n^{},x_n^{\prime \prime })(0,y_n^{\prime \prime })||(\xi _n^{},\xi _n^{\prime \prime })|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
Hence
$$\{\begin{array}{cc}x_n^{\prime \prime }\underset{๐}{\overset{}{}}x_0^{\prime \prime },\hfill & \\ x_n^{}\underset{๐}{\overset{}{}}0,\hfill & \\ \xi _n^{\prime \prime }\underset{๐}{\overset{}{}}\xi _0^{\prime \prime },\hfill & \\ |x_n^{}||\xi _n|\underset{๐}{\overset{}{}}0\hfill & \end{array}$$
and $`(x_0^{\prime \prime };\xi _0^{\prime \prime })j^\mathrm{\#}(\mathrm{\Lambda })`$.
Conversely, let $`pj^\mathrm{\#}(\mathrm{\Lambda })`$, $`p=(x_0^{\prime \prime };\xi _0^{\prime \prime })`$. Then there exists a sequence $`\{(x_n^{},x_n^{\prime \prime };\xi _n^{},\xi _n^{\prime \prime })\}_n`$ in $`\mathrm{\Lambda }`$ such that
$$\{\begin{array}{cc}(x_n^{\prime \prime };\xi _n^{\prime \prime })\underset{๐}{\overset{}{}}(x_0^{\prime \prime };\xi _0^{\prime \prime }),\hfill & \\ x_n^{}\underset{๐}{\overset{}{}}0,\hfill & \\ |x_n^{}||\xi _n|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
The sequences $`\{(x_n^{},x_n^{\prime \prime };\xi _n^{},\xi _n^{\prime \prime })\}_n`$ in $`\mathrm{\Lambda }`$ and $`\{(0,x_n^{\prime \prime };\xi _n^{},0)\}`$ in $`T_M^{}X`$ satisfy the necessary conditions so that $`(0,x_0^{\prime \prime };0,\xi _0^{\prime \prime })\mathrm{\Lambda }\widehat{+}T_M^{}X`$, hence $`(x_0^{\prime \prime };\xi _0^{\prime \prime })j_dj_\pi ^1(\mathrm{\Lambda }\widehat{+}T_M^{}X)`$. q.e.d.
###### Lemma 3.3.
Let $`\mathrm{\Lambda }`$ be a closed conic subset of $`T^{}X`$ and $`M`$ a closed submanifold of $`X`$. One has:
$$(\mathrm{\Lambda }\widehat{+}T_M^{}X)\widehat{+}T_M^{}X=\mathrm{\Lambda }\widehat{+}T_M^{}X.$$
###### Proof.
Let $`(x^{},x^{\prime \prime })`$ be a system of local coordinates on $`X`$ such that $`M=\{(x^{},x^{\prime \prime });x^{}=0\}`$ and let $`(x^{},x^{\prime \prime };\xi ^{},\xi ^{\prime \prime })`$ be the associated coordinates on $`T^{}X`$.
Let $`(x_0,;\xi _0)(\mathrm{\Lambda }\widehat{+}T_M^{}X)\widehat{+}T_M^{}X`$, then there exists sequences $`\{(x_n;\xi _n)\}_n`$ and $`\{(y_n;\eta _n)\}_n`$ in $`\mathrm{\Lambda }\widehat{+}T_M^{}X`$ and $`T_M^{}X`$, respectively, such that
$$\{\begin{array}{cc}x_n,y_n\underset{๐}{\overset{}{}}x_0,\hfill & \\ \xi _n+\eta _n\underset{๐}{\overset{}{}}\xi _0,\hfill & \\ |x_ny_n||\xi _n|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
For each $`n`$, since $`(x_n;\xi _n)\mathrm{\Lambda }\widehat{+}T_M^{}X`$ , there exist sequences $`\{(x_m^n;\xi _m^n)\}_m`$ in $`\mathrm{\Lambda }`$ and $`\{(y_m^n;\eta _m^n)\}_m`$ in $`T_M^{}X`$ such that
$$\{\begin{array}{cc}x_m^n,y_m^n\underset{๐}{\overset{}{}}x_n,\hfill & \\ \xi _m^n+\eta _m^n\underset{๐}{\overset{}{}}\xi _n,\hfill & \\ |x_m^ny_m^n||\xi _m^n|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
Hence we can find subsequences $`\{(x_k;\xi _k)\}_k`$ and $`\{(y_k;\eta _k))\}_k`$ of $`\{(x_m^n;\xi _m^n)\}_{n,m}`$ and $`\{(y_m^n;\eta _m^n+\eta _n)\}_{m,n}`$, respectively, such that
$$\{\begin{array}{cc}x_k,y_k\underset{๐}{\overset{}{}}x_0,\hfill & \\ \xi _k+\eta _k\underset{๐}{\overset{}{}}\xi _0,\hfill & \\ |x_ky_k||\xi _k|\underset{๐}{\overset{}{}}0,\hfill & \end{array}$$
which gives $`(x_0;\xi _0)\mathrm{\Lambda }\widehat{+}T_M^{}X`$.
Conversely, since $`\pi (\mathrm{\Lambda }\widehat{+}T_M^{}X)M`$, $`(\mathrm{\Lambda }\widehat{+}T_M^{}X)(\mathrm{\Lambda }\widehat{+}T_M^{}X)+T_M^{}X(\mathrm{\Lambda }\widehat{+}T_M^{}X)\widehat{+}T_M^{}X.`$ q.e.d.
Let us now assume that $`X`$ is an open subset of $`^n`$ with the coordinates $`(x)=(x_1,\mathrm{},x_n)`$ and that $`M`$ is the submanifold $`\{(x^{},x^{\prime \prime })X;(x^{})=(x_1,\mathrm{},x_d)=0\}`$. Let $`\delta >0`$ and let $`\gamma `$ be the closed convex proper cone given by:
$$\gamma =\{(x^{},x^{\prime \prime });x_n\frac{1}{\delta }|(x^{},x_{d+1},\mathrm{}x_{n1})|\}.$$
Hence
$$\gamma ^a=\{(\xi ^{},\xi ^{\prime \prime });\xi _n\delta |(\xi ^{},\xi _{d+1},\mathrm{}\xi _{n1})|\}.$$
Therefore $`(x+\gamma )M=x+(\gamma M)`$, for each $`xM`$. Let $`^+`$ denote the set of real positive numbers and let us introduce the following notation: for any $`\lambda ^+`$
$$\gamma _\lambda =\{(x^{},x^{})X;(\lambda ^1x^{},x^{\prime \prime })\gamma \}$$
$$V_\lambda =\{(x^{},x^{\prime \prime });(\lambda ^1x^{},x^{\prime \prime })V\}.$$
Remark that if $`\lambda <1`$, $`\text{Int}(\gamma _\lambda ^a)\gamma ^a`$.
###### Lemma 3.4.
Let $`\mathrm{\Lambda }`$ be a conic closed subset of $`T^{}X`$.
Let $`x^{\prime \prime }M\pi (\mathrm{\Lambda })`$ and assume that there is a compact neighborhood $`V`$ of $`x^{\prime \prime }`$ such that
$$(V\times \gamma ^a)(\mathrm{\Lambda }\widehat{+}T_M^{}X)T_X^{}X.$$
Then, there exists a real positive number $`C`$ such that for any $`\lambda `$ and $`ฯต`$ satisfying $`0<\lambda ,ฯต<C,`$
$$(V_{\lambda ฯต}\times \gamma _\lambda ^a)\mathrm{\Lambda }T_X^{}X.$$
###### Proof.
We shall argue by contradiction. Therefore, we can find sequences $`(\lambda _l)_l,(ฯต_l)_l`$ of positive numbers converging to $`0`$, $`(x_l^{},x_l^{\prime \prime };\xi _l^{},\xi _l^{\prime \prime })_t`$ in $`\mathrm{\Lambda }`$, $`(\xi _l^{},\xi _l^{\prime \prime })(0,0)`$, such that $`|x_l^{}|C^{}ฯต_l\lambda _l`$, for some positive constant $`C^{}`$ only depending on $`V`$, and $`(0,x_l^{\prime \prime };\lambda _l\xi _l^{},\xi _l^{\prime \prime })V\times \gamma ^a`$.
Since the $`n`$-component $`(\xi _l)_l`$ is positive, after dividing $`(\xi _l^{},\xi _l^{\prime \prime })`$ by $`\xi _{l,n}`$, we may assume that $`\xi _{l,n}=1`$ and that $`(\lambda _l\xi _l^{},\xi _l^{\prime \prime })`$ is a bounded sequence. Since $`|\xi _l^{}||x_l^{}|C^{}ฯต_l\lambda _l|\xi _l^{}|`$ we get that $`(|\xi _l^{}||x_l^{}|)_l`$ converges to $`0`$. Moreover, since $`x_l^{\prime \prime }`$ is bounded, we may assume that $`x_l^{\prime \prime }`$ converges to some $`\stackrel{ห}{x^{\prime \prime }}VM`$ and that $`(\lambda _l\xi _l^{},\xi _l^{\prime \prime })`$ converges to some $`(\xi _0^{},\xi _0^{\prime \prime })\gamma ^a`$, with $`\xi _{0n}=1`$. Considering the sequences $`(x_l^{},x_l^{\prime \prime };\xi _l^{},\xi _l^{\prime \prime })_l\mathrm{\Lambda }`$ and $`(0,x_l^{\prime \prime };\xi _l^{}+\lambda _l\xi _l^{},0)T_M^{}X`$ we get that $`(0,\stackrel{ห}{x^{\prime \prime }};\xi _0^{},\xi _0^{\prime \prime })(V\times \gamma ^a)(\mathrm{\Lambda }\widehat{+}T_M^{}X)`$, which entails $`\xi _{0n}=0`$, a contradiction. q.e.d.
Let $`\mathrm{\Omega }`$ be an open subset of $`X`$. We shall now recall the notion of conormal cone to $`\mathrm{\Omega }`$, $`N^{}(\mathrm{\Omega })`$. It is the subset of $`T^{}X`$ defined as follows:
Given $`xX`$, we denote by $`N_x(\mathrm{\Omega })`$ the subset of $`T_xX`$ consisting of vectors $`v0`$ such that, in a local chart in a neighborhood of $`x`$, there exist an open cone $`\gamma `$ containing $`v`$ and a neighborhood $`U`$ of $`x`$ such that
$$U((\mathrm{\Omega }U)+\gamma )\mathrm{\Omega }.$$
Note that, in particular, $`N_x(\mathrm{\Omega })=T_xX`$ if and only if $`x\overline{\mathrm{\Omega }}`$ or $`x\mathrm{\Omega }`$. We denote by $`N(\mathrm{\Omega })`$ the open convex cone of $`TX`$:
$$N(\mathrm{\Omega })=\underset{xX}{}N_x(\mathrm{\Omega }),$$
and call it the strict normal cone to $`\mathrm{\Omega }`$.
Finally $`N^{}(\mathrm{\Omega })`$, the conormal cone to $`\mathrm{\Omega }`$, is given by
$$N^{}(\mathrm{\Omega })=\underset{xX}{}(N_x^{}(\mathrm{\Omega })),$$
where, for each $`x\mathrm{\Omega }`$, $`N_x^{}(\mathrm{\Omega })=(N_x(\mathrm{\Omega }))^{}`$.
## 4 Review on the truncated microsupport
We shall now recall equivalent definitions of the truncated microsupport, following .
Given $`(x_0,\xi _0)^n\times (^n)^{}`$ and $`\epsilon `$ we set:
$$H_\epsilon (x_0,\xi _0)=\{x^n;xx_0,\xi _0>\epsilon \},$$
and if there is no risk of confusion we will write $`H_\epsilon `$ instead of $`H_\epsilon (x_0,\xi _0)`$.
###### Proposition 4.1.
Let $`X`$ be a real analytic manifold and let $`pT^{}X`$. Let $`FD^b(๐ค_X)`$, $`k`$ and $`\alpha \{\mathrm{},\omega \}`$. Then the following conditions are equivalent:
$`(i)_k`$ There exist $`F^{}D^{>k}(๐ค_X)`$ and an isomorphism $`FF^{}`$ in $`D^b(๐ค_X;p)`$;
$`(ii)_k`$ There exist $`F^{}D^{>k}(๐ค_X)`$ and a morphism $`F^{}F`$ in $`D^b(๐ค_X)`$ which is an isomorphism in $`D^b(๐ค_X;p)`$;
$`(iii)_{k,\alpha }`$ There exists an open conic neighborhood $`U`$ of $`p`$ such that for any $`x\pi (U)`$ and any $``$-valued $`C^\alpha `$-function $`\phi `$ defined on a neighborhood of $`x`$ such that $`\phi (x)=0`$, $`d\phi (x)U`$, one has
$`H_{\{\phi 0\}}^j(F)_x=0,`$ for any $`jk.`$
When $`X`$ is an open subset of $`^n`$ and $`p=(x_0,\xi _0)`$, the above conditions are also equivalent to:
$`(iv)_k`$ There exist a proper closed convex cone $`\gamma ^n`$, $`\epsilon >0`$ and an open neighborhood $`W`$ of $`x_0`$ with $`\xi _0Int(\gamma ^{})`$ such that $`(W+\gamma ^a)\overline{H_\epsilon }X`$ and
$`H^j(X;๐ค_{(x+\gamma ^a)H_\epsilon }F)=0,`$ for any $`jk,xW.`$
###### Remark 4.2.
Note that when $`X`$ is an open subset of $`^n`$ and $`p=(x_0,\xi _0)`$, the equivalent conditions of Proposition 4.1 are also equivalent to:
There exists some $`F^{}D^b(๐ค_X)`$ isomorphic to $`F`$ in a neighborhood of $`x_0`$ and a closed proper convex cone $`\gamma `$ in $`E`$, with $`0\gamma `$ and $`\xi _0Int\gamma ^a`$, such that $`R\varphi _\gamma (F^{})D^{>k}(๐ค_{X_\gamma })`$.
###### Definition 4.3.
Let $`FD^b(๐ค_X)`$. We define the closed conic subset $`SS_k(F)`$ of $`T^{}X`$ by: $`pSS_k(F)`$ if and only if $`F`$ satisfies the equivalent conditions in the preceding Proposition.
We shall need the following properties of the truncaded microsupport also proved in :
(i) Given a distinguished triangle $`F^{}FF^{\prime \prime }\stackrel{+1}{}`$, one has
$$SS_k(F)SS_k(F^{})SS_k(F^{\prime \prime }),$$
(1)
$$(SS_k(F^{})\backslash SS_{k1}(F^{\prime \prime }))(SS_k(F^{\prime \prime })\backslash SS_{k+1}(F^{}))SS_k(F).$$
(2)
(ii) For any $`FD^b(๐ค_X)`$, one has
$$SS_k(F)T_X^{}X=\pi (SS_k(F))=supp(\tau ^k(F)).$$
(3)
###### Proposition 4.4.
Let $`X`$ and $`Y`$ be two manifolds. Then for $`FD^b(๐ค_X)`$, $`GD^b(๐ค_Y)`$ and $`k`$, one has:
$$SS_k(FG)=\underset{i+j=k}{}SS_i(F)\times SS_j(G).$$
###### Proposition 4.5.
Let $`Y`$ and $`X`$ be two manifolds, let $`f:YX`$ be a morphism and let $`GD^b(๐ค_Y)`$ such that $`f`$ is proper on the support of $`G`$. Then for any $`k`$,
$$SS_k(Rf_{}(G))f_\pi f_d^1(SS_k(G)).$$
(4)
The equality holds in the case f is a closed embedding.
###### Proposition 4.6.
Let $`Y`$ and $`X`$ be two manifolds and let $`f:YX`$ be a smooth morphism. Let $`FD^b(๐ค_X)`$. Then
$$SS_k(f^1F)=f_df_\pi ^1(SS_k(F)).$$
(5)
To end this section, we shall prove the following characterizations of the truncated microsupport not included in , which will be useful in the sequel.
###### Lemma 4.7.
Let $`E`$ be a real finite-dimensional vector space, $`X`$ an open subset of $`E`$ and let $`FD^b(๐ค_X)`$. Let $`U`$ be an open subset of $`X`$ and $`\gamma `$ be a closed convex proper cone in $`E`$ with $`0\gamma `$. Assume that
$$SS_k(F)(U\times \text{Int}(\gamma ^a))=\mathrm{}.$$
Then, given $`(x_0,\xi _0)U\times Int(\gamma ^a)`$, $`\epsilon >0`$ and an open subset $`WX`$ such that $`(W+\gamma )H_\epsilon U`$, one has:
$$H^j(X;๐ค_{(x+\gamma )H_\epsilon }F)=0,foranyxW+\gamma andjk.$$
(6)
###### Proof.
We may assume that $`X`$ is an open subset of $`^n`$.
Let $`(x_0;\xi _0)U\times Int(\gamma ^a)`$, $`\epsilon >0`$ and $`WX`$ be an open subset such that $`(W+\gamma )H_\epsilon U`$. Let us prove (6).
By the microlocal cut-off lemma (Proposition 5.2.3 of ), we have a distinguished triangle
$$\varphi _\gamma ^1R\varphi _\gamma FFG\stackrel{+1}{},$$
with $`SS(G)(X\times \text{Int}(\gamma ^a))=\mathrm{}`$. Therefore, setting $`F^{}=\varphi _\gamma ^1R\varphi _\gamma F`$, one has
$$H^j(X;๐ค_{(x+\gamma )H_\epsilon }F)H^j(X;๐ค_{(x+\gamma )H_\epsilon }F^{}),$$
for any $`xW+\gamma `$ and $`j`$, and $`SS_k(F^{})(U\times \text{Int}(\gamma ^a))=\mathrm{}`$. Hence we may replace $`F`$ by $`F^{}`$ to prove condition (6).
Arguing by induction on $`k`$, we may assume that (6) holds for $`k1`$ and hence $`FD^k(๐ค_X)`$. Hence, given $`xW+\gamma `$,
$$H^k(X;๐ค_{(x+\gamma )H_\epsilon }F)\mathrm{\Gamma }(X;๐ค_{(x+\gamma )H_\epsilon }H^k(F)).$$
Given $`s\mathrm{\Gamma }(X;๐ค_{(x+\gamma )H_\epsilon }H^k(F))`$ we can extend $`s`$ to a section
$$\stackrel{~}{s}\mathrm{\Gamma }(\mathrm{\Omega };๐ค_{H_\epsilon }H^k(F))\mathrm{\Gamma }(\mathrm{\Omega };H^k(F)),$$
where $`\mathrm{\Omega }`$ is a $`\gamma `$-open neighborhood of $`x+\gamma `$ such that $`\mathrm{\Omega }H_\epsilon U`$.
Set $`S=`$ supp$`(\stackrel{~}{s})\mathrm{\Omega }H_\epsilon `$. Since $`H_{\{\phi 0\}}^k(F)\mathrm{\Gamma }_{\{\phi 0\}}(H^k(F))`$, for any real analytic function $`\phi `$ defined on $`^n`$, we get $`S=\mathrm{}`$ from the following Lemma, and hence $`H^k(X;๐ค_{(x+\gamma )H_\epsilon }F)=0.`$ q.e.d.
###### Lemma 4.8 ).
Let $`\gamma `$ be a proper closed convex cone in $`^n`$. Let $`\mathrm{\Omega }`$ be a $`\gamma `$-open subset of $`^n`$ and let $`S`$ be a closed subset of $`\mathrm{\Omega }`$ such that $`S^n`$. Assume the following condition: for any $`x^n`$ and any real analytic function $`\phi `$ defined on $`^n`$, the three conditions $`S\{x;\phi (x)<0\}=\mathrm{}`$, $`\phi (x)=0`$ and $`d\phi (x)Int(\gamma ^a)`$ imply $`xS`$. Then $`S`$ is an empty set.
###### Corollary 4.9.
Let $`E`$ be a real finite dimensional vector space, $`X`$ an open subset of $`E`$ and let $`FD^b(๐ค_X)`$. Let $`U`$ be an open subset of $`X`$ and $`\gamma `$ be a closed convex proper cone in $`E`$ with $`0\gamma `$. Assume that
$$SS_k(F)(U\times \text{Int}(\gamma ^a))=\mathrm{}.$$
Then, for each $`(x_0,\xi _0)U\times \text{Int}(\gamma ^a)`$ there exists an open neighborhood $`V`$ of $`x_0`$ in $`U`$ such that
$$R\varphi _\gamma (R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(F))D^{>k}(๐ค_{X_\gamma })$$
for every $`\gamma `$-open subsets $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_0`$ with $`\mathrm{\Omega }_0\mathrm{\Omega }_1`$, $`\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0V`$ and $`x_0\text{Int}(\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0)`$.
###### Proof.
We may assume assume $`X=^n`$. Let us consider an open neighborhood $`V^{}`$ of $`x_0`$ such that $`\overline{V^{}}`$ is a compact subset of $`U`$.
For each $`x\overline{V^{}}`$ there exists $`\epsilon (x)>0`$ and an open neighborhood $`V(x)`$ of $`x`$ such that $`(V(x)+\gamma )H_{\epsilon (x)}U`$. Since $`\overline{V^{}}`$ is compact, we can find a finite covering of $`\overline{V^{}}`$ by open subsets, let us say, $`\{V(x_1),\mathrm{},V(x_l)\}`$, with $`l`$. Let us choose $`\epsilon =min\{\epsilon (x_i);i=1,\mathrm{},l\}`$. Then, $`(V^{}+\gamma )H_\epsilon U`$, and by Lemma 4.7,
$`H^j(X;๐ค_{(x+\gamma )H_\epsilon }F)=0`$, for each $`xV^{}`$ and $`jk`$.
Since
$$H^j(X;๐ค_{(x+\gamma )H_\epsilon }F)H^j(\varphi _\gamma ^1R\varphi _\gamma F_{H_\epsilon })_x,$$
for all $`xV^{}`$ and $`jk`$, one gets $`\varphi _\gamma ^1R\varphi _\gamma F_{H_\epsilon }D^{>k}(๐ค_{V_\gamma ^{}})`$.
Let $`V=V^{}H_\epsilon `$ and consider $`\gamma `$-open subsets $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_0`$ such that $`\mathrm{\Omega }_0\mathrm{\Omega }_1`$, $`x_0\text{Int}(\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0)`$ and $`\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0V`$.
From $`\varphi _\gamma ^1R\varphi _\gamma F_{H_\epsilon }D^{>k}(๐ค_{V_\gamma })`$ one gets
$$R\varphi _\gamma (R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(\varphi _\gamma ^1R\varphi _\gamma F_{H_\epsilon }))D^{>k}(๐ค_{X_\gamma }).$$
Since
$$R\varphi _\gamma (R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(F_{H_\epsilon }))R\varphi _\gamma (R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(\varphi _\gamma ^1R\varphi _\gamma F_{H_\epsilon })),$$
we obtain $`R\varphi _\gamma (R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(F_{H_\epsilon }))D^{>k}(๐ค_{X_\gamma })`$.
Finally, we conclude that $`R\varphi _\gamma (R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(F))D^{>k}(๐ค_{X_\gamma }),`$ since $`R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(F_{H_\epsilon })`$ is isomorphic to $`R\mathrm{\Gamma }_{\mathrm{\Omega }_1\backslash \mathrm{\Omega }_0}(F)`$.
q.e.d.
## 5 Complements on functorial properties of the truncated microsupport
In order to prove the main results we need further functorial properties of the truncated microsupport similar to those of the microsupport itself but requiring adapted proofs.
###### Lemma 5.1.
Let $`X`$ be a finite dimensional real vector space, $`\gamma `$ be a closed convex proper cone of $`X`$ with $`0\gamma `$, and $`\mathrm{\Omega }`$ a $`\gamma ^a`$-open subset of $`X`$ such that, for any compact $`K`$ of $`X`$, $`\mathrm{\Omega }(K+\gamma )`$ is relatively compact. Let $`FD^b(๐ค_X)`$ and assume $`R\varphi _{\gamma }^{}{}_{}{}^{}FD^{>k}(๐ค_{X_\gamma })`$. Then we have
$$R\varphi _{\gamma }^{}{}_{}{}^{}F_\mathrm{\Omega }D^{>k}(๐ค_{X_\gamma }).$$
###### Proof.
The proof is contained in the proof of Lemma 5.4.3 (i) of . q.e.d.
###### Proposition 5.2.
Let $`X`$ be a manifold, $`FD^b(๐ค_X)`$ and $`\mathrm{\Omega }`$ be an open subset of $`X`$.
(i) Assume $`SS_k(F)N^{}(\mathrm{\Omega })^aT_X^{}X`$. Then
$$SS_k(R\mathrm{\Gamma }_\mathrm{\Omega }(F))N^{}(\mathrm{\Omega })+SS_k(F).$$
(ii) Assume $`SS_k(F)N^{}(\mathrm{\Omega })T_X^{}X`$. Then
$$SS_k(F_\mathrm{\Omega })N^{}(\mathrm{\Omega })^a+SS_k(F).$$
###### Proof.
The proof is an adaptation of the proof of Proposition 5.4.8 (i) and (ii) of , using Corollary 4.9 and Remark 4.2 instead of Propositions 5.2.1 and 5.1.1 of , respectively. q.e.d.
###### Proposition 5.3.
Let $`\mathrm{\Omega }`$ be an open subset of $`X`$ and let $`j`$ be the embedding $`\mathrm{\Omega }X`$. Let $`FD^b(๐ค_\mathrm{\Omega })`$. Then:
(i) $`SS_k(Rj_{}F)SS_k(F)\widehat{+}N^{}(\mathrm{\Omega }).`$
(ii) $`SS_k(Rj_!F)SS_k(F)\widehat{+}N^{}(\mathrm{\Omega })^a.`$
###### Proof.
The proof is the stepwise adaptation of the proof of Proposition 6.3.1 of , using Propositions 5.2, 4.1 and Corollary 4.9 instead of Propositions 5.4.8, 5.1.1 and 5.2.1 of , respectively. q.e.d.
## 6 Proofs of the main results
### 6.1 Proofs of Theorems 1.1, 1.2 and Corollaries
*Proof of Theorem 1.2* Let us first consider the case of the embedding of a closed submanifold of $`X`$:
###### Proposition 6.1.
Let $`M`$ be a closed submanifold of $`X`$ and $`FD^b(๐ค_X)`$. Then
$$SS_k(F|_M)j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X),$$
where $`j`$ is the embedding of $`M`$ in $`X`$.
###### Proof.
Let $`d`$ denote the codimension of $`M`$. Let $`(x_1,\mathrm{},x_n)`$ be a system of local coordinates on $`X`$ such that $`M=\{(x_1,\mathrm{},x_n);x_1=\mathrm{}=x_d=0\}`$ and let $`(x;\xi )`$ denote the associated coordinates on $`T^{}X`$. Set $`x^{}=(x_1,\mathrm{},x_d)`$, $`x^{\prime \prime }=(x_{d+1},\mathrm{},x_n)`$.
Let $`(x_0^{\prime \prime };\xi _0^{\prime \prime })T^{}M`$ such that $`(x_0^{\prime \prime };\xi _0^{\prime \prime })j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X)`$. We shall prove that $`(x_0^{\prime \prime };\xi _0^{\prime \prime })SS_k(F|_M)`$.
By the assumption, $`(0,x_0^{\prime \prime };\xi ^{},\xi _0^{\prime \prime })SS_k(F)\widehat{+}T_M^{}X`$ for any $`\xi ^{}^d`$. In particular, $`(0,x_0^{\prime \prime };0,\xi _0^{\prime \prime })SS_k(F)\widehat{+}T_M^{}X`$. We may assume that $`(0,x_0^{\prime \prime })\pi (SS_k(F))M`$ and by (3), that $`\xi _0^{\prime \prime }0`$.
Setting $`x_0=(0,x_0^{\prime \prime })`$, $`\xi _0=(0,\xi _0^{\prime \prime })`$ and $`p=(x_0,\xi _0)`$, there exists a closed convex proper cone $`\gamma `$ such that Int$`(\gamma )\mathrm{}`$ and
$$\{\begin{array}{cc}\xi _0\text{Int}(\gamma ^a),\hfill & \\ (\{x_0\}\times \gamma ^a)(SS_k(F)\widehat{+}T_M^{}X)T_X^{}X.\hfill & \end{array}$$
(7)
Therefore we may find a neighborhood $`V`$ of $`x_0`$ such that
$$(V\times \gamma ^a)(SS_k(F)\widehat{+}T_M^{}X)T_X^{}X.$$
(8)
In particular, $`(\{x_0\}\times \gamma ^a)SS_k(F)\{(x_0;0)\}`$. Therefore,
$$(\{x_0\}\times \text{Int}(\gamma ^a))SS_k(F)=\mathrm{},$$
and we may choose $`V`$ such that
$$(V\times \text{Int}(\gamma ^a))SS_k(F)=\mathrm{}.$$
(9)
and
$$(V\times \gamma ^a)T_M^{}XT_X^{}X.$$
(10)
After changing the local coordinates on $`X`$ if necessary, we may also assume:
$$\{\begin{array}{cc}\xi _0=(0,\mathrm{},0,1),\hfill & \\ \gamma ^a=\{(\xi ^{},\xi ^{\prime \prime });\xi _n\delta |(\xi ^{},\xi _{d+1},\mathrm{}\xi _{n1})|\},\hfill & \end{array}$$
for some $`\delta >0`$. Hence,
$$\gamma =\{(x^{},x^{\prime \prime });x_n\frac{1}{\delta }|(x^{},x_{d+1},\mathrm{}x_{n1})|\},$$
and, for any $`xM,`$ $`(x+\gamma )M=x+(\gamma M)`$. For $`\epsilon >0`$ let us denote by $`H_\epsilon `$ the open half-space $`H_\epsilon =\{xX;\text{Re}xx_0,\xi _0>\epsilon \}`$. Let us choose $`\epsilon >0`$ and an open neighborhood $`WV`$ of $`x_0`$ such that $`(W+\gamma )H_\epsilon V`$. Set $`\gamma ^{}=\gamma M`$, $`V^{}=VM`$, $`W^{}=WM`$ and $`H_\epsilon ^{}=\{x^{\prime \prime }M;x^{\prime \prime }x_0^{\prime \prime },\xi _0^{\prime \prime }>\epsilon \}`$. Since $`\gamma ^{}`$ is a closed convex proper cone in $`M`$ such that $`\xi _0^{}\text{Int}(\gamma ^a)`$ and $`W^{}`$ is an open neighborhood of $`x_0^{}`$ in $`M`$, by Proposition 4.1 its enough to prove that there exists $`\epsilon ^{}>0`$ such that $`(W^{}+\gamma ^{})\overline{H_\epsilon ^{}^{}}M`$ and $`H^j(M;๐ค_{(x+\gamma ^{})H_\epsilon ^{}^{}}F|_M)=0,`$ for all $`jk`$ and $`xW^{}`$.
This will be a consequence of Lemma 3.4 with $`\mathrm{\Lambda }=SS_k(F)`$. We shall use the notation $`\gamma _\lambda ,V_\lambda `$ introduced in Section 3. Let $`C`$ be given by Lemma 3.4 and let us choose sequences $`(\lambda _l)_l`$, $`(ฯต_l)_l`$ of real positive numbers, satisfying $`0<ฯต_l,\lambda _l<C`$, such that $`(\lambda _l)_l`$ converges to $`0`$ and $`(ฯต_l)_l`$ converges to $`C`$.
Remark that $`\gamma _{\lambda _l}^a\gamma ^a`$ and that
$$V_{\lambda _lฯต_l}M=VM=V^{},$$
$$W_{\lambda _lฯต_l}M=WM=W^{},$$
$$W_{\lambda _lฯต_l}+\gamma _{\lambda _lฯต_l}H_\epsilon V_{\lambda _lฯต_l},$$
$$(W_{\lambda _lฯต_l}+\gamma _{\lambda _l})H_{ฯต_l\epsilon }V_{\lambda _lฯต_l}.$$
Let $`x^{\prime \prime }MW`$ be given, choose a sequence $`x_l^{\prime \prime }`$ in $`W`$ converging to $`x^{\prime \prime }`$ such that $`x^{\prime \prime }\text{Int}(x_l^{\prime \prime }+\gamma ^{})`$ and note $`H^{}=HM`$. Then, for any $`jk`$, we have
$$H^j(M;๐ค_{(x^{\prime \prime }+\gamma ^{})H_{C\epsilon }^{}}F|_M)\underset{\begin{array}{c}\\ l\end{array}}{lim}H^j(X;๐ค_{(x_l^{\prime \prime }+\gamma _{\lambda _l})H_{ฯต_l\epsilon }}F)=0,$$
thanks to Lemma 4.7. Hence $`(x_0^{\prime \prime };\xi _0^{\prime \prime })SS_k(F|_M)`$. q.e.d.
###### End of the proof of Theorem 1.2.
Let us decompose $`f`$ by the graph map
$`Y\underset{๐}{\overset{}{}}Y\times X\underset{}{\overset{}{}}X`$, $`f=hg`$
where $`g(y)=(y,f(y))`$ and $`h`$ is the second projection on $`Y\times X`$.
Identifying $`Y`$ with the graph of $`f`$, we may assume that $`Y`$ is a closed subvariety of $`Y\times X`$, and we get by Proposition 6.1 and Proposition 4.6,
$$SS_k(f^1F)=SS_k(g^1(h^1F))$$
$$g_dg_\pi ^1(h_dh_\pi ^1(SS_k(F))\widehat{+}T_Y^{}(Y\times X)).$$
We shall prove that
$$g_dg_\pi ^1(h_dh_\pi ^1(SS_k(F))\widehat{+}T_Y^{}(Y\times X))=f^\mathrm{\#}(SS_k(F)).$$
Let $`(y)`$ be a system of local coordinates on $`Y`$, $`(x)`$ a system of local coordinates on on $`X`$ and let $`(y;\xi )`$, $`(x;\eta )`$ be the associated coordinates on $`T^{}Y`$ and $`T^{}X`$, respectively.
Let $`(y_0;\xi _0)g_dg_\pi ^1(h_dh_\pi ^1(SS_k(F))\widehat{+}T_Y^{}(Y\times X)),`$ then there exists $`\xi ,\eta `$ such that $`(y_0,f(y_0);\xi ,\eta )h_dh_\pi ^1(SS_k(F))\widehat{+}T_Y^{}(Y\times X)`$ and $`\xi _0=\xi +{}_{}{}^{t}f_{}^{}(y_0)\eta `$. Hence we may find sequences $`\{(y_n,x_n;\xi _n,\eta _n)\}_n`$ in $`h_dh_\pi ^1(SS_k(F))`$ and $`\{(y_n^{},f(y_n^{});\xi _n^{},\eta _n^{})\}_n`$ in $`T_Y^{}(Y\times X)`$ such that
$$\{\begin{array}{cc}(y_n,x_n),(y_n^{},f(y_n^{}))\underset{๐}{\overset{}{}}(y_0,f(y_0)),\hfill & \\ (\xi _n,\eta _n)+(\xi _n^{},\eta _n^{})\underset{๐}{\overset{}{}}(\xi ,\eta ),\hfill & \\ |(y_n,x_n)(y_n^{},f(y_n^{}))||(\xi _n,\eta _n)|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
One has $`(x_n;\eta _n)SS_k(F),`$ $`\xi _n=0`$ and $`\xi _n^{}+{}_{}{}^{t}f_{}^{}(y_n^{})\eta _n^{}=0`$; hence we obtain $`{}_{}{}^{t}f_{}^{}(y_n^{})(\eta _n+\eta _n^{})\underset{๐}{\overset{}{}}{}_{}{}^{t}f_{}^{}(y_0)\eta =\xi _0\xi `$ and then $`{}_{}{}^{t}f_{}^{}(y_n^{})\eta _n\underset{๐}{\overset{}{}}\xi _0.`$
Therefore we have sequences $`\{(x_n;\eta _n)\}_nSS_k(F)`$ and $`\{y_n^{}\}_n`$ in $`Y`$ such that
$$\{\begin{array}{cc}y_n\underset{๐}{\overset{}{}}y_0,x_n\underset{๐}{\overset{}{}}f(y_0),\hfill & \\ {}_{}{}^{t}f_{}^{}(y_n^{})\eta _n\underset{๐}{\overset{}{}}\xi _0\hfill & \\ |x_nf(y_n^{})||\eta _n|\underset{๐}{\overset{}{}}0.\hfill & \end{array}$$
This gives $`(y_0;\xi _0)f^\mathrm{\#}(SS_k(F))`$ and also the converse thanks to Proposition 3.1. q.e.d.
###### Corollary 6.2.
Let $`M`$ be a closed submanifold of $`X`$ and $`FD^b(๐ค_X)`$. Then
$$SS_k(F_M)SS_k(F)\widehat{+}T_M^{}X.$$
###### Proof.
Let $`j:MX`$ denote the embedding of $`M`$ on $`X`$. Then $`F_Mj_{}(F|_M)`$ and by Proposition 4.5 and Proposition 6.1
$$SS_k(F_M)=j_\pi j_d^1(SS_k(F|_M))$$
$$j_\pi j_d^1j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X)SS_k(F)\widehat{+}T_M^{}X.$$
q.e.d.
###### Proposition 6.3.
Let $`M`$ be a closed submanifold of $`X`$, $`U=XM`$, $`j`$ the embedding $`UX`$, $`\iota `$ the embedding of $`M`$ in $`X`$ and let $`FD^b(๐ค_U)`$. Then:
(i) $`SS_k(Rj_{}F)\pi ^1(M)SS_k(F)\widehat{+}T_M^{}X,`$
(ii) $`SS_k(Rj_!F)\pi ^1(M)SS_k(F)\widehat{+}T_M^{}X,`$
(iii) $`SS_k((Rj_{}F)|_M)\iota _d\iota _\pi ^1(SS_k(F)\widehat{+}T_M^{}X).`$
###### Proof.
The proof of the two first conditions is analogous to the proof of the two first conditions of Proposition 6.3.2 of , replacing Proposition 5.4.4 and Theorem 6.3.1 of , by Propositions 4.5 and 5.3, respectively.
Let us now prove the third inequality. By Proposition 6.1 and $`(i)`$,
$$SS_k((Rj_{}F)|_M)\iota _d\iota _\pi ^1(SS_k(Rj_{}F)\widehat{+}T_M^{}X)$$
$$\iota _d\iota _\pi ^1((SS_k(F)\widehat{+}T_M^{}X)\widehat{+}T_M^{}X).$$
By Lemma 3.3
$$(SS_k(F)\widehat{+}T_M^{}X)\widehat{+}T_M^{}X=SS_k(F)\widehat{+}T_M^{}X.$$
Hence
$$SS_k((Rj_{}F)|_M)\iota _d\iota _\pi ^1(SS_k(F)\widehat{+}T_M^{}X).$$
q.e.d.
Note that, with Lemma 3.2 and Proposition 6.3 in hand, we obtain the analogue of Proposition 6.3.2 of .
###### Corollary 6.4.
Let $`M`$ be a closed submanifold of $`X`$ and $`FD^b(๐ค_X)`$. Then
$$SS_k(R\mathrm{\Gamma }_M(F))SS_k(F)\widehat{+}T_M^{}X.$$
###### Proof.
This is a consequence of Proposition 6.3, together with the distinguished triangle
$$R\mathrm{\Gamma }_M(F)FR\mathrm{\Gamma }_{X\backslash M}(F)\stackrel{+1}{}.$$
q.e.d.
###### Corollary 6.5.
Let $`M`$ be a closed submanifold of $`X`$ and $`FD^b(๐ค_X)`$. Assume that
$$SS_k(F)T_M^{}XT_X^{}X.$$
Then we have a natural isomorphism
$$\tau ^k(F_M\omega _{M|X})\tau ^k(R\mathrm{\Gamma }_M(F)).$$
In particular
$$SS_k(F_M\omega _{M|X})=SS_k(R\mathrm{\Gamma }_M(F)).$$
###### Proof.
Let $`\dot{\pi }`$ be the restriction of $`\pi `$ to the cotangent bundle deprived of the zero section. We have a distinguished triangle
$$\tau ^k(F_M\omega _{M|X})\tau ^k(R\mathrm{\Gamma }_M(F))\tau ^k(\dot{\pi }(\mu _M(F)|_{\dot{T}_M^{}X}))\underset{+1}{}.$$
Since by Theorem 5.1 of $`supp(\tau ^k(\mu _M(F)))SS_k(F)T_M^{}X`$, the third term vanishes hence the result. q.e.d.
###### Corollary 6.6.
Let $`M`$ be a closed submanifold of $`X`$ and $`FD^b(๐ค_X)`$. Let $`j`$ denote the embedding of $`M`$ in $`X`$.Then
$$SS_k(j^!F)j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X).$$
###### Proof.
This is a consequence of Corollary 6.4 and Proposition 6.1 together with Lemma 3.3. q.e.d.
###### Proof of Theorem 1.1.
The proof is the adaptation step by step of the proof of Theorem 6.4.1 of , applying Proposition 4.6, and Proposition 6.3 instead of Proposition 5.4.5 and Proposition 6.3.2, respectively of . q.e.d.
Let now $`Y`$ be a complex closed smooth hypersurface of $`X`$ defined as the zero locus of a holomorphic function $`f`$. Let $`\psi _Y`$ denote the functor of nearby cycles associated to $`Y`$. Then $`Y`$ may be regarded as a submanifold of $`T_YX`$ by a canonical section $`s`$ such that $`\psi _Y(F)s^1\nu _Y(F)`$. Once more we identify $`T_{T_Y^{}X}T^{}X`$, $`T^{}(T_Y^{}X)`$ and $`T^{}(T_YX)`$ (cf. Proposition 5.5.1 of ).
Recall that, in a system of linear coordinates $`x=(x_1,\mathrm{},x_n)`$ on $`X`$ such that $`Y`$ is defined by $`x_1=0`$, $`s:YT_YX`$ is the section $`s(x_2,\mathrm{},x_n)=(x_2,\mathrm{},x_n;1)`$. With the local coordinates described above, and $`A`$ being a conic closed subset of $`T^{}(T_Y^{}X)`$, we have:
$$s_ds_\pi ^1(A)=\{(x_2,\mathrm{},x_n;\xi _2,\mathrm{},\xi _n);\xi _1,(x_2,\mathrm{},x_n,1;\xi _2,\mathrm{},\xi _n,\xi _1)A\}.$$
Corollary 1.4 is an immediate consequence of Proposition 6.1.
The following estimate for the tensor product can be seen as a generalization of Proposition 6.1:
###### Proposition 6.7.
Let $`F`$ and $`G`$ belong to $`D^b(๐ค_X)`$. Then:
$$SS_k(F^๐G)\underset{i+j=k}{}(SS_i(F)\widehat{+}SS_j(G)).$$
###### Proof.
Let $`\delta _X:XX\times X`$ be the diagonal embedding.
Since $`F^๐G\delta _X^1(F^๐G)`$, the result follows from Proposition 6.1 and Proposition 4.4. q.e.d.
### 6.2 Application to $`๐`$-modules
Let $`X`$ be a complex finite dimensional manifold. One of the important problems in the theory of $`๐`$-modules is the relation between the caracteristic variety of a system $``$ and that of its induced system $`_Y`$ along a closed submanifold $`Y`$, which was completely solved in the non charateristic case by M. Kashiwara as well as in a more general situation treated in , which includes the case where $``$ is regular along $`Y`$ in the sense of . Similarly, in the case of a smooth complex hypersurface, it is interesting to relate $`\text{Char}()`$ and $`\text{Char}(\psi _Y())`$, where $`\psi _Y`$ denotes the functor of nearby cycles.
Let $`d`$ be the codimension of $`Y`$, denote by $`j`$ the embedding $`YX`$ and by $`\pi ^{}`$ the projection $`T^{}YY`$. Given an homogeneous involutive subvariety $`V`$ of $`T^{}X`$ of codimension $`d`$, we shall say that $`Y`$ is orthogonal to $`V`$ if there exists a smooth involutive submanifold $`V^{}`$ containing $`V`$ such that $`Y`$ and $`V^{}`$ are orthogonal. More precisely, there exist a set $`\{f_1,\mathrm{},f_d\}`$ of homogeneous functions of degree zero vanishing on $`\pi ^1(Y)`$, such that the differential $`df_i`$ are linearly independent on $`\pi ^1(Y)`$, and a set $`\{g_1,\mathrm{},g_p\},pd,`$ of homogeneous functions of degree one linearly independent on $`V^{}`$ such that the matrix of the Poisson brackets $`[\{f_i,g_j\}]|_V^{}`$ has everywhere rank $`d`$.
As before, $`F`$ will denote the complex $`Rom_{๐_X}(,๐ช_X)`$. Let $`SS(F)=_\alpha V_\alpha `$ be the decomposition of $`SS(F)`$ in its irreducible involutive components in a neighborhood of $`pT^{}X`$. Let us denote by $`Y_\alpha `$ the variety $`\pi (V_\alpha )`$.
Recall that in Theorem 6.7 of it is proved that, for any $`k`$, $`SS_k(F)=(\underset{codimY_\alpha <k}{}V_\alpha )(\underset{codimY_\alpha =k}{}T_{Y_\alpha }^{}X)`$.
###### Proof of Theorem 1.4.
The first assertion is an immediate consequence of Theorem 1.1 and the second follows from the regularity of $``$. q.e.d.
###### Proof of Corollary 1.5.
It is a consequence of Corollary 1.3 and the regularity of $``$. q.e.d.
###### Proof of Theorem 1.6.
Since $``$ is regular along $`Y`$, one has the isomorphism
$$Rom_{๐_Y}(_Y,๐ช_Y)Rom_{๐_X}(,๐ช_X)|_Y,$$
and the first part is an immediate consequence of Proposition 6.1. Let us now prove the second assertion. It will be a consequence of the Lemma below:
###### Lemma 6.8.
Assume that the homogeneous involutive variety $`V`$ is irreducible and that $`Y`$ is orthogonal to $`V`$.
Then :
(i) $`V^{}=j_d(j_\pi ^1(V))`$ is an irreducible homogeneous involutive subvariety of $`T^{}Y`$.
(ii) The codimension of $`\pi ^{}(V^{})`$ is equal to the codimension of $`\pi (V)`$.
(iii) When $`V`$ is the characteristic variety of a coherent $`๐_X`$-module, the orthogonality of $`Y`$ implies that $`Y`$ is non characteristic for $``$.
###### Proof.
Let $`V^{}`$ be a smooth involutive manifold containing $`V`$ such that $`Y`$ is orthogonal to $`V^{}`$. Since the assertions can be checked locally, by a standard reasoning we may consider a system $`(x;\xi )`$ of local symplectic coordinates on $`T^{}X`$ in a neighborhood of $`pV\pi ^1(Y)=j_\pi ^1(V)`$, such that $`Y`$ is the sumanifold $`\{(x)=(x_1,\mathrm{}x_n);x_1=\mathrm{}=x_d=0\}`$ and $`V^{}`$ is defined in $`T^{}X`$ by the equations $`\xi _1=\mathrm{}=\xi _d=g_{d+1}(x^{\prime \prime };\xi ^{\prime \prime })=\mathrm{}=g_p(x^{\prime \prime };\xi ^{\prime \prime })=0`$, where we set $`(x^{})=(x_,\mathrm{},x_d)`$ (resp. $`(\xi ^{})=(\xi _1,\mathrm{}\xi _d)`$), $`(x^{\prime \prime })=(x_{d+1},\mathrm{},x_n)`$ (resp. $`(\xi ^{\prime \prime })=(\xi _{d+1},\mathrm{},\xi _n)`$). Therefore, the irreducible ideal of definition $`I(V)`$ is generated by a set of functions
$$\{\xi _1,\mathrm{},\xi _d,g_{d+1}(x^{\prime \prime };\xi ^{\prime \prime }),\mathrm{},g_p(x^{\prime \prime };\xi ^{\prime \prime }),h_{p+1}(x\mathrm{"};\xi \mathrm{"}),\mathrm{},h_{p+l}(x\mathrm{"};\xi \mathrm{"})\},$$
for some $`l0`$. Hence $`I(V^{})`$ is generated in $`๐ช_{T^{}Y}`$ by the set of functions
$$\{g_{d+1}(x^{\prime \prime };\xi ^{\prime \prime }),\mathrm{},g_p(x^{\prime \prime };\xi ^{\prime \prime }),h_{p+1}(x\mathrm{"};\xi \mathrm{"}),\mathrm{},h_{p+l}(x\mathrm{"};\xi \mathrm{"})\},$$
which entails $`(i)`$, $`(ii)`$ and $`(iii)`$. q.e.d.
Since $`Y`$ is non characteristic, we have
$$SS(F|_Y)=\text{Char}(_Y)=j_dj_\pi ^1(SS(F)).$$
On the other side, since $`SS_k(F)T_Y^{}XT_X^{}X`$, we get from the first assertion that $`SS_k(F|_Y)j_dj_\pi ^1(SS_k(F))`$, for any k. Moreover, setting $`V_\alpha ^{}=j_dj_\pi ^1(V_\alpha )`$, by the preceding Lemma, for any $`\alpha `$ such that $`codimY_\alpha k`$, $`V_\alpha ^{}`$ is an irreducible component of $`SS(F|_Y)`$. Therefore by Theorem 6.7 of , for any $`ik,`$
$$SS_i(F|_Y)j_dj_\pi ^1(SS_i(F)).$$
q.e.d.
###### Example 6.9.
Let $`X=^n`$, with $`n3`$, endowed with the coordinates $`(x_1,\mathrm{},x_n)`$. Let $`Y`$ be the hypersurface $`\{x_n=0\}`$ and $`\mathrm{\Omega }=\{xX;\text{Re}(x_1x_{n1})<0\}`$. Let $`๐ฅ`$ be a coherent left ideal of $`๐_X`$ and set $`=๐_X/๐ฅ`$. Assume that there exist in $`๐ฅ`$ an operator $`P`$ in the Weierstrass form with respect to the derivation $`D_{x_n}`$ and an operator $`Q`$ such that the principal symbol of $`Q`$, $`\sigma (Q)`$, is of the form
$$\sigma (Q)=x_1q(x_1,\mathrm{},x_{n1};\xi _1,\mathrm{},\xi _{n1}),$$
and $`q`$ does not vanish on $`T_{\delta \mathrm{\Omega }}^{}X`$. Then,$`T_{\delta (\omega )}^{}XSS_1()\{0\}`$ and, setting $`\mathrm{\Omega }^{}=\mathrm{\Omega }Y`$, $`\mathrm{\Omega }^{}`$ has smooth boundary $`\delta \mathrm{\Omega }^{}`$. By Theorem 1.3, $`T_{\delta \mathrm{\Omega }^{}}^{}YSS_1(_Y)\{0\}.`$ Therefore
$$\text{om}_{๐_Y}(_Y,_{\{\text{Re}(x_1x_{n1})0\}}^1(๐ช_Y))|_{\delta \mathrm{\Omega }^{}}=0.$$
###### Proof of Corollary 1.5.
As proved in , we have the isomorphism
$$Rom_{๐_Y}(\psi _Y(),๐ช_Y)\psi _Y(Rom_{๐_X}(,๐ช_X)).$$
It is then enough to use Proposition 6.1. q.e.d.
Let $`M`$ be a real analytic manifold of dimension $`n`$, $`X`$ a complex analytic manifold complexifying $`M`$ and $``$ a coherent $`๐_X`$-module.
Let $`๐_M`$ denote the sheaf of real analytic functions on $`M`$. Remark that $`๐_M=๐ช_X|_M`$. Let $`_M`$ denote the sheaf of Satoโs hyperfunctions on $`M`$. Recall that
$$_MR\mathrm{\Gamma }_M(๐ช_X)or_{M/X}.$$
###### Proof of Proposition 1.7.
One has
$$Rom_{๐_X}(,๐_M)F|_M.$$
Therefore, by Proposition 6.1
$$SS_k(Rom_{๐_X}(,๐_M))j_dj_\pi ^1(SS_k(F)\widehat{+}T_M^{}X)$$
q.e.d.
Let us remark that a variant of the preceding result was obtained in for $`k=1`$ using directly the properties of holomorphic functions.
###### Proof of Corollary 1.8.
The first part is an immediate consequence of Corollary 6.5. The second follows from Proposition 1.7 and Theorem 6.7 of . q.e.d.
###### Example 6.10.
Let $`M=^n`$, with $`n2`$, endowed with the coordinates $`x=(x_1,\mathrm{},x_n)`$. Let $`\mathrm{\Omega }=\{xM;\varphi (x)<0\}`$ for some real $`C^1`$-function. Let $`X=^n`$ and $``$ be a coherent $`๐_X`$-module defined by $`=๐_X/๐_XP`$ where $`P`$ is a differential operator. Assume that the principal symbol $`\sigma (P)`$ is of the form
$$\sigma (P)=a(x)q(x;\xi ),$$
where $`a(x)`$ is an holomorphic function and $`q`$ does not vanish on $`T^{}M`$, more precisely, $`q`$ is the principal symbol of an elliptic operator. Recall that $`SS_1(F)\overline{\{(x;\xi );a(x)=0,\xi da(x)\}}q^1(0).`$This entails that $`T_M^{}XSS_1(F)T_X^{}X`$ hence
$$SS_1(Rom_{๐_X}(,๐_M))=SS_1(Rom_{๐_X}(,_M))j_dj_\pi ^1(SS_1(F)).$$
Assume that $`d\varphi (x)`$ is not in $`\overline{da(x)}`$ for any $`x\delta \mathrm{\Omega }a^1(0)`$. Hence $`T_{\delta \mathrm{\Omega }}^{}MSS_1(Rom_{๐_X}(,๐_M))T_M^{}M.`$ In other words
$$\text{om}_{๐_X}(,_{\{\varphi (x)0\}}^1(๐_M))|_{\delta \mathrm{\Omega }}$$
$$=\text{xt}_{๐_X}^1(,\mathrm{\Gamma }_{\{\varphi (x)0\}}(_M))|_{\delta \mathrm{\Omega }}=0.$$
###### Remark 6.11.
In general we do not have an interesting estimate for
$`SS_k(Rom_{๐_X}(,_M))`$. Let $`j`$ denote the inclusion $`MX`$. Then $`_Mj^!๐ช_Xor_{M/X}`$ and $`Rom_{๐_X}(,_M)j^!(Rom_{๐_X}(,๐ช_X))[n]`$. By Corollary 6.6, one gets
$$SS_k(Rom_{๐_X}(,_M))=SS_{k+n}(j^!Rom_{๐_X}(,๐ช_X))$$
$$SS_{k+n}(Rom_{๐_X}(,๐ช_X))\widehat{+}T_M^{}X.$$
By Theorem 6.7 of ,
$$SS_{k+n}(Rom_{๐_X}(,๐ช_X))=SS(Rom_{๐_X}(,๐ช_X))=\text{Char}().$$
Hence we get
$$SS_k(Rom_{๐_X}(,_M))\text{Char}()\widehat{+}T_M^{}X\text{for any k}0,$$
in other words, if $`M`$ is hyperbolic for $``$ then
$$SS_{k+n}(Rom_{๐_X}(,๐ช_X))T_X^{}X.$$
But this is well known and is an example that the notion of truncated microsupport does not work well under Fourier Transform.
Let $`D_c^b(๐ค_X)`$ denote the full subcategory of $`D^b(๐ค_X)`$ consisting of objects with $``$-constructible cohomology, that is, the objects $`FD^b(๐ค_X)`$ for which there exists a complex analytic stratification $`X=X_\alpha `$ such that the sheaf $`H^j(F)|_{X_\alpha }`$ is locally constant of finite rank, for every $`j`$ and $`\alpha `$.
A perverse sheaf is an object $`F`$ of $`D_c^b(๐ค_X)`$ satisfying the following two conditions:
*(a) for any complex submanifold $`Y`$ of $`X`$ of codimension $`d`$, $`H_Y^j(F)|_Y`$ is zero for $`j<d`$;*
*(b) for any $`j`$, $`H^j(F)`$ is supported by a complex analytic subset of codimension $`j`$.*
P. Schapira proved in that, when $`F`$ is a perverse object of $`D_c^b(๐ค_X)`$,
$`H^j(R\mathrm{\Gamma }_S(F))_x=0`$, for $`j2n`$,
for any closed subanalytic subset $`S`$ of $`X`$ and any $`xX`$ being non isolated in $`S`$.
###### Proposition 6.12.
Let $``$ be a coherent $`๐_X`$-module Then
$$SS_{n1}(Rom_{๐_X}(,_M))=SS(Rom_{๐_X}(,_M)).$$
###### Proof.
Let $`\phi `$ be a real analytic function defined on $`X`$ and $`x_0X`$ such that $`\phi (x_0)=0`$. Then the set $`\{xX;\phi (x)0\}`$ is a closed subanalytic subset of $`X`$. Assume that $``$ is holonomic. By the Riemann-Hilbert correspondence ( ), $`F`$ is perverse.
Hence,
$$H^j(R\mathrm{\Gamma }_{\{\phi 0\}}Rom_{๐_X}(,_M))_{x_0}H^{j+n}(R\mathrm{\Gamma }_{\{\phi 0\}M}(F))_{x_0}=0,$$
for every $`jn`$. By Proposition 4.1,
$$SS_{n1}(Rom_{๐_X}(,_M))=SS(Rom_{๐_X}(,_M)),$$
under the assumption that $``$ is an holonomic $`๐_X`$-module.
To treat the general case, we argue as in the proof of Theorem 2 of . Let us denote by the functor $`๐ฉ๐ฉ^{}=xt_{๐_X}^n(๐ฉ,๐_X).`$ Recall that Kashiwara proved in that if $``$ is coherent, then $`^{}`$ is holonomic, $`^{}^{}`$ and $`^{}`$ is a submodule of $``$. Defining the coherent $`๐_X`$-module $``$ by the exact sequence:
$$0^{}0,$$
one gets $`^{}=0`$ and so $``$ locally admits a projective resolution of lenght $`n1`$. Therefore,
$$H^j(R\mathrm{\Gamma }_{\{\phi 0\}}Rom_{๐_X}(,_M))_{x_0}H^j(R\mathrm{\Gamma }_{\{\phi 0\}}Rom_{๐_X}(^{},_M))_{x_0}=0.$$
for $`jn`$.
This proves
$$SS_{n1}(Rom_{๐_X}(,_M))=SS(Rom_{๐_X}(,_M)),$$
for every coherent $`๐_X`$-module $``$. q.e.d.
Ana Rita Martins
Centro de รlgebra da Universidade de Lisboa, Complexo 2,
2 Avenida Prof. Gama Pinto, 1699 Lisboa Portugal
arita@pcmat.fc.ul.pt
Teresa Monteiro Fernandes
Centro de รlgebra da Universidade de Lisboa, Complexo 2,
2 Avenida Prof. Gama Pinto, 1699 Lisboa Portugal
tmf@ptmat.fc.ul.pt |
warning/0506/math0506277.html | ar5iv | text | # ARITHMETIC PROPERTIES OF PROJECTIVE VARIETIES OF ALMOST MINIMAL DEGREE
## 1. Introduction
Let $`_k^r`$ denote the projective $`r`$-space over an algebraically closed field $`k.`$ Let $`X_k^r`$ be an irreducible non-degenerate projective variety of dimension $`d`$. The degree $`\mathrm{deg}X`$ of $`X`$ is defined as the number of points of $`X๐,`$ where $`๐`$ is a linear subspace defined by generically chosen linear forms $`\mathrm{}_1,\mathrm{},\mathrm{}_d`$. It is well known that
$$\mathrm{deg}X\mathrm{codim}X+1,$$
(cf e.g. ), where $`\mathrm{codim}X=rd`$ is used to denote the codimension of $`X`$. In case equality holds, $`X`$ is called a variety of minimal degree. Varieties of minimal degree are classified and well understood. A variety $`X`$ of minimal degree is either a quadric hypersurface, a (cone over a) Veronese surface in $`_k^5,`$ or a (cone over a smooth) rational normal scroll (cf \[21, Theorem 19.9\]). In particular these varieties are arithmetically Cohen-Macaulay and arithmetically normal.
The main subject of the present paper is to investigate varieties of almost minimal degree, that is irreducible, non-degenerate projective varieties $`X_k^r`$ with $`\mathrm{deg}X=\mathrm{codim}X+2`$. From the point of view of polarized varieties, Fujita , , has studied extensively such varieties in the framework of varieties of $`\mathrm{\Delta }`$-genus $`1`$. Nevertheless, in our investigation we take a purely arithmetic point of view and study our varieties together with a fixed embedding in a projective space.
A natural approach to understand a variety $`X_k^r`$ of almost minimal degree is to view it (if possible) as a birational projection of a variety of minimal degree $`\stackrel{~}{X}_k^{r+1}`$ from a point $`p_k^{r+1}\stackrel{~}{X}`$. If sufficiently many varieties of almost minimal degree can be obtained by such projections, we may apply to them the program of: โclassifying by projections of classified varietiesโ. It turns out, that this classification scheme can indeed be applied to an interesting class of varieties of almost minimal degree $`X_k^r`$, namely those, which are not arithmetically normal or equivalently, to all those which are not simultaneously normal and arithmetically Gorenstein. More precisely, we shall prove the following result, in which $`\mathrm{Sec}_p(\stackrel{~}{X})`$ is used to denote secant cone of $`\stackrel{~}{X}`$ with respect to $`p`$:
###### Theorem 1.1.
Let $`X_k^r`$ be a non-degenerate irreducible projective variety and let $`t\{1,2,\mathrm{},dimX+1\}`$. Then, the following conditions are equivalent:
* $`X`$ is of almost minimal degree, of arithmetic depth $`t`$ and not arithmetically normal.
* $`X`$ is of almost minimal degree and of arithmetic depth $`t`$, where either $`tdimX`$ or else $`t=dimX+1`$ and $`X`$ is not normal.
* $`X`$ is of almost minimal degree and of arithmetic depth $`t`$, where either $`X`$ is not normal and $`t>1`$ or else $`X`$ is normal and $`t=1`$.
* $`X`$ is a (birational) projection of a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree from a point $`p_k^{r+1}\stackrel{~}{X}`$ such that $`dim\mathrm{Sec}_p(\stackrel{~}{X})=t1`$.
For the proof of this result see Theorem 5.6 (if $`tdimX`$) resp. Theorem 6.9 (if $`t=dimX+1`$). In the spirit of Fujita we say that a variety of almost minimal degree is maximally Del Pezzo if it is arithmetically Cohen-Macaulay (or โ equivalently โ arithmetically Gorenstein). Then, as a consequence of Theorem 1.1 we have:
###### Theorem 1.2.
A variety $`X_k^r`$ of almost minimal degree is either maximally Del Pezzo and normal or a (birational) projection of a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree from a point $`p_k^{r+1}\stackrel{~}{X}`$.
In this paper, our interest is focussed on those varieties $`X_k^r`$ of almost minimal degree which are birational projections of varieties of minimal degree. As already indicated by Theorem 1.1 and in accordance with our arithmetic point of view, the arithmetic depth of $`X`$ is the key invariant of our investigation. It turns out, that this arithmetic invariant is in fact closely related to the geometric nature of our varieties. Namely, the picture sketched in Theorem 1.1 can be completed as follows:
###### Theorem 1.3.
Let $`X_k^r`$ be a variety of almost minimal degree and of arithmetic depth $`t`$, such that $`X=\varrho (\stackrel{~}{X})`$, where $`\stackrel{~}{X}_k^{r+1}`$ is a variety of minimal degree and $`\varrho :_k^{r+1}\{p\}_k^r`$ is a birational projection from a point $`p_k^{r+1}\stackrel{~}{X}.`$ Then:
* $`\nu :=\varrho :\stackrel{~}{X}X`$ is the normalization of $`X`$.
* The secant cone $`\mathrm{Sec}_p(\stackrel{~}{X})_k^{r+1}`$ is a projective subspace $`_k^{t1}_k^{r+1}`$.
* The singular locus $`\mathrm{Sing}(\nu )=\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})X`$ of $`\nu `$ is a projective subspace $`_k^{t2}_k^r`$ and coincides with the non-normal locus of $`X`$.
* If $`tdimX,\mathrm{Sing}(\nu )`$ coincides with the non $`S_2`$-locus and the non-Cohen-Macaulay locus of $`X`$ and the generic point of $`\mathrm{Sing}(\nu )`$ in $`X`$ is of Goto-type.
* The singular fibre $`\nu ^1(\mathrm{Sing}(\nu ))=\mathrm{Sec}_p(\stackrel{~}{X})\stackrel{~}{X}`$ is a quadric in $`_k^{t1}=\mathrm{Sec}_p(\stackrel{~}{X})`$.
For the proves of these statements see Theorem 5.6 and Corollary 6.10.
Clearly, the projecting variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree plays a crucial rรดle for $`X`$. We thus may distinguish the exceptional case in which $`X`$ is a cone over the Veronese surface and the general case in which $`\stackrel{~}{X}`$ is a cone over a rational normal scroll. In this latter case, we have the following crucial result, in which we use the convention $`dim\mathrm{}=1:`$
###### Theorem 1.4.
Let $`X_k^r`$ be a variety of almost minimal degree which is a birational projection of a (cone over a) rational normal scroll $`\stackrel{~}{X}_k^{r+1}`$ from a point $`p_k^{r+1}\stackrel{~}{X}`$. Then, there is a (cone over a) rational normal scroll $`Y_k^r`$ such that $`XY`$ and $`\mathrm{codim}_X(Y)=1`$. Moreover, if the vertex of $`\stackrel{~}{X}`$ has dimension $`h`$, the dimension $`l`$ of the vertex of $`Y`$ satisfies $`hlh+3.`$ In addition, the arithmetic depth $`t`$ satisfies $`th+5.`$
For a proof of this result see Theorem 7.3 and Corollary 7.5. It should be noticed, that there are varieties of almost minimal degree, which cannot occur as a $`1`$-codimensional subvariety of a variety of minimal degree (cf Example 9.4 and Remark 6.5).
Our paper is built up following the idea, that the arithmetic depth $`t:=0ptA`$ of a variety $`X_k^r=\mathrm{Proj}(S),S=k[x_0,\mathrm{},x_r],`$ of almost minimal degree with homogeneous coordinate ring $`A=A_X`$ is a key invariant. In Section 2 we present a few preliminaries and discuss the special case where $`X`$ is a curve.
In Section 3 we consider the case where $`t=1`$. We show that the total ring of global sections $`_nH^0(X,๐ช_X(n))`$ of $`X`$ โ that is the $`S_+`$-transform $`D(A)`$ of $`A`$ โ is the homogeneous coordinate ring of a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree. In geometric terms: $`X`$ is isomorphic to $`\stackrel{~}{X}`$ by means of a projection from a generic point $`p_k^{r+1}`$ and hence normal but not arithmetically normal (cf Propositions 3.1 and 3.4).
In Section 4 we begin to investigate the case $`(1<)tdimX`$, that is the case in which $`X`$ is not arithmetically Cohen-Macaulay. First, we prove some vanishing statements for the cohomology of $`X`$ and describe the structure of the $`t`$-th deviation module $`K^t(A)`$ of $`A`$. Moreover we determine the Hilbert series of $`A`$ and the number of defining quadrics of $`X`$ (cf Theorem 4.2 and Corollary 4.4).
In Section 5 we aim to describe $`X`$ as a projection if $`tdimX`$. As a substitute for the $`S_+`$-transform $`D(A)`$ of the homogeneous coordinate ring $`A`$ (which turned out to be useful in the case $`t=1`$) we now consider the endomorphism ring $`B:=\text{End}_A(K(A))`$ of the canonical module of $`A`$ (cf Theorem 5.3). It turns out that $`B`$ is the homogeneous coordinate ring of variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree, and this allows to describe $`X`$ as a projection of $`\stackrel{~}{X}`$ (cf Theorem 5.6). Endomorphism rings of canonical modules have been studied extensively in a purely algebraic setting (cf , ). The striking point is the concrete geometric meaning of these rings in the case of varieties of almost minimal degree.
In Section 6 we study the case where $`t=dimX+1`$, that is the case where $`X`$ is arithmetically Cohen-Macaulay. Now, $`X`$ is a Del Pezzo variety in the sense of Fujita . According to our arithmetic point of view we shall speak of maximal Del Pezzo varieties in order to distinguish them within the larger class of polarized Del Pezzo varieties. We shall give several equivalent characterizations of these varieties (cf Theorem 6.2). We shall in addition introduce the notion of Del Pezzo variety and show among other things that this notion coincides with Fujitaโs definition for the polarized pair $`(X,๐ช_X(1))`$ (cf Theorem 6.8). Finally we shall prove that the graded integral closure $`B`$ of the homogeneous coordinate ring $`A`$ of a non-normal maximal Del Pezzo variety $`X_k^r`$ is the homogeneous coordinate ring of a variety of minimal degree $`\stackrel{~}{X}_k^{r+1}`$ (cf Theorem 6.9) and describe $`X`$ as a projection of $`\stackrel{~}{X}`$ (cf Corollary 6.10). Contrary to the case in which $`tdimX`$, we now cannot characterize $`B`$ as the endomorphism ring of the canonical module $`K(A)`$, simply as $`A`$ is a Gorenstein ring. We therefore study $`B`$ by geometric arguments, which rely essentially on the fact that we know already that the non-normal locus of $`X`$ is a linear subspace (cf Proposition 5.8). It should be noticed that on turn these geometric arguments seem to fail if $`tdimX`$.
In Section 7 we assume that $`X`$ is a (birational) projection of a (cone over a) rational normal scroll $`\stackrel{~}{X}_k^{r+1}`$. We then prove what is claimed by the previous Theorem 1.4. Here, we extensively use the determinantal description of rational normal scrolls (cf ). As an application we give some constraints on the arithmetic depth $`t`$ of $`X`$ (cf Corollary 7.5 and Corollary 7.6).
In Section 8 we study the Betti numbers of the homogeneous coordinate ring $`A`$ of our variety of almost minimal degree $`X_k^r.`$ We focus on those cases, which after all merit a particular interest, that is the situation where $`tdimX`$ and $`X`$ is a projection of a rational normal scroll. Using what has been shown in Section 7, we get a fairly good and detailed view on the behaviour of the requested Betti numbers.
Finally, in Section 9 we present various examples that illustrate the results proven in the previous sections. In several cases we calculated the Betti numbers of the vanishing ideal of the occuring varieties on use of the computer algebra system Singular .
## 2. Preliminaries
We first fix a few notation, which we use throughout this paper. By $`_0`$ (resp. $`)`$ we denote the set of non-negative (resp. positive) integers.
###### Notation 2.1.
A) Let $`k`$ be an algebraically closed field, let $`S:=k[x_0,\mathrm{},x_r]`$ be a polynomial ring, where $`r2`$ is an integer. Let $`X_k^r=\mathrm{Proj}(S)`$ be a reduced irreducible projective variety of positive dimensions $`d`$. Moreover, let $`๐ฅ=๐ฅ_X๐ช_{_k^r}`$ denote the sheaf of vanishing ideals of $`X`$, let $`I=I_X=_nH^0(_k^r,๐ฅ(n))S`$ denote the vanishing ideal of $`X`$ and let $`A=A_X:=S/I`$ denote the homogeneous coordinate ring of $`X`$.
B) If $`M`$ is a finitely generated graded $`S`$-module and if $`i`$, we use $`H^i(M)=H_{S_+}^i(M)`$ to denote the $`i`$-th local cohomology module of $`M`$ with respect to the irrelevant ideal $`S_+=_nS_n`$ of $`S`$. Let $`D(M)=D_{S_+}(M)`$ denote the $`S_+`$-transform $`\underset{}{\mathrm{lim}}\mathrm{Hom}_S(S_+^n,M)`$ of $`M`$. Moreover, let us introduce the $`i`$-th deficiency module of $`M`$:
$$K^i(M)=K_S^i(M):=\mathrm{Ext}_S^{r+1i}(M,S(r1)).$$
The $`S`$-modules $`D(M),H^i(M)`$ and $`K^i(M)`$ are always furnished with their natural gradings.
###### Reminder 2.2.
A) Let $`i`$. If $`U=_nU_n`$ is a graded $`S`$-module, we denote by $`{}_{}{}^{}\mathrm{Hom}_{k}^{}(U,k)`$ the graded $`S`$-module $`_n\mathrm{Hom}_k(U_n,k)`$. If $`M`$ is a finitely generated graded $`S`$-module, by graded local duality, we have isomorphisms of graded $`S`$-modules
(2.1) $`K^i(M)`$ $`{}_{}{}^{}\mathrm{Hom}_{k}^{}(H^i(M),k)\text{ and}`$
(2.2) $`H^i(M)`$ $`{}_{}{}^{}\mathrm{Hom}_{k}^{}(K^i(M),k)\mathrm{Hom}_S(K^i(M),E),`$
where $`E`$ denotes the graded injective envelope of the $`S`$-module $`k=S/S_+`$.
B) By $`0ptM`$ we denote the depth of the finitely generated graded $`S`$-module $`M`$ (with respect to the irrelevant ideal $`S_+`$ of $`S`$), so that
(2.3) $`\begin{array}{cc}\hfill 0ptM& =inf\{i|H^i(M)0\}\hfill \\ & =inf\{i|K^i(M)0\},\hfill \end{array}`$
(with the usual convention that $`inf\mathrm{}=\mathrm{}`$). Here $`0ptA`$ is called the arithmetic depth of the variety $`X_k^r`$. If we denote the Krull dimension of $`M`$ by $`dimM`$, we have
(2.4) $`\begin{array}{cc}\hfill dimM& =sup\{i|H^i(M)0\}\hfill \\ & =sup\{i|K^i(M)0\},\hfill \end{array}`$
(with the conventions that $`sup\mathrm{}=\mathrm{}`$ and $`dim0=\mathrm{}`$).
C) For a graded $`S`$-module $`U=_nU_n`$, let $`\text{end }U:=sup\{nU_n0\}`$ and $`\mathrm{beg}U:=inf\{nU_n0\}`$ denote the end resp. the beginning of $`U`$. In these notation, the Castelnuovo-Mumford regularity of the finitely generated graded $`S`$-module $`M`$ is defined by
(2.5)
$$\mathrm{reg}M=sup\{\text{ end }H^i(M)+i|i\}=inf\{\mathrm{beg}K^i(M))+i|i\}.$$
Keep in mind that the Castelnuovo-Mumford regularity of the variety $`X_k^r`$ is defined as
(2.6)
$$\mathrm{reg}X=\mathrm{reg}I=\mathrm{reg}A+1.$$
We are particularly interested in the canonical module of $`A`$, that is in the graded $`A`$-module
(2.7)
$$K(A):=K^{dim(A)}(A)=K^{d+1}(A).$$
###### Remark 2.3.
A) Let $`0<i<dim(A)=d+1`$ and let $`๐ญ\mathrm{Spec}S`$ with $`dimS/๐ญ=i`$. Then, the $`S_๐ญ`$-module $`A_๐ญ`$ has positive depth and hence vanishes or is of projective dimension $`<dimS_๐ญ=r+1i`$. Therefore $`K^i(A)_๐ญ\mathrm{Ext}^{r+1i}(A_๐ญ,S_๐ญ)=0`$. So
(2.8)
$$dimK^i(A)<i\text{ for }\mathrm{\hspace{0.17em}0}<i<dim(A)=d+1.$$
B) Let $`n`$ and let $`fA_n\backslash \{0\}.`$ Then, $`f`$ is $`A`$-regular and the short exact sequence $`0A(n)\stackrel{๐}{}AA/fA0`$ yields an epimorphism of graded $`A`$-modules $`f:H^{d+1}(A)(n)H^{d+1}(A)`$. So, by the isomorphisms (2.1) of Reminder 2.2, the multiplication map $`f:K^{d+1}(A)K^{d+1}(A)(n)`$ is injective. Moreover, localizing at the prime ideal $`IS`$ we get
$$K^{d+1}(A)_A\text{Quot}(A)K^{d+1}(A)_I\mathrm{Ext}_{S_I}^{rd}(S_I/IS_I,S_I)S_I/IS_I=\text{Quot }(A).$$
So, we may resume:
(2.9) The canonical module $`K(A)`$ of $`A`$ is torsion free and of rank $`1`$.
C) Let $`\mathrm{}S_1\backslash \{0\}`$ be a linear form. We write $`T:=S/\mathrm{}S`$ and consider $`T`$ as a polynomial ring in $`r`$ indeterminates. For the $`T`$-deficiency modules $`K_T^i(A/\mathrm{}A)`$ of $`A/\mathrm{}A`$, the isomorphisms (2.1) of Reminder 2.2 together with the base ring independence of local cohomology furnish the following isomorphisms of graded $`A/\mathrm{}A`$-modules
$$K_T^i(A/\mathrm{}A){}_{}{}^{}\mathrm{Hom}_{k}^{}(H_{T_+}^i(A/\mathrm{}A),k){}_{}{}^{}\mathrm{Hom}_{k}^{}(K_{S_+}^i(A/\mathrm{}A),k)K_S^i(A/\mathrm{}A).$$
So for all $`i`$ we obtain
(2.10)
$$K_T^i(A/\mathrm{}A)K_S^i(A/\mathrm{}A)=\mathrm{Ext}_S^{r+1i}(A/\mathrm{}A,S(r1)).$$
D) Let $`\mathrm{}`$ be as above. If we apply $`\mathrm{Ext}_S^{r+1i}(,S(r1))`$ to the short exact sequence $`0A(1)\stackrel{}{}AA/\mathrm{}A0`$ and keep in mind the isomorphisms (2.10), we get for each $`i`$ an exact sequence of graded $`A/\mathrm{}A`$-modules
(2.11)
$$0(K_S^{i+1}(A)/\mathrm{}K_S^{i+1}(A))(1)K_T^i(A/\mathrm{}A)0:_{K_S^i(A)}\mathrm{}0.$$
Correspondingly, applying local cohomology, we get for each $`i`$ an exact sequence of graded $`A/\mathrm{}A`$-modules
(2.12)
$$0H_{S_+}^i(A)/\mathrm{}H_{S_+}^i(A)H_{T_+}^i(A/\mathrm{}A)(0:_{H_{S_+}^{i+1}(A)}\mathrm{})(1)0.$$
E) We keep the above notation. In addition, we assume that $`\mathrm{}S_1\backslash \{0\}`$ is chosen generically. Then, according to Bertiniโs Theorem (cf ) the hyperplane section
$`Y:`$ $`=X\mathrm{Proj}(T)=\mathrm{Proj}(T/IT)\mathrm{Proj}(A/\mathrm{}A)\mathrm{Proj}(T)=_k^{r1}`$
is reduced and irreducible if $`dimA>2.`$ The homogeneous coordinate ring of $`Y`$ is
$$A^{}=A/(\mathrm{}A)^{\text{sat}}T/(IT)^{\text{sat}},$$
where $`^{\text{sat}}`$ is used to denote the saturation of a graded ideal in a homogeneous $`k`$-algebra. Observe that we have the following isomorphisms of graded $`A/\mathrm{}A`$-modules (cf (2.1), (2.10)).
(2.13) $`H_{S_+}^i(A/\mathrm{}A)`$ $`H_{T_+}^i(A/\mathrm{}A)H_T^i(A^{})\text{ for all }i>0;`$
(2.14) $`K_S^i(A/\mathrm{}A)`$ $`K_T^i(A/\mathrm{}A)K_T^i(A^{})\text{ for all }i>0.`$
On use of (2.12) and (2.13) we now easily get
(2.15)
$$H_{T_+}^i(A^{})_m=0H_{S_+}^{i+1}(A)_{m1}=0\text{ for all }i>0\text{ and all }m,$$
where, for a graded $`S`$-module $`U=_nU_n`$, we use $`U_n`$ to denote the $`m`$-th left truncation $`_{nm}U_n`$ of $`U`$. Finally, if $`0ptA>1`$, we have $`A^{}=A/\mathrm{}A`$. If $`0ptA=1`$, we know that $`H_{S_+}^1(A)`$ is a finitely generated non-zero $`A`$-module so that, by Nakayama, $`\mathrm{}H_{S_+}^1(A)H_{S_+}^1(A)`$ and hence $`H_{T_+}^1(A^{})0`$ (cf (2.12) and (2.13)). So, the arithmetic depth of $`Y`$ behaves as follows
(2.16)
$$0ptA^{}=\{\begin{array}{cc}0ptA1,\hfill & \text{if depth }A>1,\hfill \\ 1,\hfill & \text{if depth }A=1.\hfill \end{array}$$
The aim of the present paper is to investigate the case in which the degree of $`X`$ exceeds the codimension of $`X`$ by $`2`$. Keep in mind, that the degree of $`X`$ always exceeds the codimension of $`X`$ by 1. Therefore, we make the following convention.
###### Convention 2.4.
We write $`dimX,\mathrm{codim}X`$ and $`\mathrm{deg}X`$ for the dimension, the codimension and the degree of $`X`$ respectively, so that $`d=dimX=dimA1,\mathrm{codim}X=0ptI=rdimX=rd`$. Keep in mind that
$$\mathrm{deg}X\mathrm{codim}X+1$$
(cf e.g. ). We say that $`X`$ is of almost minimal degree, if $`\mathrm{deg}X=\mathrm{codim}X+2=rd+2.`$ Note that $`X`$ is called of minimal degree (cf ) whenever $`\mathrm{deg}X=\mathrm{codim}X+1.`$
We now discuss the case in which $`X`$ is a curve of almost minimal degree.
###### Remark 2.5.
A) We keep the hypotheses and notations of Remark 2.3 and assume that $`dimX=1`$ and that $`\mathrm{deg}X=\mathrm{codim}X+2=r+1`$. Then, for a generic linear form $`\mathrm{}S_1\backslash \{0\}`$ and in the notation of part E) of Remark 2.3, the generic hyperplane section
$`Y:`$ $`=\mathrm{Proj}(T/IT)\mathrm{Proj}(A/\mathrm{}A)=\mathrm{Proj}(A^{})\mathrm{Proj}(T)=_k^{r1}`$
is a scheme of $`r+1`$ points in semi-uniform position in $`_k^{r1}`$ (cf , ). Consequently, by (cf ) we can say that $`IT`$ is generated by quadrics. Therefore we may conclude: The homogeneous $`T`$-module
(2.17)
$$H_{T_+}^0((IT)^{\text{sat}}/IT)H_T^0(A/\mathrm{}A)\text{ is generated in degree }2\text{.}$$
Moreover, (cf \[4, (2.4) a)\])
(2.18)
$$dim_kH_{T_+}^1(A/\mathrm{}A)_n=dim_kH_{T_+}^1(A^{})_n=\{\begin{array}{cc}r+1,\hfill & \text{ if }n<0\hfill \\ r,\hfill & \text{ if }n=0\hfill \\ 1,\hfill & \text{ if }n=1\hfill \\ 0,\hfill & \text{ if }n>1\hfill \end{array}$$
So, by the exact sequences (2.12) and by statement (2.15) we get
(2.19) $`H^1(A)/\mathrm{}H^1(A)`$ $`k(1)\text{ and}`$
(2.20) $`\text{end }H^2(A)`$ $`0.`$
B) Assume first, that $`A`$ is a Cohen-Macaulay ring. Then $`H^1(A)=0`$ and so, by (2.20), the Hilbert polynomial $`P_A(x)[x]`$ of $`A`$ satisfies $`P_A(n)=dim_kA_n`$ for all $`n>0`$. As $`dim_kA_1=r+1`$ it follows $`P_A(x)=(r+1)x`$ and hence $`H^2(A)_nA_n`$ for all $`n`$. So, by (2.1) $`K(A)_n=K^2(A)_nA_n`$ for all $`n`$. As $`K(A)`$ is torsion-free of rank $`1`$ (cf (2.9)), we get an isomorphism of graded $`A`$-modules $`K(A)A(0)`$. Therefore, $`A`$ is a Gorenstein ring.
If $`A`$ is normal, $`X_k^r`$ is a smooth non-degenerate curve of genus $`dim_kK(A)_0=1`$ and of degree $`r+1`$, hence an elliptic normal curve: we are in the case $`\overline{I}`$ of \[4, (4.7) B)\].
C) Yet assume that $`A`$ is a Cohen-Macaulay (and hence a Gorenstein) ring. Assume that $`A`$ is not normal. Let $`B`$ denote the graded normalization of $`A`$. Then, there is a short exact sequence of graded $`S`$-modules $`0AB\stackrel{๐}{}C0`$ with $`dimC=1`$. As $`H_{S_+}^0(B)=H_{S_+}^1(B)=H_{S_+}^0(A)=H_{S_+}^1(A)=0`$ we get $`H_{S_+}^0(C)=0`$ and an exact sequence of graded $`S`$-modules
$$0H_{S_+}^1(C)H^2(A)H_{S_+}^2(B)0.$$
As $`dimC=1`$ and $`H_{S_+}^0(C)=0`$, there is some $`c`$ such that
$$dim_kC_n+dim_kH_{S_+}^1(C)_n=c\text{ for all }n.$$
By (2.20) and the above sequence $`dim_kC_n=c`$ for all $`n>0`$. As $`C_0=0`$ and $`dim_kH^2(A)_0=dim_kA_0=1`$ (cf part B) ) it follows $`c=1`$. As $`H_{S_+}^0(C)=0`$, there exits a $`C`$-regular element $`hS_1\{0\},`$ and choosing $`tC_1\{0\}`$ we get
(2.21)
$$C=k[h]tk[h](1).$$
Choose $`\overline{y}B_1`$ such that $`\pi (\overline{y})=t.`$ Then we get $`B/A=(\overline{y}A+A)/A`$ and hence $`B=A[\overline{y}]`$. So, if $`y`$ is an indeterminate, there is a surjective homomorphism of homogeneous $`k`$-algebras
$$S[y]=k[x_0,\mathrm{},x_r,y]\stackrel{๐ฝ}{}B,y\overline{y},$$
which occurs in the commutative diagram
where $`\alpha `$ is the natural map. Thus, the normalization $`\stackrel{~}{X}:=\mathrm{Proj}(B)`$ of $`X`$ is a curve of degree $`r+1`$ in $`\mathrm{Proj}(S[y])=_k^{r+1}`$ โ a rational normal curve โ and the normalization morphism $`\nu :\stackrel{~}{X}X`$ is induced by a simple projection $`\varrho :^{r+1}\backslash \{p\}_k^r`$ with center $`\{p\}=|\mathrm{Proj}(S[y]/S_+S[y])|`$.
Moreover, by (2.21) we have $`\nu _{}๐ช_{\stackrel{~}{X}}/๐ช_X\stackrel{~}{C}k`$, so that $`\nu _{}๐ช_{\stackrel{~}{X}}/๐ช_X`$ is supported in a single point $`qX`$, โ the unique singularity of $`X`$ โ a double point. That is, we are in the case $`\overline{\text{III}}`$ of \[4, (4.7) B)\].
D) We keep the notations and hypotheses of part A). But contrary to what we did in parts B) and C) we now assume that $`A`$ is not Cohen-Macaulay, so that $`H^1(A)0`$. Then, by (2.19) and by Nakayama it follows $`H^1(A)/\mathrm{}H^1(A)=k(1)`$. In particular $`H^1(A)_1k`$ and the multiplication map $`\mathrm{}:H^1(A)_nH^1(A)_{n+1}`$ is surjective for all $`n1`$.
Now we claim that $`H^1(A)_n=0`$ for all $`n>1`$. Assuming the opposite, we would have an isomorphism $`H^1(A)_1\stackrel{}{}H^1(A)_2`$ and the exact sequence of graded $`S`$-modules $`0H_{S_+}^0(A/\mathrm{}A)H^1(A)(1)\stackrel{}{}H^1(A)`$ would imply that $`H_{T_+}^0(A/\mathrm{}A)_2H_{S_+}^0(A/\mathrm{}A)_2=0`$ and hence $`H_{T_+}^0(A/\mathrm{}A)=0`$ (cf (2.17)). This would imply $`0ptA>1`$, a contradiction. This proves our claim and shows (cf (2.1))
(2.22)
$$H^1(A)k(1)\text{ and }K^1(A)k(1).$$
By (2.18) (applied for $`n=1`$) it follows that the natural map $`H^1(A)_1H^1(A/\mathrm{}A)_1`$ is an isomorphism. So $`H^2(A)_0=0`$. In particular, we get
(2.23)
$$P_A(x)=(r+1)x+1,\text{ end }H^2(A)=1,$$
where $`P_A(x)[x]`$ is used to denote the Hilbert polynomial of $`A`$.
As $`K(A)`$ is torsion-free over the $`2`$-dimensional domain $`A`$ (cf (2.9)) and satisfies the second Serre property $`S_2`$ (cf \[28, 3.1.1\]), in view of the second statement of (2.23) we get :
(2.24)
$$K(A)\text{ is a }CM\text{-module with beg }K(A)=1.$$
According to (2.22), the $`S_+`$-transform $`D(A)`$ of $`A`$ is a domain which appears in a short exact sequence $`0AD(A)k(1)0`$. Choosing $`\overline{y}D(A)_1A_1`$ we obtain $`D(A)=A[\overline{y}]`$. So, if $`y`$ is an indeterminate, there is a surjective homomorphism of homogeneous $`k`$-algebras $`S[y]=k[x_0,\mathrm{},x_r,y]\stackrel{๐พ}{}D(A)`$, sending $`y`$ to $`\overline{y}`$ and extending the natural map $`\alpha :SA`$ (cf part C) ). In particular $`\stackrel{~}{X}:=\mathrm{Proj}(D(A))`$ is a curve of degree $`r+1`$ in $`\mathrm{Proj}(S[y])=_k^{r+1}`$ โ a rational normal curve. Moreover, the natural morphism $`\epsilon :\stackrel{~}{X}X`$ is an isomorphism induced by the simple projection $`\varrho :_k^{r+1}\backslash \{p\}_k^r`$ with $`\{p\}=|\mathrm{Proj}(S[y]S_+S[y])|(\stackrel{~}{X}).`$ That is, we are in the case $`\overline{\text{II}}`$ of \[4, (4.7) B)\].
## 3. The case โArithmetic Depth $`=1`$
In this section we study varieties of almost minimal degree and arithmetic depth one. In particular, we shall extend the results of part D) of Remark 2.5 from curves to higher dimensions.
###### Proposition 3.1.
Let $`X_k^r`$ be a projective variety of almost minimal degree such that $`0ptA=1`$ and $`dimX=d`$. Then
* $`H^i(A)=K^i(A)=0`$ for all $`i1,d+1`$;
* $`\mathrm{end}H^{d+1}(A)=\mathrm{beg}K(A)=d`$;
* $`H^1(A)k(1),K^1(A)k(1)`$;
* $`K(A)`$ is a torsion-free $`CM`$-module of rank one;
* $`D(A)`$ is a homogeneous $`CM`$ integral domain with $`\mathrm{reg}D(A)=1`$ and $`dim_kD(A)_1=r+2`$.
###### Proof.
(Induction on $`d=dimX`$). For $`d=1`$ all our claims are clear by the results of part D) of Remark 2.5.
So, let $`d>1`$. Let $`\mathrm{}S_1\backslash \{0\}`$ be generic. Then in the notation of part E) of Remark 2.3 we have $`dimA^{}=d`$ and $`0ptA^{}=1`$ (cf (2.16)). By induction $`H_{T_+}^i(A^{})=0`$ for all $`i1,d`$. So, by (2.12) we obtain
(3.1)
$$H^i(A)=0,\text{ for all }i1,2,d+1.$$
Moreover, by induction and in view of (2.13) we get $`H_{T_+}^1(A/\mathrm{}A)k(1)`$. As $`H^1(A)`$ is a non-zero and finitely generated graded $`S`$-module, we have $`\mathrm{}H^1(A)H^1(A)`$. So, by (2.12) we obtain
(3.2)
$$H^1(A)/\mathrm{}H^1(A)k(1)$$
and
(3.3)
$$H^2(A)=0.$$
Combining (3.1), (3.3) and (2.1), we get claim (a). By induction
$$\text{ end }H^d(A/\mathrm{}A)=\text{ end }H_{T_+}^d(A^{})=d+1.$$
As $`H^d(A)=0`$, (2.12) gives $`\text{ end }H^{d+1}(A)=d`$. In view of (2.1) we get claim (b). Also, by induction $`0ptK_T^d(A^{})=d`$. As $`d>0`$, we have $`H_{T_+}^d(A^{})H_{T_+}^d(A/\mathrm{}A)`$ and hence $`K_T^d(A/\mathrm{}A)K_T^d(A^{})`$, (cf (2.1)). As $`K^d(A)=0`$, (2.9) and (2.11) prove statement (d). Moreover $`D(A)`$ is a positively graded finite integral extension domain of $`A`$ such that $`H_{S_+}^1(D(A))=0`$ and $`H_{S_+}^i(D(A))H^i(A)`$ for all $`i>1`$, it follows from statements (a) and (b), that $`D(A)`$ is a $`CM`$-ring with $`\mathrm{reg}D(A)=1`$. In view of (3.2) and the natural exact sequence
(3.4)
$$0A\stackrel{๐}{}D(A)\stackrel{๐}{}H^1(A)0$$
there is some $`\delta D(A)_1\backslash A`$ such that $`D(A)=A+\delta A`$. In particular we have $`D(A)=A[\delta ]`$ and $`D(A)_1Ak.`$ Therefore statement (e) is proved.
It remains to show statement (c). In view of (2.1) it suffices to show that $`H^1(A)k(1)`$. By (3.2) and as $`H^1(A)_n=0`$ for all $`n0`$, there is an isomorphism of graded $`S`$-modules $`H^1(A)S/๐ฎ(1)`$, where $`๐ฎS`$ is a graded $`S_+`$-primary ideal. We have to show that $`๐ฎ=S_+`$. There is a minimal epimorphism of graded $`S`$-modules
$$\pi :SS(1)D(A)0$$
such that $`\pi _S`$ coincides with the natural map $`\alpha :SA`$ and $`\pi (S(1))=\delta A=\delta S`$. As $`\mathrm{beg}(\text{Ker }(\alpha )=I)2`$ and $`\pi (S_1)\pi (S(1)_1)=A_1\delta k=0`$, it follows $`\mathrm{beg}\text{Ker }(\pi )2`$. Moreover, by statement (e) we have $`\mathrm{reg}D(A)=1`$. Therefore a minimal free presentation of $`D(A)`$ has the form
(3.5)
$$S^\beta (2)SS(1)\stackrel{\pi }{}D(A)0$$
with $`\beta _0`$. It follows $`\mathrm{Tor}_1^S(k,D(A))k^\beta (2)`$. As $`1=\eta (1)`$ is a minimal generator of the $`S`$-module $`D(A)`$, the sequence (3.4) induces an epimorphism of graded $`S`$-modules
$$\mathrm{Tor}_1^S(k,D(A))\mathrm{Tor}_1^S(k,H^1(A))0.$$
Therefore
$$(๐ฎ/S_+๐ฎ)(1)\mathrm{Tor}_0^S(k,๐ฎ(1))\mathrm{Tor}_1^S(k,(S/๐ฎ)(1))\mathrm{Tor}_1^S(k,H^1(A))$$
is concentrated in degree $`2`$. So, by Nakayama, $`๐ฎ`$ is generated in degree one, thus $`๐ฎ=S_+`$. โ
Varieties of almost minimal degree and arithmetic depth one can be characterized as simple generic projections from varieties of minimal degree.
###### Reminder 3.2.
A) Recall that an irreducible reduced non-degenerate projective variety $`\stackrel{~}{X}_k^s`$ is said to be of minimal degree if $`\text{deg }\stackrel{~}{X}=\text{codim }\stackrel{~}{X}+1`$.
B) Projective varieties of minimal degree are rather well understood, namely (cf e.g. \[21, Theorem 19.9\]): A projective variety $`\stackrel{~}{X}_k^s`$ of minimal degree is either
(3.6) $`\text{a quadric hypersurface },`$
(3.7) a (cone over a) Veronese surface in $`_k^5`$ or
(3.8) $`\text{a (cone over a) rational normal scroll}.`$
C) In particular, a variety $`\stackrel{~}{X}_k^s`$ of minimal degree is arithmetically Cohen-Macaulay and arithmetically normal.
###### Remark 3.3.
A) Let $`\stackrel{~}{X}_k^s`$ be an irreducible reduced projective variety, let $`p_k^s\backslash \stackrel{~}{X}`$, let $`\varrho :_k^s\backslash \{p\}_k^{s1}`$ be a projection with center $`p`$ and let $`X:=\varrho (\stackrel{~}{X})_k^{s1}`$. Then, the induced morphism $`\varrho :\stackrel{~}{X}X`$ is finite. Moreover, we have $`\mathrm{deg}X=\mathrm{deg}\stackrel{~}{X}`$ if and only if $`\varrho `$ is birational, hence if and only if there is a line $`\overline{\mathrm{}}_k^s`$ with $`p\overline{\mathrm{}}`$ and such that the scheme $`\overline{\mathrm{}}\stackrel{~}{X}`$ is non-empty, reduced and irreducible. It is equivalent to say that there are lines $`\overline{\mathrm{}}_k^s`$ which join $`p`$ and $`\stackrel{~}{X}`$ and are not secant lines of $`\stackrel{~}{X}`$.
But this means precisely that the join $`\text{Join}(p,\stackrel{~}{X})`$ of $`p`$ and $`\stackrel{~}{X}`$ is not contained in the secant cone $`\mathrm{Sec}_p(\stackrel{~}{X})`$ of $`\stackrel{~}{X}`$ of $`p`$. Observe that here $`\mathrm{Sec}_p(\stackrel{~}{X})`$ is understood as the union of $`p`$ with all lines $`\overline{\mathrm{}}_k^s`$ such that $`p\overline{\mathrm{}}`$ and $`\overline{\mathrm{}}\stackrel{~}{X}`$ is a scheme of dimension $`0`$ and of degree $`>1`$.
Also, $`\varrho `$ is an isomorphism if and only if for any line $`\overline{\mathrm{}}_k^s`$ with $`p\overline{\mathrm{}}`$, the scheme $`\overline{\mathrm{}}\stackrel{~}{X}`$ is either empty or reduced and irreducible. It is equivalent to say that $`p\mathrm{Sec}(\stackrel{~}{X})`$, where the secant variety $`\mathrm{Sec}(\stackrel{~}{X})`$ of $`\stackrel{~}{X}`$ is understood as the union of all lines $`\overline{\mathrm{}}_k^s`$ such that $`\overline{\mathrm{}}\stackrel{~}{X}`$ is a scheme of dimension $`0`$ and of degree $`>1`$, or else $`\overline{\mathrm{}}\stackrel{~}{X}`$.
B) Assume now in addition that $`\stackrel{~}{X}_k^s`$ is of minimal degree. Then by the above observation we can say:
(3.9)
$$X_k^{s1}\text{ is of almost minimal degree if and only if }\text{Join}(p,\stackrel{~}{X})\mathrm{Sec}_p(\stackrel{~}{X})\text{.}$$
(3.10) $`\begin{array}{cc}& X_k^{s1}\text{ is of almost minimal degree and }\varrho :\stackrel{~}{X}X\text{an isomorphism }\hfill \\ & \text{if and only if }p\mathrm{Sec}(\stackrel{~}{X}),\text{ thus if and only if }\mathrm{Sec}_p(\stackrel{~}{X})=\{p\}.\hfill \end{array}`$
Moreover, by the classification of Reminder 3.2 B) it follows that $`\mathrm{Sec}(\stackrel{~}{X})=_k^s`$ whenever $`\stackrel{~}{X}`$ is not smooth.
Now, we can give the announced geometric characterization of varieties of almost minimal degree and arithmetic depth one.
###### Proposition 3.4.
The following statements are equivalent:
* $`X`$ is of almost minimal degree and of arithmetic depth $`1`$.
* $`X`$ is the projection $`\varrho (\stackrel{~}{X})`$ of a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree from a point $`p_k^{r+1}\backslash \mathrm{Sec}(\stackrel{~}{X})`$.
###### Proof.
(i) $``$ (ii): Assume that $`X`$ is of almost minimal degree with $`0ptA=1`$. Then, by statement e) Proposition 3.1, there is some $`\overline{y}D(A)_1\backslash A_1`$ such that $`D(A)=A[\overline{y}]`$. Now, as in the last paragraph of part D) in Remark 2.5, we may view $`\stackrel{~}{X}:=\text{ Proj}(D(A))`$ as a non-degenerate irreducible projective variety in $`_k^{r+1}`$ such that $`\mathrm{deg}(\stackrel{~}{X})=\mathrm{deg}(X)`$ and a projection $`\varrho :_k^{r+1}\backslash \{p\}_k^r`$ from an appropriate point $`p_k^{r+1}\backslash \stackrel{~}{X}`$ induces an isomorphism $`\varrho :\stackrel{~}{X}X`$. In view of Remark 3.3 B) this proves statement (ii).
(ii) $``$ (i): Assume that there is a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree and a projection $`\varrho :_k^{r+1}\backslash \{p\}_k^r`$ from a point $`p\text{ Sec}(\stackrel{~}{X})`$ with $`X=\varrho (\stackrel{~}{X})`$. Then, by (LABEL:3.10), $`X`$ is of almost minimal degree and $`\varrho :\stackrel{~}{X}X`$ is an isomorphism. It remains to show that $`0ptA=1`$, hence that $`H^1(A)0`$. So, let $`B`$ denote the homogeneous coordinate ring of $`\stackrel{~}{X}_k^{r+1}`$. Then, the isomorphism $`\varrho :\stackrel{~}{X}X`$ leads to an injective homomorphism of graded integral domains $`AB`$ such that $`B/A`$ is an $`A`$-module of finite length. Therefore $`B_n(A:_BS_+^n)=D(A)`$. As $`\stackrel{~}{X}_k^{r+1}`$ is non-degenerate, we have $`dim_kB_1=r+2>dim_kA_1`$, hence $`ABD(A)`$ so that $`H^1(A)D(A)/A0`$. โ
## 4. The non-arithmetically Cohen-Macaulay case
In this section we study projective varieties of almost minimal degree which are not arithmetically Cohen-Macaulay. So, we are interested in the case where $`\mathrm{deg}X=\mathrm{codim}X+2`$ and $`10ptAdimX`$.
Our first aim is to generalize Proposition 3.1. In order to do so, we prove the following auxiliary result, in which $`\mathrm{NZD}_S(M)`$ is used to denote the set of non-zero divisors in $`S`$ with respect to the $`S`$-module $`M`$.
###### Lemma 4.1.
Let $`M`$ be a finitely generated graded $`S`$-module, let $`m\{0,\mathrm{},r\}`$ and $`n`$. Let $`z_m,\mathrm{},z_rS_1`$ be linearly independent over $`k`$ such that $`z_m\mathrm{NZD}_s(M)`$ and such that there is an isomorphism of graded $`S`$-modules $`M/z_mM(S/(z_m,\mathrm{},z_r))(n)`$.
Then, there are linearly independent elements $`y_{m+1},\mathrm{},y_rS_1`$ such that there is an isomorphism of graded $`S`$-modules $`M(S/(y_{m+1},\mathrm{},y_r))(n).`$
###### Proof.
By Nakayama there is an isomorphism of graded $`S`$-modules $`M(S/๐ฎ)(n)`$, where $`๐ฎS`$ is a homogeneous ideal. In particular
(4.1)
$$z_m\mathrm{NZD}_S(S/๐ฎ).$$
As $`S/(z_m,๐ฎ)(M/z_mM)(n)S/(z_m,\mathrm{},z_r)`$, we have
$$(z_m,๐ฎ)=(z_m,z_{m+1},\mathrm{},z_r).$$
Also, by (4.1), we have $`z_m๐ฎ_1`$, so that $`๐ฎ_1`$ becomes a $`k`$-vector space of dimension $`rm`$. Let $`y_{m+1},\mathrm{},y_rS_1`$ form a $`k`$-basis of $`๐ฎ_1`$. As
$$(y_{m+1},\mathrm{},y_r)๐ฎ(z_m,y_{m+1},\mathrm{},y_r)$$
and in view of (4.1) we obtain
$$๐ฎ=(y_{m+1},\mathrm{},y_r)+๐ฎz_mS=(y_{m+1},\mathrm{},y_r)+z_m๐ฎ.$$
So, by Nakayama $`๐ฎ=(y_{m+1},\mathrm{},y_r)`$. โ
Now, we are ready to prove the first main result of this section, which recover results of written down in the context of the modules of deficiency.
###### Theorem 4.2.
Assume that $`X_k^r`$ is of almost minimal degree and that $`t:=0ptAdimX=:d`$. Then
* $`H^i(A)=K^i(A)=0`$ for all $`it,d+1`$;
* $`\text{end }H^{d+1}(A)=\mathrm{beg}K(A)=d`$;
* There are linearly independent forms $`y_{t1},\mathrm{},y_rS_1`$ such that there is an isomorphism of graded $`S`$-modules
$$K^t(A)(S/(y_{t1},\mathrm{},y_r))(2t);$$
* $`K(A)`$ is a torsion-free CM-module of rank one.
###### Proof.
(Induction on $`t`$). The case $`t=1`$ is clear by Proposition 3.1.
So, let $`t>1`$ and $`\mathrm{}S_1\backslash \{0\}`$ be generic. Then, in the notation of Remark 2.3 E) we have $`A^{}=A/\mathrm{}A`$, $`dimA^{}=dimY+1=d`$ and $`0ptA^{}=t1d1=dimY`$ (cf (2.16)).
(a): By induction, $`H_{T_+}^j(A^{})=0`$ for all $`jt1,d`$. So, (2.15) gives $`H^i(A)=0`$ for all $`i0,1,t,d+1`$. As $`t>1`$, we have $`H^0(A)=H^1(A)=0`$. In view of (2.1) this proves statement (a).
(b): By induction $`\text{end }H_{T_+}^d(A/\mathrm{}A)=d+1`$. As $`H^d(A)=0`$, (2.12) implies that end $`H^{d+1}`$ $`(A)=d`$ and (2.1) gives our claim.
(c): By induction there are forms $`z_{t1},\mathrm{},z_rS_1`$ whose images $`\overline{z}_{t1},\mathrm{},\overline{z}_rT_1`$ are linearly independent over $`k`$ and such that there is an isomorphism of graded $`T`$-modules $`K^{t1}(A/\mathrm{}A)=(T/(\overline{z}_{t1},\mathrm{},\overline{z}_r))(2(t1))`$. Let $`z_{t2}:=\mathrm{}`$. Then $`z_{t2},\mathrm{},z_rS_1`$ are linearly independent and
$$K_T^{t1}(A/\mathrm{}A)(S/(z_{t2},z_{t1},\mathrm{},z_r))(3t).$$
By statement (a) we have $`K^{t1}(A)=0.`$ So, the sequence (2.11) gives an isomorphism of graded $`S`$-modules
(4.2)
$$K^t(A)/z_{t2}K^t(A)(S/(z_{t2},\mathrm{},z_r))(2t).$$
Assume first, that $`t<d`$. By induction $`K_T^t(A/z_{t2}A)=K_T^t(A^{})`$ vanishes and hence (2.11) yields $`0:_{K^t(A)}z_{t2}=0`$, thus $`z_{t2}\text{NZD}_S(K^t(A))`$. So, (4.2) and Lemma 4.1 imply statement (c).
Now, let $`t=d`$. Then (4.2) implies $`dimK^d(A)/z_{t2}K^d(A)=d2`$. Our first aim is to show that $`dimK^d(A)=d1`$. If $`d>2,`$ this follows by the genericity of $`z_{t2}=\mathrm{}.`$ So, let $`t=d=2.`$ Then $`A/\mathrm{}A`$ is a domain of depth $`1`$ which is the coordinate ring of a curve $`Y_k^{r1}`$ of almost minimal degree (cf Remark 2.5 A)). So, according to Remark 2.5 D) we have $`H^1(A/\mathrm{}A)k(1)`$ and $`H^2(A/\mathrm{}A)_n=0`$ for all $`n0.`$ If we apply cohomology to the exact sequence $`0A(1)\stackrel{\mathrm{}}{}AA/\mathrm{}A0`$ and keep in mind that $`H^1(A)=0`$ we thus get $`\mathrm{}:H^2(A)_1\stackrel{}{}H^2(A)_0k.`$ According to Remark 2.5 D) there is an isomorphism $`\stackrel{~}{Y}\stackrel{}{}Y,`$ where $`\stackrel{~}{Y}_k^{r1}`$ is a rational normal curve, so that $`Y_k^1`$ is smooth. As $`Y`$ is a hyperplane section of $`X`$ it follows that the non-singular locus of $`X`$ is finite. So, if we apply \[1, Proposition 5.2\] to the ample sheaf of $`๐ช_X`$-modules $`:=๐ช_X(1)`$ and observe that $`H^2(A)_nH^1(X,^n)`$ for all $`n,`$ we get that $`H^2(A)_nk`$ for all $`n0.`$ Consequently, $`K^2(A)_n0`$ for all $`n0,`$ hence $`dimK^2(A)>0=dimK^2(A)/z_0K^2(A).`$ Therefore $`dimK^2(A)=1,`$ which concludes the case $`t=d.`$
According to (4.2) the $`S`$-module $`K^d(A)/z_{t2}K^d(A)`$ is generated by a single homogeneous element of degree $`d2`$. By Nakayama, $`K^d(A)`$ has the same property. So, there is a graded ideal $`๐ฎS`$ with $`K^d(A)(S/๐ฎ)(2d)`$. In particular we have $`dimS/๐ฎ=d1`$.
Now, another use of (4.2) yields
$$\begin{array}{c}S/(๐ฎ,z_{d2})(S/๐ฎ)/z_{d2}(S/๐ฎ)K^t(A)(d2)/z_{d2}K^t(A)(d2)\hfill \\ \hfill (K^t(A)/z_{d2}K^t(A))(d2)S/(z_{t2},\mathrm{},z_r),\end{array}$$
so that $`(๐ฎ,z_{d2})=(z_{t2},\mathrm{},z_r)`$ is a prime ideal. As
$$dimS/๐ฎ=d1>dim(S/(z_{t2},\mathrm{},z_r))$$
it follows, that $`๐ฎ`$ is a prime ideal. Moreover, as $`z_{d2}๐ฎ`$, we obtain $`z_{d2}\text{NZD}_S(S/๐ฎ)=\text{NZD}_S(K^d(A))`$.
Now, our claim follows from (4.2) and Lemma 4.1.
(d): In view (2.9) it remains to show that $`0ptK(A)=d+1`$. By (2.9) and by induction we have
(4.3)
$$\mathrm{}\mathrm{NZD}_S(K(A))\text{ and }0ptK_T^d(A/\mathrm{}A)=d.$$
So, by the sequence (2.11), applied with $`i=d`$, it suffices to show that $`0:_{K^d(A)}\mathrm{}=0`$. If $`t<d`$, this last equality follows from statement (a). If $`t=d`$, statement (c) yields $`0ptK^d(A)=d1>0`$ and by the genericity of $`\mathrm{}`$ we get $`\mathrm{}\text{NZD}_S(K^d(A))`$. โ
###### Remark 4.3.
Keep the notations and hypotheses of Theorem 4.2. Then, by statement (c) of Theorem 4.2 and in view of (2.1) we get $`\text{end }H^t(A)=2t`$. So, by statements (a) and (b) of Theorem 4.2 we obtain
(4.4)
$$\mathrm{reg}(A)=2\text{ and }dim_kA_n=P_A(n),\text{ for all }n>2t.$$
###### Corollary 4.4.
Let $`X_k^r`$ be of almost minimal degree with $`dimX=d`$ and $`0ptA=t`$. Then:
* The Hilbert series of $`A`$ is given by
$$F(\lambda ,A)=\frac{1+(r+1d)\lambda }{(1\lambda )^{d+1}}\frac{\lambda }{(1\lambda )^{t1}}.$$
* The Hilbert polynomial of $`A`$ is given by
$$P_A(n)=(rd+2)\left(\genfrac{}{}{0pt}{}{n+d1}{d}\right)+\left(\genfrac{}{}{0pt}{}{n+d1}{d1}\right)\left(\genfrac{}{}{0pt}{}{n+t2}{t2}\right).$$
* The number of independent quadrics in $`I`$ is given by
$$dim_k(I_2)=t+\left(\genfrac{}{}{0pt}{}{r+1d}{2}\right)d2.$$
###### Proof.
(a): (Induction on $`t`$). If $`t=1`$, $`D(A)`$ is a CM-module of regularity $`1`$ (cf Proposition 3.1 (e) ) and of multiplicity $`\mathrm{deg}X=rd+2`$. Therefore
(4.5)
$$F(\lambda ,D(A))=\frac{1+(r+1d)\lambda }{(1\lambda )^{d+1}}.$$
In view of statement c) of Proposition 3.1 we thus get
$$F(\lambda ,A)=F(\lambda ,D(A))\lambda =\frac{1+(r+1d)\lambda }{(1\lambda )^{d+1}}\lambda $$
and hence our claim.
So, let $`t>1`$. Then, as $`t^{}=t1`$ (cf (2.16)) and $`A^{}=A/\mathrm{}A`$ we get by induction
$`F(\lambda ,A)={\displaystyle \frac{F(\lambda ,A^{})}{1\lambda }}`$ $`=[{\displaystyle \frac{1+((r1)+1(d1))\lambda }{(1\lambda )^d}}{\displaystyle \frac{\lambda }{(1\lambda )^{t2}}}](1\lambda )^1`$
$`={\displaystyle \frac{1+(r+1d)\lambda }{(1\lambda )^{d+1}}}{\displaystyle \frac{\lambda }{(1\lambda )^{t1}}}.`$
(b), (c): These are purely arithmetical consequences of statement (a). โ
###### Remark 4.5.
Observe that Corollary 4.4 also holds if $`X`$ is arithmetically Cohen-Macaulay. In this case, the shape of the Hilbert series $`F(\lambda ,A)`$ (cf statement (a) ) yields that $`A`$ is a Gorenstein ring (cf ) which says that a projective variety of almost minimal degree which is arithmetically Cohen-Macaulay is already arithmetically Gorenstein. For $`dimX=0`$ this may be found in .
Finally, by Remark 2.5 B), by statement (2.9) and the exact sequence (2.12) it follows immediately by induction on $`d=dimX`$ that $`K(A)A(1d)`$ if $`X`$ is is arithmetically Cohen-Macaulay. This shows again that $`X`$ is arithmetically Gorenstein.
## 5. Endomorphism Rings of Canonical Modules
Or next aim is to extend the geometric characterization of Proposition 3.4 to arbitrary non-arithmetically Cohen-Macaulay varieties of almost minimal degree.
We attack this problem via an analysis of the properties of the endomorphism ring of the canonical module $`K(A)`$ of $`A`$, which in the local case has been studied already in . The crucial point is, that this ring has a geometric meaning in the context of varieties of almost minimal degree.
###### Notation 5.1.
We write $`B`$ for the endomorphism ring of the canonical module of $`A`$, thus
$$B:=\mathrm{Hom}_S(K(A),K(A)).$$
Observe that $`B`$ is a finitely generated graded $`A`$-module and that
(5.1)
$$B=\mathrm{Hom}_A(K(A),K(A)).$$
In addition we have a homomorphism of graded $`A`$-modules
(5.2)
$$\epsilon :AB,aa\text{ id}_{K(A)}.$$
Keep in mind, that $`B`$ carries a natural structure of (not necessarily commutative) ring and that $`\epsilon `$ is a homomorphism of rings.
The homomorphism $`\epsilon :AB`$ occurs to be of genuine interest for its own. So we give a few properties of it.
###### Proposition 5.2.
Let $`d:=dimX1`$. Then
* $`B=kB_1B_2\mathrm{}`$ is a positively graded commutative integral domain of finite type over $`B_0=k`$.
* $`\epsilon :AB`$ is a finite injective birational homomorphism of graded rings.
* There is a (unique) injective homomorphism $`\stackrel{~}{\epsilon }`$ of graded rings, which occurs in the commutative diagram
* If $`๐ญ\mathrm{Spec}(A)`$, the ring $`A_๐ญ`$ has the second Serre property $`S_2`$ if and only if the localized map $`\epsilon _๐ญ:A_๐ญB_๐ญ`$ is an isomorphism.
* $`\epsilon :AB`$ is an isomorphism if and only if $`A`$ satisfies $`S_2`$.
* $`\stackrel{~}{\epsilon }:D(A)B`$ is an isomorphism if and only if $`X`$ satisfies $`S_2`$.
* $`B`$ satisfies $`S_2`$ (as an $`A`$-module and as a ring).
* If the $`A`$-module $`K(A)`$ is Cohen-Macaulay, then $`B`$ is Cohen-Macaulay (as an $`A`$-module and as a ring).
###### Proof.
(a), (b): By (2.9) (cf Remark 2.3) we know that $`K(A)`$ is torsion-free and of rank one. From this it follows easily that $`B`$ is a commutative integral domain. Also the map $`\epsilon :AB`$ is a homomorphism of $`A`$-modules, and so becomes injective by the torsion-freeness of the $`A`$-module $`K(A)`$. The intrinsic $`A`$-module structure on $`B`$ and the $`A`$-module structure induced by $`\epsilon `$ are the same. As $`B`$ is finitely generated as an $`A`$-module it follows that $`\epsilon `$ is a finite homomorphism of rings.
It is easy to verify that the natural grading of the $`A`$-module $`B`$ respects the ring structure on $`B`$ and thus turns $`B`$ into a graded ring. In particular $`\epsilon `$ becomes a homomorphism of graded rings. As $`A`$ is positively graded, $`\epsilon `$ is finite and $`B`$ is a domain, it follows that $`B`$ is finite. As $`k`$ is algebraically closed and $`B_0`$ is a domain, we get $`B_0k`$. As $`A`$ is of finite type over $`k`$ and $`\epsilon `$ is finite, $`B`$ is of finite type over $`k`$, too.
(c): As $`dimA>1`$ we know that $`H_{S_+}^1(A)`$ is of finite length. Therefore
$$\mathrm{Ext}_S^j(H_{S_+}^1(A),S)=0\text{ for all }jr+1.$$
So, the short exact sequence $`0AD(A)H^1(A)0`$ yields an isomorphism of graded $`A`$-modules
$$K(A)=\mathrm{Ext}_S^{rd}(A,S(r1))\mathrm{Ext}_S^{rd}(D(A),S(r1)).$$
Therefore, $`K(A)`$ carries a natural structure of graded $`D(A)`$-module. As $`D(A)`$ is a birational extension ring of $`A`$, we can write
$$B=\mathrm{Hom}_A(K(A),K(A))=\mathrm{Hom}_{D(A)}(K(A),K(A))$$
and hence consider $`B`$ as a graded $`D(A)`$-module in a natural way. In particular, there is a homomorphism of rings
$$\stackrel{~}{\epsilon }:D(A)B,cc\text{ id}_{K(A)},$$
the unique homomorphism of rings $`\stackrel{~}{\epsilon }`$ which appears in the commutative diagram
As $`A,D(A)`$ and $`B`$ are domains and as $`\epsilon `$ is injective, $`\stackrel{~}{\epsilon }`$ is injective, too. Clearly $`\stackrel{~}{\epsilon }`$ is finite, and respects gradings.
(d): Let $`๐ญ\mathrm{Spec}(A)`$. Then, by the chain condition in $`\mathrm{Spec}(S)`$, $`K(A)_๐ญ`$ is nothing else than the canonical module $`K_{A_๐ญ}`$ of the local domain $`A_๐ญ`$. In particular we may identify $`B_๐ญ=\mathrm{Hom}_A(K(A),K(A))_๐ญ\mathrm{Hom}_{A_๐ญ}(K_{A_๐ญ},K_{A_๐ญ})`$. Then, the natural map $`\epsilon _๐ญ:A_๐ญB_๐ญ`$ induced by $`\epsilon `$ coincides with the natural map
$$A_๐ญ\mathrm{Hom}_{A_๐ญ}(K_{A_๐ญ},K_{A_๐ญ}),bb\text{ id}_{K_{A_๐ญ}}.$$
But this latter map is an isomorphism if and only if $`A_๐ญ`$ satisfies $`S_2`$ (cf \[28, 3.5.2\]).
(e): Is clear by statement (d).
(f): By statement (d), $`X`$ satisfies $`S_2`$ if and only if $`\epsilon _๐ญ:A_๐ญB_๐ญ`$ is an isomorphism for all $`๐ญ\mathrm{Proj}(A)`$. But this latter statement is equivalent to the fact that $`B/\epsilon (A)`$ has finite length, thus to the fact that $`B\epsilon (A):_BA_+^n`$ for some $`n`$, hence to $`B\stackrel{~}{\epsilon }(D(A))`$.
(g): Let $`๐ญ\mathrm{Spec}(A)`$ of $`0pt2`$. Then, the canonical module $`K_{A_๐ญ}`$ is of $`0pt2`$ (cf \[28, 3.1.1\]). So, there is a $`K_{A_๐ญ}`$-regular sequence $`x,yA_๐ญ`$. By the left-exactness of the functor $`\mathrm{Hom}_{A_๐ญ}(K_{A_๐ญ},)`$ it follows that $`x,y`$ is a regular sequence with respect to $`\mathrm{Hom}_{A_๐ญ}(A_๐ญ,A_๐ญ)=B_๐ญ`$, so that $`0pt_{A_๐ญ}B_๐ญ2`$. This shows that the $`A`$-module $`B`$ satisfies $`S_2`$. As $`B`$ is finite over $`A`$, it satisfies $`S_2`$ as a ring.
(h): Assume that $`K(A)`$ is a Cohen-Macaulay module. Consider the exact sequence of graded $`A`$-modules $`0A\stackrel{๐}{}BB/A0`$. By statement (d) we have $`(B/A)_๐ญ=0`$ as soon as $`๐ญ\mathrm{Spec}(A)`$ is of height $`1`$. So $`dimB/Ad1`$ and $`\epsilon `$ induces an isomorphism of graded $`S`$-modules $`H^{d+1}(A)H^{d+1}(B)`$. By (2.1) we get an isomorphism of graded $`S`$-modules $`K^{d+1}(B)K(A)`$. Therefore, the $`A`$-module $`K^{d+1}(B)`$ is Cohen-Macaulay. In view of \[28, 3.2.3\] we thus get $`H_{A_+}^i(B)=0`$ for $`i=2,\mathrm{},d`$. By statement g) we have $`H_A^i(B)=0`$ for $`i=0,1`$. So, the $`A`$-module $`B`$ is Cohen-Macaulay. As $`B`$ is finite over $`A`$, it becomes a Cohen-Macaulay ring. โ
We now apply the previous result in the case of varieties of almost minimal degree. We consider $`B`$ as a graded extension ring of $`A`$ by means of $`\epsilon :AB`$.
###### Theorem 5.3.
Assume that $`X_k^r`$ is of almost minimal degree and that $`t:=0ptAdimX=:d`$. Then
* $`B`$ is a finite graded birational integral extension domain of $`A`$ and $`CM`$.
* There are linearly independent linear forms $`y_{t1},\mathrm{},y_rS_1`$ and an isomorphism of graded $`S`$-modules $`B/A(S/(y_{t1},\mathrm{},y_r))(1)`$.
* The Hilbert polynomial of $`B`$ is given by
$$P_B(x)=(rd+2)\left(\genfrac{}{}{0pt}{}{x+d1}{d}\right)+\left(\genfrac{}{}{0pt}{}{x+d1}{d1}\right).$$
* If $`t=1`$, then $`B=D(A)`$.
###### Proof.
(a): According to statement (d) of Theorem 4.2 the $`A`$-module $`K(A)`$ is CM. So our claim follows form statements (a), (b) and (h) of Proposition 5.2.
(d): Let $`t=1`$. According to Proposition 3.4 and statement (LABEL:3.10) of Remark 3.3, $`X`$ is isomorphic to a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree and thus CM (cf Reminder 3.2, part B) ). So, by statement (f) of Proposition 5.2 we have $`D(A)=B`$.
(b): We proceed by induction on $`t`$. If $`t=1`$, statement (d) gives $`B/AH^1(A)`$ and so we may conclude by statement (c) of Proposition 3.1. Let $`t>1`$. We write $`C=B/A`$ and consider the exact sequence of graded $`S`$-modules
$$0A\stackrel{๐}{}BC0.$$
Let $`\mathrm{}S_1\backslash \{0\}`$ be generic. As $`t>1`$ and $`B`$ is CM (as an $`A`$-module) we have $`0pt_A(C)>0`$. Therefore $`\mathrm{}\text{NZD}(C)`$. We write $`A^{}=A/\mathrm{}A`$ and $`T=S/\mathrm{}S`$. Then $`A^{}`$ is a domain and $`Y:=\mathrm{Proj}(A^{})\mathrm{Proj}(T)=_k^{r1}`$ is a variety of almost minimal degree (cf Remark 2.3).
Let $`K(A^{}):=K_T^d(A^{}),B^{}:=\mathrm{Hom}_T(K(A^{}),K(A^{}))`$ and $`C^{}=B^{}/A^{}`$. By induction there are linearly independent linear forms $`\overline{z}_{t1},\mathrm{},\overline{z}_rT_1`$ and an isomorphism of graded $`T`$-modules $`B^{}/A^{}(T/(\overline{z}_{t1},\mathrm{},\overline{z}_r)T)(1)`$. We write $`\mathrm{}=z_{t2}`$. Then, there are linear forms $`z_{t1},\mathrm{},z_rS_1`$ such that $`z_{t2},z_{t1},\mathrm{},z_r`$ are linearly independent and such that there is an isomorphism of graded $`S`$-modules $`C^{}(S/(z_{t2},\mathrm{},z_r)S)(1)`$. As $`z_{t2}=\mathrm{}\text{NZD}(C)`$, it suffices to show that there is an isomorphism of graded $`S`$-modules $`C/\mathrm{}CC^{}`$ (cf Lemma 4.1). As $`\mathrm{}\text{NZD}(C)`$ there is an exact sequence of graded $`S`$-modules
$$0A^{}\stackrel{๐ผ}{}B/\mathrm{}BC/\mathrm{}C0$$
in which $`\alpha `$ is induced by $`\epsilon `$. It thus suffices to construct an isomorphism of graded $`A`$-modules $`\overline{\gamma }`$, which occurs in the commutative diagram
where $`\epsilon ^{}`$ is used to denote the natural map. As $`A^{}`$ is a domain and as the $`A`$-module $`B`$ is CM and torsion free of rank $`1`$ (by statement (a)), the $`A^{}`$-module $`B/\mathrm{}B`$ is torsion-free and again of rank $`1`$ (as $`\mathrm{}S_1`$ is generic). By statement (a) the $`A^{}`$-module $`B^{}`$ is also torsion-free of rank $`1`$. So, it suffices to find an epimorphism $`\gamma :BB^{}`$, which occurs in the commutative diagram
and hence such that $`\gamma (1_B)=1_B^{}`$.
By our choice of $`\mathrm{}`$ and in view of statements (a), (c) and (d) of Theorem 4.2 we have $`\mathrm{}\text{NZD}(K^d(A))\text{NZD}(K(A))`$. So, by (2.11) of Remark 2.3 we get an exact sequence of graded $`S`$-modules
$$0K(A)(1)\stackrel{}{}K(A)\stackrel{๐}{}K(A^{})(1)0.$$
Let $`U:=\mathrm{Ext}_S^1(K(A),K(A))(1)`$. If we apply the functors
$$\mathrm{Hom}_S(K(A),)\text{ and }\mathrm{Hom}_S(,K(A^{}))(1)$$
to the above exact sequence, we get the following diagram of graded $`S`$-modules with exact rows and columns
(5.3)
where $`\mu :=\mathrm{Hom}_S(id_{K(A)},\pi ),\nu :=\mathrm{Hom}_S(\pi ,id_{K(A^{})})`$. With $`\gamma :=\nu ^1\mu `$, it follows
$`\begin{array}{c}\hfill \gamma (1_B)=\gamma (id_{K(A)})=\nu ^1(\mu (id_{K(A)}))=\nu ^1(\pi id_{K(A)})=\nu ^1(\pi )\\ \hfill =\nu ^1(id_{K(A^{})(1)}\pi )=\nu ^1(\nu (id_{K(A^{})}))=id_{K(A^{})}=1_B^{}.\end{array}`$
So, it remains to show that $`\mu `$ is surjective. It suffices to show that $`(0:_U\mathrm{})=0`$. Assume to the contrary, that $`0:_U\mathrm{}0`$. Then $`\mathrm{}`$ belongs to some associated prime ideal $`๐ญ\mathrm{Ass}_SU.`$ As $`\mathrm{}`$ is generic, this means that $`S_1๐ญ`$ so that $`S_+๐ญ`$ and hence $`0:_US_+0`$, thus $`0:_{(0:_U\mathrm{})}S_+0`$. Therefore $`0pt_S(0:_U\mathrm{})=0`$. In view of statement (a) we have $`0pt_S(B/\mathrm{}B)=0pt_S(B^{})=d`$. Moreover the above diagram (5.3) yields an exact sequence of graded $`S`$-modules
$$0B/\mathrm{}BB^{}0:_U\mathrm{}0,$$
which shows that $`0pt_S(0:_U\mathrm{})d1>0`$, a contradiction.
(c): By statement (b) we have $`P_{B/A}(x)=\left(\genfrac{}{}{0pt}{}{x+t3}{t2}\right)`$ if $`t>1`$ and $`P_{B/A}(x)=0`$ if $`t=1`$. In view of statement (a) of Corollary 4.4 we get our claim. โ
Now, we are ready to draw a few conclusions about the geometric aspect.
###### Notation 5.4.
A) We convene that $`_k^1=\mathrm{}`$ and we use CM$`(X)`$, $`S_2(X)`$ and $`\mathrm{Nor}(X)`$ to denote respectively the locus of Cohen-Macaulay points, $`S_2`$-points and normal points of $`X`$.
B) If $`\nu :\stackrel{~}{X}X`$ is a morphism of schemes, we denote by $`\mathrm{Sing}(\nu )`$ the set
$$\{xX\nu _x^{\mathrm{}}:๐ช_{X,x}\stackrel{\simeq ฬธ}{}(\nu _{}๐ช_{\stackrel{~}{X}})_x\}$$
of all points $`xX`$ over which $`\nu `$ is singular.
###### Definition 5.5.
We say that $`xX`$ is a Goto or $`G`$-point, if the local ring $`๐ช_{X,x}`$ is of โGoto typeโ (cf ) thus if $`dim๐ช_{X,x}>1`$ and
$$H_{๐ช_{X,x}}^i(๐ช_{X,x})=\{\begin{array}{cc}0,\hfill & \text{if }i1,dim๐ช_{X,x},\hfill \\ \kappa (x),\hfill & \text{if }i=1.\hfill \end{array}$$
###### Theorem 5.6.
Assume that $`X_k^r`$ is of almost minimal degree and that $`t:=0ptAdimX=:d`$. Then
* $`B`$ is the homogeneous coordinate ring of a $`d`$-dimensional variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree.
* $`B`$ is the normalization of $`A`$.
* The normalization $`\nu :\stackrel{~}{X}X`$ given by the inclusion $`\epsilon :AB`$ is induced by a projection $`\varrho :_k^{r+1}\{p\}_k^r`$ from a point $`p_k^{r+1}\stackrel{~}{X}`$.
* The secant cone $`\mathrm{Sec}_p(\stackrel{~}{X})_k^{r+1}`$ is a projective subspace of dimension $`t1`$ and $`\mathrm{Sing}(\nu )=\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})X`$ is a projective subspace $`_k^{t2}_k^r`$.
* The generic point $`xX`$ of $`\text{ }\text{Sing}(\nu )`$ is a $`G`$-point.
* $`\mathrm{Nor}(X)=S_2(X)=\text{CM}(X)=X\text{Sing }(\nu )`$.
###### Proof.
(a): By Proposition 5.2 (a), (b) we see that $`B`$ is an integral, positively graded $`k`$-algebra of finite type and with $`dimB=dimA=d+1`$. By Theorem 5.3 (b) and on use of Nakayama we have in addition $`B=k[B_1]`$ with $`dim_kB_1=dim_kA_1+1=r+2`$. So, $`B`$ is the homogeneous coordinate ring of a non-degenerate projective variety $`\stackrel{~}{X}_k^{r+1}`$ of dimension $`d`$. By Theorem 5.3 (c) we have $`\mathrm{deg}\stackrel{~}{X}=r+1`$.
(b): By statement (a) the ring $`B`$ is normal (cf Reminder 3.2 C) ). In addition, $`B`$ is a birational integral extension ring of $`A`$.
(c): This follows immediately from the fact that $`dim_kB_1=dim_kA_1+1=r+2`$ (cf part C) of Remark 2.5).
(e): According to Theorem 5.3 (b) we have an exact sequence of graded $`S`$-modules
$$0AB(S/P)(1)0,$$
where $`P:=(y_{t1},\mathrm{},y_r)S`$ with appropriate independent linear forms $`y_{t1},\mathrm{},y_rS_1`$. In particular $`0:_SB/A=P`$ and hence $`IP`$. It follows that
$$\mathrm{Sing}(\nu )=\mathrm{Supp}\mathrm{Coker}(\nu :๐ช_X\nu _{}๐ช_{\stackrel{~}{X}})=\mathrm{Supp}((B/A)^{})=\mathrm{Proj}(S/P)$$
and that $`x:=P/I\mathrm{Proj}(A)=X`$ is the generic point of $`\mathrm{Sing}(\nu )`$. Localizing the above sequence at $`x`$ we get an exact sequence of $`๐ช_{X,x}`$-modules
(5.4)
$$0๐ช_{X,x}(\nu _{}๐ช_{\stackrel{~}{X}})_x\kappa (x)0$$
in which $`๐ช_{X,x}`$ has dimension $`dt+1>1`$. As $`B`$ is a CM-module over $`A`$ (cf Theorem 5.3 (a) ), $`(\nu _{}๐ช_{\stackrel{~}{X}})_x\stackrel{~}{B}_x`$ is a CM-module over the local domain $`๐ช_{X,x}`$. So, the sequence (5.4) shows that $`H_{๐ช_{X,x}}^1(๐ช_{X,x})\kappa (x)`$ and $`H_{๐ช_{X,x}}^i(๐ช_{X,x})=0`$ for all $`i1,dim(๐ช_{X,x})`$.
(d): Let $`PS`$ be as above. We already know that $`\mathrm{Sing}(\nu )=\mathrm{Proj}(S/P)=_k^{t2}_k^r`$. Moreover, a closed point $`qX`$ belongs to $`\mathrm{Sing}(\nu )`$ if and only if the line $`\varrho ^1(q)_k^{r+1}`$ is a secant line of $`\stackrel{~}{X}`$. So
$$\mathrm{Sec}_p(\stackrel{~}{X})=\{p\}\varrho ^1(\mathrm{Sing}(\nu ))=\{p\}\varrho ^1(_k^{t1})=_k^{t1}_k^{r+1}$$
and $`\mathrm{Sing}(\nu )=\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})`$.
(f): $`\nu :\stackrel{~}{X}\nu ^1(\mathrm{Sing}(\nu ))X\mathrm{Sing}(\nu )`$ is an isomorphism. As $`\stackrel{~}{X}`$ is a normal CM-variety (cf Reminder 3.2 C) ) it follows that $`\mathrm{Nor}(X),\text{CM}(X)X\mathrm{Sing}(\nu )`$. As $`S_2(X)\mathrm{Nor}(X)\text{CM}(X)`$ it remains to show that $`S_2(X)X\mathrm{Sing}(\nu )`$. As the generic point $`x`$ of $`\mathrm{Sing}(\nu )`$ is not an $`S_2`$-point (cf statement (e) ), our claim follows. โ
Now, we may extend Proposition 3.4 to arbitrary non arithmetically Cohen-Macaulay varieties.
###### Corollary 5.7.
Let $`1tdim(X)`$. Then, the following statements are equivalent:
* $`X`$ is of almost minimal degree and of arithmetic depth $`t`$.
* $`X`$ is the projection $`\varrho (\stackrel{~}{X})`$ of a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree from a point $`p_k^{r+1}\stackrel{~}{X}`$ such that $`dim\mathrm{Sec}_p(\stackrel{~}{X})=t1`$.
###### Proof.
(i) $``$ (ii): Clear by Theorem 5.6 (d).
(ii) $``$ (i): Let $`\stackrel{~}{X}_k^{r+1}`$, $`p`$ and $`\varrho :_k^{r+1}\stackrel{~}{X}X`$ be as in statement (ii). Let $`\stackrel{~}{A}`$ be the homogeneous coordinate ring of $`\stackrel{~}{X}`$. Then $`\stackrel{~}{A}`$ is normal (cf Reminder 3.2 C) ). Observe that $`\varrho :\stackrel{~}{X}X`$ is a finite morphism (cf Remark 3.3 A) ) so that $`dim\stackrel{~}{X}=d`$.
As $`dim\mathrm{Sec}_p(\stackrel{~}{X})=t1<d=dim\stackrel{~}{X}<dim\text{Join}(p,\stackrel{~}{X})`$ there are lines joining $`p`$ and $`\stackrel{~}{X}`$ which are not secant lines of $`\stackrel{~}{X}`$. So, in view of Remark 3.3 A) and statement (3.9) of Remark 3.3 B), the morphism $`\varrho :\stackrel{~}{X}X`$ is birational and $`X_k^r`$ is of almost minimal degree.
Moreover, the finite birational morphism $`\varrho :\stackrel{~}{X}X`$ is induced by a finite injective birational homomorphism $`\delta :A\stackrel{~}{A}`$ of graded rings. Thus, by Theorem 5.6 (b) we get an isomorphism of graded rings $`\iota ,`$ which occurs in the commutative diagram
Now Theorem 5.6 d) shows that $`0ptA=dim\mathrm{Sec}_p(\stackrel{~}{X})+1=t`$. โ
As a further application of Theorem 5.6 we now have a glance at arithmetically Cohen-Macaulay varieties of almost minimal degree and show that their non-normal locus is either empty or a linear space. More precisely
###### Proposition 5.8.
Assume that $`X_k^r`$ is of almost minimal degree, $`S_2`$ and not normal. Let $`dimX=:d`$. Then $`X`$ is arithmetically Cohen-Macaulay and the non-normal locus $`X\mathrm{Nor}(X)`$ is a linear space $`_k^{d1}_k^r`$.
###### Proof.
(Induction on $`d`$) Let $`d=1`$. Then $`X`$ is a curve of degree $`r+1`$ and thus may have at most one singular point (cf \[4, (4.7) (B)\]). So, let $`d>1`$. Then $`\mathrm{Nor}(X)X=S_2(X)`$ shows that $`X`$ is arithmetically Cohen-Macaulay (cf Theorem 5.6 f) ). Moreover, as $`X`$ is $`S_2`$ and not normal, the Serre criterion for normal points shows that the non-normal locus $`X\mathrm{Nor}(X)`$ of $`X`$ is of pure codimension $`1`$.
Let $`ZX`$ be the reduced and purely $`(d1)`$-dimensional closed subscheme supported by $`X\mathrm{Nor}(X)`$. It suffices to show that $`\mathrm{deg}Z=1`$. Now, let $`\mathrm{}S_1`$ be a generic linear form and consider the hyperplane $`_k^{r1}:=\mathrm{Proj}(S/\mathrm{}S)`$. Then, the hyperplane section $`X^{}:=_k^{r1}X=\mathrm{Proj}(A/\mathrm{}A)`$ is an arithmetically Cohen-Macaulay-variety of almost minimal degree in $`_k^{r1}`$ with $`dimX^{}=d1`$, and $`Z^{}:=_k^{r1}ZX^{}`$ is a reduced purely $`1`$-codimensional subscheme with $`\mathrm{deg}Z^{}=\mathrm{deg}Z`$. Now, let $`z^{}`$ be one of the generic points of $`Z^{}`$. Then $`z^{}Z`$ shows that $`๐ช_{X,z^{}}`$ is not normal and hence not regular. As $`๐ช_{X^{},z^{}}`$ is a hypersurface ring of $`๐ช_{X,z}`$, the ring $`๐ช_{X^{},x^{}}`$ is not regular either. As $`dim๐ช_{X^{},z^{}}=1`$ it follows that $`๐ช_{X^{},z^{}}`$ is not normal and hence $`z^{}X^{}\mathrm{Nor}(X^{})`$. By induction we have $`X^{}\mathrm{Nor}(X^{})=_k^{d2}`$ for some linear subspace $`_k^{d2}_k^r`$. It follows that $`\{\overline{z^{}}\}=_k^{d2}`$. This shows that the closed reduced subschemes $`Z^{}`$ and $`_k^{d2}`$ of $`_k^r`$ coincide, hence $`\mathrm{deg}Z=\mathrm{deg}Z^{}=1`$. โ
## 6. Del Pezzo Varieties and Fujitaโs Classification
In this section we shall treat projective varieties of almost minimal degree which are arithmetically Cohen-Macaulay. We call these varieties maximal Del Pezzo varieties and make sure that this is in coincidence with Fujitaโs notion of Del Pezzo variety . We also briefly discuss the link with Fujitaโs classification of varieties of $`\mathrm{\Delta }`$-genus $`1`$.
###### Remark 6.1.
A) Let $`d:=dim(X)>0`$ and let $`\omega _X=K(A)^{}`$ denote the dualizing sheaf of $`X`$. Keep in mind that a finitely generated graded $`A`$-module of $`0pt>1`$ is determined (up to a graded isomorphism) by the sheaf of $`๐ช_X`$-modules $`\stackrel{~}{M}`$ induced by $`M`$. So, as $`K(A)`$ satisfies the second Serre property $`S_2`$ (cf \[28, 3.1.1\]), we have for each $`r`$:
(6.1) $`\omega _X๐ช_X(r)\text{ if and only if }K(A)D(A)(r).`$
(6.2) $`\text{If depth }A>1,\text{ then }\omega _X๐ช_X(r)\text{ if and only if }K(A)A(r).`$
B) $`X_k^r`$ is said to be linearly complete if the inclusion morphism $`X_k^r`$ is induced by the complete linear system $`|๐ช_X(1)|`$. It is equivalent to say that the natural monomorphism $`\eta :A_1H^0(X,๐ช_X(1))=D(A)_1`$ is an isomorphism hence โ equivalently โ that $`H^1(A)_1=0`$.
###### Theorem 6.2.
The following statements are equivalent:
* $`X`$ is arithmetically Gorenstein and of almost minimal degree.
* $`X`$ is arithmetically Cohen-Macaulay and of almost minimal degree.
* $`X`$ is $`S_2`$, linearly complete and of almost minimal degree.
* $`\omega _X๐ช_X(1d)`$ and $`X`$ is of almost minimal degree.
* $`K(A)A(1d)`$ and $`X`$ is arithmetically Cohen-Macaulay.
* $`\omega _X๐ช_X(1d)`$ and $`X`$ is arithmetically Cohen-Macaulay.
* $`H^{d+1}(A)_{1d}k`$ and $`H^1(A)_1=H^i(A)_n=0`$ for $`2id`$ and $`1dn1`$.
###### Proof.
The implications (i) $``$ (ii) $``$ (iii) and (v) $``$ (vi) are obvious. The implication (ii) $``$ (i) follows by Remark 4.5, the implication (vi) $``$ (v) by (6.2). It remains to prove the implications (v) $``$ (vii) $``$ (ii) and (ii) $``$ (v) $``$ (iv) $``$ (iii) $``$ (ii). The implication (v) $``$ (vii) is easy.
(vii) $``$ (ii): $`H^{d+1}(A)_{1d}k`$ implies that $`H^{d+1}(A)_n=0`$ for all $`n>1d`$. As $`P_A(n)=dim_kA_k_{i=1}^{d+1}(1)^idim_kH^i(A)_n`$, statement (vii) implies that
$$P_A(n)=\{\begin{array}{cc}(1)^d,\hfill & \text{if }n=1d,\hfill \\ 0,\hfill & \text{if }1d<n<0,\hfill \\ 1,\hfill & \text{if }n=0,\hfill \\ r+1,\hfill & \text{if }n=1.\hfill \end{array}$$
So, we may write $`P_A(n)=(rd+2)\left(\genfrac{}{}{0pt}{}{n+d1}{d}\right)+\left(\genfrac{}{}{0pt}{}{n+d2}{d2}\right)`$. In particular, $`X`$ is of almost minimal degree. But then, by Corollary 4.4 b) we see that $`X`$ must be arithmetically Cohen-Macaulay.
(ii) $``$ (v): Assume that $`X`$ is arithmetically Cohen-Macaulay and of almost minimal degree. Then, the shape of the Hilbert series given in Corollary 4.4 (a) allows to conclude that $`K(A)A(1d)`$.
(v) $``$ (iv): We know that (v) $``$ (vii) and (vii) $``$ (ii). So, (v) implies that $`X`$ is of almost minimal degree and hence induces (iv).
(iv) $``$ (iii): By (6.2), statement (iv) implies $`K(A)D(A)(1d)`$ and hence by Theorem 4.2 (b) that $`X`$ is arithmetically Cohen-Macaulay.
(iii) $``$ (ii): Assume that $`X`$ is of almost minimal degree, $`S_2`$ and linearly complete. We have to show that $`X`$ is arithmetically Cohen-Macaulay. Assume to the contrary that $`t=0ptAd`$. As $`H^1(A)_1=0`$, Proposition 3.1 (c) yields $`t>1`$. But now, by Theorem 5.6 (d) and (e), $`X`$ contains a $`G`$-point and thus cannot be $`S_2`$ โ a contradiction. โ
###### Definition 6.3.
A) $`X_k^r`$ is called a maximal Del Pezzo variety if it satisfies the equivalent conditions (i) โ (vii) of Theorem 6.2.
B) $`X_k^r`$ is called a Del Pezzo variety, if there is an integer $`r^{}r`$, a maximal Del Pezzo variety $`X^{}_k^r^{}`$ and a linear projection $`\pi :_k^r^{}_k^{r^{}r1}_k^r`$ with $`X^{}^{r^{}r1}=\mathrm{}`$ and such that $`\pi `$ gives rise to an isomorphism $`\pi :X^{}\stackrel{}{}X`$. So, $`X_k^r`$ is Del Pezzo if and only if it is obtained by a non-singular projection of a maximal Del Pezzo variety.
###### Remark 6.4.
A) As a linearly complete variety $`X_k^r`$ cannot be obtained by a proper non-singular projection of a non-degenerate variety $`X^{}_k^r^{}`$ we can say that $`X_k^r`$ is maximally Del Pezzo if and only if it is Del Pezzo and linearly complete.
B) Keep the previous notation and let $`r^{}=h^0(X,๐ช_X(1))1`$, whence
$$r^{}=dim_kD(A)_11=r+dim_kH^1(A)_k.$$
Moreover let $`\phi :X_k^r^{}`$ be the closed immersion defined by the complete linear system $`|๐ช_X(1)|`$ and set $`X^{}:=\phi (X)`$. Then $`X^{}_k^r^{}`$ is linearly complete with homogeneous coordinate ring $`A^{}=k[D(A)_1]D(A)`$, whereas the isomorphism $`\phi :X\stackrel{}{}X^{}`$ is induced by the inclusion $`AA^{}`$ and inverse to an isomorphism $`\pi :X^{}\stackrel{}{}X`$ which is the restriction of a linear projection $`\pi :_k^r^{}_k^{r^{}r1}_k^r`$ with $`X^{}_k^{r^{}r1}=\mathrm{}`$. Now, clearly $`X^{}_k^r^{}`$ is linearly complete and moreover
$`D(A^{})`$ $`=D(A),H^i(A^{})H^i(A),i1,K(A^{})K(A)\text{ and }P_A^{}(x)=P_A(x).`$
Note that $`X^{}_k^r^{}`$ is called the linear completion of $`X_k^r`$.
C) Observe that the linear completion of $`X_k^r`$ is the maximal non-degenerate projective variety $`X^{}_k^r^{}`$ which can be projected non-singularly onto $`X.`$ More precisely: If $`\stackrel{~}{X}_k^{\overline{r}}`$ is a non-degenerate projective variety and $`\overline{\pi }:\stackrel{~}{X}\stackrel{}{}X`$ is an isomorphism induced by a linear projection $`\overline{\pi }:_k^{\overline{r}}_k^r`$, then $`\overline{r}r^{}`$ and the isomorphism $`\overline{\pi }^1\pi :X^{}\stackrel{}{}\stackrel{~}{X}`$ comes from a linear projection $`\varrho :_k^r^{}_k^{\overline{r}}`$. In particular, if $`\stackrel{~}{X}_k^{\overline{r}}`$ is linearly complete we have $`r^{}=\overline{r}`$ and $`\varrho `$ becomes an isomorphism so that we may identify $`X^{}`$ with $`\stackrel{~}{X}`$. Consequently, by what we said in part A) it follows that $`X_k^r`$ is Del Pezzo if and only if its linear completion $`X^{}_k^r^{}`$ is (maximally) Del Pezzo.
We now shall tie the link to Fujitaโs classification of polarized Del Pezzo varieties.
###### Remark 6.5.
(see ). A) A polarized variety over $`k`$ is a pair $`(V,)`$ consisting of a reduced irreducible projective variety $`V`$ over $`k`$ and an ample invertible sheaf of $`๐ช_V`$-modules $``$.
B) Let $`(V,)`$ be a polarized $`k`$-variety. For a coherent sheaf of $`๐ช_V`$-modules $``$ and $`i_0`$ let $`h^i(V,)`$ denote the $`k`$-dimension of the $`i`$-th Serre cohomology group $`H^i(V,)`$ of $`V`$ with coefficients in $``$. Then, the function
$$n\chi _{(V,)}(n):=\underset{i=0}{\overset{dimV}{}}(1)^ih^i(V,^n)$$
is a polynomial of degree $`dimV`$, the so called Hilbert polynomial of the polarized variety $`(V,)`$.
C) Let $`(V,)`$ be a polarized variety of dimension $`d`$. Then, there are uniquely determined integers $`\chi _i(V,),i=0,\mathrm{},d`$ such that
$$\chi _{(V,)}(n)=\underset{i=0}{\overset{d}{}}\chi _i(V,)\left(\genfrac{}{}{0pt}{}{n+i1}{i}\right).$$
Clearly $`\chi _d(V,)>0`$. The degree, the $`\mathrm{\Delta }`$-genus and the sectional genus of the polarized variety $`(V,)`$ are defined respectively by
$`\mathrm{deg}(V,):=`$ $`\chi _d(V,);`$
$`\mathrm{\Delta }(V,):=`$ $`d+\mathrm{deg}(V,)h^0(V,);`$
$`g_s(V,):=`$ $`1\chi _{d1}(V,).`$
D) According to Fujita (cf ) the polarized variety $`(V,)`$ is called a Del Pezzo variety, if it satisfies the following conditions
(6.3) $`\mathrm{\Delta }(V,)=1,`$
(6.4) $`g_s(V,)=1,`$
(6.5) $`V\text{ has only Gorenstein singularities and }\omega _V^{(1dimV)},`$
(6.6) $`\text{For all }i0,dimV\text{ and all }n\text{ it holds }H^i(V,^n)=0.`$
###### Remark 6.6.
A) We consider the polarized variety $`(X,๐ช_X(1))`$. For all $`n`$ we have $`H^0(X,๐ช_X(1)^n)=H^0(X,๐ช_X(n))=D(A)_n`$. Thus:
(6.7)
$$\chi _{(X,๐ช_X(1))}(n)=P_A(n),$$
where $`P_A`$ is the Hilbert polynomial of $`A.`$ Therefore
(6.8) $`\mathrm{deg}(X,๐ช_X(1))`$ $`=\mathrm{deg}X,`$
(6.9) $`\mathrm{\Delta }(X,๐ช_X(1))`$ $`=\mathrm{deg}X\mathrm{codim}X1dim_kH^1(A)_1.`$
As a consequence of the last equality we obtain
(6.10)
$$\mathrm{\Delta }(X,๐ช_X(1))\mathrm{deg}X\mathrm{codim}X+1,$$
with equality if and only if $`X_k^r`$ is linearly complete.
B) Let $`X^{}_k^r^{}`$ be the linear completion of $`X_k^r`$. Then $`(X,๐ช_X(1))`$ and $`(X^{},๐ช_X^{}(1))`$ are isomorphic polarized varieties. In particular $`(X^{},๐ช_X^{}(1))`$ is Del Pezzo in the sense of Fujita if and only $`(X,๐ช_X(1))`$ is.
###### Lemma 6.7.
Let $`X_k^r`$ be of almost minimal degree. Then
* $$\mathrm{\Delta }(X,๐ช_X(1))=\{\begin{array}{cc}0,\hfill & \text{ if }\text{ }\text{depth }A=1,\hfill \\ 1,\hfill & \text{ if }\text{ }\text{depth }A>1.\hfill \end{array}$$
* $$g_s(X,๐ช_X(1))=\{\begin{array}{cc}0,\hfill & \text{ if }X\text{ is not arithmetically Cohen-Mcaulay},\hfill \\ 1,\hfill & \text{ if }X\text{ is arithmetically Cohen-Macaulay}.\hfill \end{array}$$
###### Proof.
(a): This follows immediately from Proposition 3.1 c), Theorem 4.2 b) and by (6.9).
(b): This is a consequence of Corollary 4.4 b) and (6.7). โ
###### Theorem 6.8.
Let $`X_k^r`$ be of dimension $`d>0`$. Then, the following statements are equivalent:
* $`X`$ is Del Pezzo in the sense of Definition 6.3 B).
* $`(X,๐ช_X(1))`$ is Del Pezzo in the sense of Fujita.
* $`\mathrm{\Delta }(X,๐ช_X(1))=g_s(X,๐ช_X(1))=1`$.
* $`\mathrm{\Delta }(X,๐ช_X(1))=1`$ and $`H^i(X,๐ช_X(n))=0`$ for all $`i0,d`$ and all $`n`$.
* $`g_s(X,๐ช_X(1))=1`$ and $`H^i(X,๐ช_X(n))=0`$ for all $`i0,d`$ and all $`n`$.
* $`\mathrm{\Delta }(X,๐ช_X(1))=1`$ and $`\omega _X๐ช_X(1d)`$.
* $`H^i(X,๐ช_X(n))=0`$ for all $`i0,d`$ and all $`n`$, and $`\omega _X๐ช_X(1d)`$.
* $`\mathrm{\Delta }(X,๐ช_X(1))=1`$ and $`X`$ is $`S_2`$.
###### Proof.
First let us fix a few notation. Let $`r^{}=h^0(X,๐ช_X(1))=dim_kD(A)_1`$ and let $`X^{}_k^r^{}`$ be the linear completion of $`X_k^r`$.
(i) $``$ (ii): Let $`X`$ be Del Pezzo in the sense of Definition 6.3 B). According to Remark 6.4 C) we get that $`X^{}_k^r^{}`$ is maximally Del Pezzo. According to Theorem 6.2, the pair $`(X^{},๐ช_X^{}(1))`$ thus satisfies the requirements (6.3) - (6.6) of Remark 6.5 D) and hence is Del Pezzo in the sense of Fujita. By Remark 6.6 B) the same follows for $`(X,๐ช_X(1))`$.
Clearly statement (ii) implies each of the statements (iii) - (viii). So, it remains to show that each of the statements (iii) - (viii) implies statement (i). According to Remark 6.4 C) we may replace $`X`$ by $`X^{}`$ in statement (i). As $`(X^{},๐ช_X^{}(1))`$ and $`(X,๐ช_X(1))`$ are isomorphic polarized varieties we may replace $`X`$ by $`X^{}`$ in each of the statements (iii) - (viii). So, we may assume that $`X_k^r`$ is linearly complete.
(iii) $``$ (i): According to statement (6.10) of Remark 6.6 A) the equality $`\mathrm{\Delta }(X,๐ช_X(1))=1`$ implies that $`X_k^r`$ is of almost minimal degree. But now by Lemma 6.7 b) the equality $`g_s(X,๐ช_X(1))=1`$ implies that $`X`$ is arithmetically Cohen-Macaulay.
(iv) $``$ (i): By $`\mathrm{\Delta }(X,๐ช_X(1))=1`$ we see again that $`X_k^r`$ is of almost minimal degree. As $`H^1(A)_1=0`$ we have $`\text{depth }A>1`$ (Proposition 3.1 C) ). As $`H^{i+1}(A)_nH^i(X,๐ช_X(n))=0`$ for all $`i0,d`$, it follows that $`\text{depth }A=d+1`$, so that $`X`$ is arithmetically Cohen-Macaulay.
(vi) $``$ (i): As we have proved the implication (iii) $``$ (i) it suffices to prove that statement (v) implies the equality $`\mathrm{\Delta }(X,๐ช_X(1))=1`$. We proceed by induction on $`d`$. Let $`d=1`$. As $`g_s(X,๐ช_X(1))=1`$ implies $`\chi _0(X,๐ช_X(1))=0`$ we get $`\chi _{(X,๐ช_X(1))}(n)=(\mathrm{deg}X)n`$. As $`H^1(X,๐ช_X(1))=0`$ it follows $`r+1=h^0(X,๐ช_X(1))=\chi _{(X,๐ช_X(1))}(1)=\mathrm{deg}X`$ and hence by (6.10) we get $`\mathrm{\Delta }(X,๐ช_X(1))=1`$.
So, let $`d>1`$, let $`\mathrm{}S_1`$ be a generic linear form and consider the irreducible projective variety $`Y:=\mathrm{Proj}(A/\mathrm{}A)_k^r:=\mathrm{Proj}(S/\mathrm{}S)`$ of dimension $`d1`$ and with homogeneous coordinate ring $`A^{}=(A/\mathrm{}A)/H^0(A/\mathrm{}A)`$. As
$$\chi _{(Y,๐ช_Y(1))}(n)=P_A^{}(n)=\mathrm{\Delta }P_A(x)=\mathrm{\Delta }\chi _{(X,๐ช_X(1))}(n)$$
it follows $`\mathrm{\Delta }(Y,๐ช_Y(1))=\mathrm{\Delta }(X,๐ช_X(1))`$ and $`g_s(Y,๐ช_Y(1))=g_s(X,๐ช_X(1))=1`$. So, by induction it suffices to show that
$$H^i(Y,๐ช_Y(n))=H^{i+1}(A^{})_n=0\text{ for all }i0,d1\text{ and all }n,$$
and that $`Y_k^{r1}`$ is linearly complete, hence that $`H^1(A^{})_1=0`$. As
$$H^{i+1}(A)_n=H^i(X,๐ช_X(n))=0$$
for all $`i0,d`$ and all $`n`$ and as $`H^1(A)_1=0`$ this follows immediately if we apply cohomology to the sequence
$$0A(1)\stackrel{}{}AA/\mathrm{}A0$$
and observe that $`H^j(A^{})H^j(A/\mathrm{}A)`$ for all $`j>0`$.
(vi) $``$ (i): This is immediate by (6.10) and Theorem 6.2.
(vii) $``$ (i): As $`X`$ is linearly complete, $`H^1(A)_1=0`$. Moreover by our hypothesis $`H^i(A)=0`$ for all $`i1,d+1`$. Finally, by (6.2) we have $`K(A)D(A)(1d)`$, hence $`H^{d+1}(A)_{1d}K(A)_{d1}D(A)_0k`$. So, statement (vii) of Theorem 6.2 is true.
(viii) $``$ (i): In view of (6.10), statement (viii) implies statement (iii) of Theorem 6.2. โ
Our next aim is to extend Theorem 5.3 to maximal Del Pezzo varieties.
###### Theorem 6.9.
Let $`X_k^r`$ be a maximal Del Pezzo variety of dimension $`d`$ which is non-normal. Let $`B=kB_1B_2\mathrm{}`$ be the graded normalization of $`A`$. Then:
* There are linearly independent linear forms $`y_d,y_{d+1},\mathrm{},y_rS_1`$ such that $`B/A(S/(y_d,y_{d+1},\mathrm{},y_r))(1)`$.
* $`B`$ is the homogeneous coordinate ring of a variety of minimal degree $`\stackrel{~}{X}_k^{r+1}`$. In particular, $`B`$ is a Cohen-Macaulay ring.
###### Proof.
We make induction on $`d`$. The case $`d=1`$ is clear by Proposition 5.2 C). Therefore, let $`d>1`$. Statement (b) follows easily from statement (a). So, we only shall prove this latter. According to Proposition 5.8 there are linearly independent linear forms $`y_d,y_{d+1},\mathrm{},y_rS_1`$ such that $`I(y_d,y_{d+1},\mathrm{},y_r)`$ and such that $`๐ฐ:=(y_d,y_{d+1},\mathrm{},y_r)/IS/I=A`$ defines the non-normal locus $`X\mathrm{Nor}(X)`$ of $`X`$. Observe that $`๐ฐA`$ is a prime of height $`1`$ and that $`A/๐ฐS/(y_d,y_{d+1},\mathrm{},y_r)S`$ is a polynomial ring in $`d`$ inderminates over $`k`$. Next, we consider the canonical exact sequence
(6.11)
$$0ABC0$$
in which $`C:=B/A`$ is a finitely generated graded $`A`$-module such that $`C_0=0`$ and $`\mathrm{Rad}\mathrm{Ann}_BC=๐ฐ`$. Our aim is to show that $`C(A/๐ฐ)(1)`$. Let $`\mathrm{}S_1`$ be a generic linear form.
Then, according to Bertini and as $`0ptA>1`$, the ideal $`\mathrm{}AA`$ is prime. Moreover $`(A\mathrm{}A)๐ฐ\mathrm{}`$, so that $`(A\mathrm{}A)^1B=A_\mathrm{}A`$. It follows that $`\mathrm{}B`$ has a unique minimal prime $`๐ญ`$ and that $`๐ญB_๐ญ=\mathrm{}B_๐ญ`$. As $`B`$ is $`S_2`$ we get $`\mathrm{}B=๐ญ`$, so that $`\mathrm{}B`$ is a prime ideal of $`B`$. Therefore $`B/\mathrm{}B`$ is a finite birational integral extension domain of $`A^{}:=A/\mathrm{}A`$ and hence a subring of the graded normalization $`B^{}`$ of $`A^{}`$.
As $`0ptA>2`$ and $`0ptB2,`$ the short exact sequence (6.11) yields $`0ptC2`$. In particular $`\mathrm{}`$ is $`C`$-regular. Hence we get the following commutative diagram with exact rows and columns in which $`U=\mathrm{Coker}\iota `$ is a graded $`A`$-module
(6.12) $`\begin{array}{c}\hfill \text{}\end{array}`$
Now $`X^{}:=\mathrm{Proj}(A^{})\mathrm{Proj}(S/\mathrm{}S)=_k^{r1}`$ is again a maximal Del Pezzo variety. Moreover $`X^{}`$ is non-normal, since otherwise $`B/\mathrm{}B=A/\mathrm{}A,`$ hence $`B=A`$. Let $`๐ฐ^{}A^{}`$ be the prime of height $`1`$ which defines the non-normal locus of $`X^{}`$ and keep in mind that $`A^{}/๐ฐ^{}`$ is a polynomial ring in $`d1`$ inderminates over $`k`$. By induction $`C^{}(A^{}/๐ฐ^{})(1).`$
Our next aim is to show that $`dimU0`$. As $`C0`$, we have $`C/\mathrm{}C0`$. As $`C^{}`$ is a free $`A^{}/๐ฐ^{}`$-module of rank one, it follows $`dimU<dim(A^{}/๐ฐ^{})=d1`$. As $`dim(B/\mathrm{}B)=dim(B^{})=d`$ it follows that $`\lambda `$ is an isomorphism in codimension one. As $`B^{}`$ is normal and hence satisfies the property $`R_1`$, it follows that $`B/\mathrm{}B`$ satisfies $`R_1`$, too.
Let $`s๐ฐ\{0\}.`$ As $`B`$ is normal, it satisfies the Serre property $`S_2`$ so that $`B/sB`$ satisfies $`S_1.`$ Therefore the set of generic points of the (closed) non-$`S_2`$-locus of the $`B`$-module $`B/sB`$ is given by
$$๐ฌ:=\{๐ฎV(sB):0pt(B_๐ฎ/sB_๐ฎ)=1<0pt๐ฎ1\}.$$
In particular $`๐ฌ`$ is finite, and hence $`\mathrm{}`$ avoids all members of $`๐ซ:=(\mathrm{Ass}_B(B/sB)๐ฌ)\mathrm{Proj}B.`$
Now, let $`๐ฏ\mathrm{Proj}BV(\mathrm{}B)`$ such that $`0pt๐ฏ>2.`$ If $`s๐ฏ,`$ the equality $`A_s=B_s`$ yields that $`B_๐ฏ`$ is a Cohen-Macaulay ring, so that $`0pt(B_๐ฏ/\mathrm{}B_๐ฏ)>1.`$ If $`s๐ฏ`$ the fact that $`\mathrm{}`$ avoids all members of $`๐ซ`$ implies that $`s,\mathrm{}`$ is a $`B_๐ฏ`$-sequence and $`0ptB_๐ฏ/sB_๐ฏ)>1.`$ It follows again that $`0pt(B_๐ฏ/\mathrm{}B_๐ฏ)>1.`$ This proves, that the scheme $`\mathrm{Proj}(B/\mathrm{}B)`$ is $`S_2.`$ As $`B/\mathrm{}B`$ satisfies $`R_1`$ it follows that $`\mathrm{Proj}(B/\mathrm{}B)`$ is a normal scheme, hence that $`\mathrm{Proj}(B/\mathrm{}B)=\mathrm{Proj}B^{}.`$ As a consequence, we get indeed that $`dimU0,`$ that is $`U`$ is a graded $`A`$-module of finite length.
Now, let $`๐ฑA`$ be the preimage of $`๐ฐ^{}`$ under the canonical map $`AA^{}`$. Then $`๐ฑ`$ and $`๐ฐ+\mathrm{}A`$ are primes of height 2 in $`A`$ and so
$$๐ฐ+\mathrm{}A=\mathrm{Rad}((\mathrm{Ann}_AC)+\mathrm{}A)=\mathrm{Rad}(\mathrm{Ann}_AC/\mathrm{}C)\mathrm{Ann}_AC^{}=\mathrm{Ann}_A(A^{}/๐ฐ^{}(1))=๐ฑ$$
implies that $`๐ฐ+\mathrm{}A=\mathrm{Ann}_A(C/\mathrm{}C)=๐ฑ.`$ As a consequence we get $`๐ฐC\mathrm{}C`$ and hence, by the genericity of $`\mathrm{}`$, that $`๐ฐC=0`$. It follows $`๐ฐB๐ฐ`$ and $`๐ฐ`$ becomes an ideal of $`B`$. Now, let $`aA`$ and $`cC\{0\}`$ such that $`ac=0`$. By the genericity of $`\mathrm{}`$ we may assume that $`c\mathrm{}C`$ so that $`\iota (c+\mathrm{}C)0`$ and $`a\iota (c+\mathrm{}C)=0.`$ It follows $`a๐ฑ=๐ฐ+\mathrm{}A`$ and hence, by genericity, that $`a๐ฐ.`$ This shows that $`C`$ is a torsion-free $`A/๐ฐ`$-module and hence that $`B/๐ฐ`$ is a torsion-free $`A/๐ฐ`$-module.
As $`\mathrm{rank}_{A/๐ฐ}C=e_0(C)=e_0(C/\mathrm{}C)=e_0(C^{})=\mathrm{rank}_{A^{}/๐ฐ^{}}C^{}=1`$ we get an exact sequence of graded $`A/๐ฐ`$-modules
(6.13)
$$0C(A/๐ฐ)(m)W0$$
with $`m`$ and $`dimW<dimA/๐ฐ=d.`$ We choose $`m`$ maximally. Then, there is no homogeneous element $`fA/๐ฐ`$ of positive degree with $`C(m)f(A/๐ฐ)`$, so that $`dimW=dim(W(m))=dim((A/๐ฐ)/C(m))<d1.`$
As $`0ptC2`$ we have $`0ptW1.`$ As $`\mathrm{}`$ is generic we thus get an exact sequence
$$0C/\mathrm{}C(A^{}/๐ฐ^{})(m)W/\mathrm{}W0$$
with $`dimW/\mathrm{}W=dimW1<d2=dimA^{}/๐ฐ^{}1.`$ Comparing Hilbert coefficients we get $`e_1(C/\mathrm{}C)=e_1((A^{}/๐ฐ^{})(m))=m.`$
The diagram (6.12) contains the short exact sequence $`0C/\mathrm{}C(A^{}/๐ฐ^{})(1)U0`$ with $`dimU0.`$ Assume first, that $`d3`$ so that $`dimU<dimA^{}/๐ฐ^{}1.`$ Then, we may again compare Hilbert coefficients and get $`e_1(C/\mathrm{}C)=e_1((A^{}/๐ฐ^{})(1))=1`$, hence $`m=1.`$ It follows $`dimW/\mathrm{}W=dimU0,`$ thus $`dimW1.`$ Suppose that $`dimW=1.`$ Then $`H^1(W)_n0`$ for all $`n0`$ and the sequence (6.11) yields $`H^2(C)_n0`$ for all $`n0.`$ As $`H^3(A)=0`$, the sequence (6.11) induces that $`H^2(B)_n0`$ for all $`n0;`$ but this contradicts the fact that $`B`$ is $`S_2`$. So, we have $`dimW0.`$ Now, by the sequence (6.13) we get $`W=H^0(W)H^1(C),`$ and $`0ptC>1`$ implies $`W=0.`$ Therefore $`CA/๐ฐ(1)`$.
It remains to treat the case $`d=2`$. Now $`A^{}/๐ฐ^{}k[x]`$ so that the $`A^{}/๐ฐ^{}`$-submodule $`C/\mathrm{}C`$ of $`(A^{}/๐ฐ^{})(1)`$ is generated by a single homogeneous element of degree $`m1.`$ By Nakayama it follows $`C(A/๐ฐ)(m).`$ It remains to show that $`m=1.`$
By the case $`d=1`$ we know that $`B^{}`$ is the homogeneous coordinate ring of a rational normal curve, so that $`H^2(B^{})_0=0.`$ The middle column of the diagram (6.12) now implies $`H^2(B/\mathrm{}B)_0=0`$. Applying cohomology to the exact sequence $`0B(1)\stackrel{}{}BB/\mathrm{}B0`$ we thus get an isomorphism $`H^3(B)_1H^3(B)_0`$ so that $`H^3(B)_1=0.`$ If we apply cohomology to the sequence (6.11) we thus get an exact sequence
$$0H^2(B)_1H^2(A/๐ฐ)_{1m}H^3(A)_10.$$
By Theorem 6.2 (cf statement (vii)) we have $`H^3(A)_1k`$. As $`A/๐ฐ`$ is a polynomial ring in two indeterminates over $`k`$ we have $`H^2(A/๐ฐ)_{1m}k^m.`$ It follows $`H^2(B)_1k^{m1}.`$
Moreover $`\stackrel{~}{X}:=\mathrm{Proj}B`$ is a projective normal surface and the natural morphism $`\nu :\stackrel{~}{X}X`$ is a normalization of $`X.`$ In particular $`:=B(1)^{}=\nu ^{}๐ช_X(1)`$ is an ample invertible sheaf of $`๐ช_{\stackrel{~}{X}}`$-modules and $`^n=\nu ^{}๐ช_X(n)=B(n)^{}`$ for all $`n.`$ In addition we have $`H^2(B)_1H^1(\stackrel{~}{X},^1)`$ and $`B_1H^0(\stackrel{~}{X},).`$ Moreover $`Y:=\mathrm{Proj}B/\mathrm{}B\mathrm{Proj}B^{}`$ is the effective divisor on $`\stackrel{~}{X}`$ defined by the global section $`\mathrm{}H^0(\stackrel{~}{X},)\backslash \{0\}.`$ As $`H^1(Y,๐ช_Y)=H^2(B/\mathrm{}B)=0,`$ the sectional genus $`g_s(\stackrel{~}{X},)`$ vanishes (cf \[1, (5.3) B)\]). As $`dim_kH^0(\stackrel{~}{X},)=dim_kB_1r+1>1`$ it follows $`H^1(\stackrel{~}{X},^1)=0`$ (cf \[1, Proposition (5.4)\]). Therefore $`k^{m1}H^2(B)_1=0`$, hence $`m=1.`$
Now, we may extend Theorem 5.6 as follows
###### Corollary 6.10.
Let $`X_k^r`$ be of almost minimal degree. Assume that either $`t:=0ptAdimX=:d`$ or that $`X`$ is maximally Del Pezzo (that is $`t=d+1)`$ and non-normal. Then, there is a $`d`$-dimensional variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree, a point $`p_k^{r+1}\stackrel{~}{X}`$ and a projection $`\varrho :_k^{r+1}\{p\}_k^r`$ from $`p`$ such that:
* $`\varrho (\stackrel{~}{X})=X`$ and $`\nu :=\varrho :\stackrel{~}{X}X`$ is the normalization of $`X.`$
* The secant cone $`\mathrm{Sec}_p(\stackrel{~}{X})_k^{r+1}`$ is a projective subspace $`_k^{t1}_k^{r+1}`$ and
$$X\mathrm{Nor}X=\mathrm{Sing}(\nu )=\varrho (\mathrm{Sec}_p(\stackrel{~}{X}\{p\})X$$
is a projective subspace $`_k^{t2}_k^r.`$
* The singular fibre $`\nu ^1(\mathrm{Sing}(\nu ))=\mathrm{Sec}_p(\stackrel{~}{X})\stackrel{~}{X}\stackrel{~}{X}`$ is a quadric in $`_k^{t1}=\mathrm{Sec}_p(\stackrel{~}{X}).`$
###### Proof.
Let $`B`$ be the graded normalization of $`A`$. Then, according to Theorem 5.6 resp. Theorem 6.9 we see that $`B`$ is the homogeneous coordinate ring of a variety $`\stackrel{~}{X}_k^{r+1}`$ of minimal degree. Moreover, by Theorem 5.3 resp. Theorem 6.9 there are linearly independent linear forms $`y_{t1},y_t,\mathrm{},y_rS_1`$ such that $`I(y_{t1},y_t,\mathrm{},y_r)`$ and $`B/A(S/(y_{t1},y_t,\mathrm{},y_r))(1)`$. Now, all claims except statement (c) follow as in Theorem 5.6.
To prove statement (c), we consider the prime $`๐ฐ:=(y_{t1},y_t,\mathrm{},y_r)/IA.`$ Then $`A/๐ฐ`$ is a polynomial ring in $`t1`$ indeterminates over $`k`$ and we have $`B/A(A/๐ฐ)(1).`$ In particular $`๐ฐB`$ is an ideal and we have an exact sequence $`0A/๐ฐB/๐ฐ(A/๐ฐ)(1)0`$. Therefore $`B/๐ฐ(A/๐ฐ)[z]/(f)`$ for some polynomial $`f=z^2+uz+v`$ with $`u(A/๐ฐ)_1`$ and $`v(A/๐ฐ)_2.`$ As
$$\nu ^1(\mathrm{Sing}(\nu ))=\mathrm{Proj}(B/๐ฐ)\mathrm{Sec}_p(\stackrel{~}{X})=_k^{t1}=\mathrm{Proj}((A/๐ฐ)[z])$$
the claims of (c) follows. โ
## 7. Varieties of almost minimal degree that are projections
We now wish to give a more detailed insight in the nature of those varieties of almost minimal degree which are projections of cones over rational normal scrolls. Let us first recall a few facts on such scrolls.
###### Remark 7.1.
(cf \[21, pp 94, 97, 108-110\] and ) A) Let $`n`$ and let $`l,a_1,\mathrm{},a_{l1}.`$ Let $`a_0=1,a_l=n`$ and assume that $`a_ia_{i1}>1`$ for $`i=1,\mathrm{},l.`$ Then, up to projective equivalence, the numbers $`a_1,a_2,\mathrm{},a_{l1}`$ define a unique rational normal $`l`$-fold scroll $`S_{a_1\mathrm{}a_{l1}}_k^n.`$ Keep in mind that $`S_{a_1\mathrm{}a_{l1}}`$ is smooth, rational, arithmetically Cohen-Macaulay and of dimension $`l.`$
B) Keep the notation of part A). After an appropriate linear coordinate transformation we may assume that the vanishing ideal of $`S_{a_1\mathrm{}a_{l1}}`$ in the polynomial ring $`k[x_0,\mathrm{},x_r]`$ is the ideal generated by the $`2\times 2`$-minors of the $`2\times (nl+1)`$-matrix
$$M_{a_1\mathrm{}a_{l1}}:=\left(\begin{array}{ccccc}x_0\mathrm{}x_{a_11}\hfill & |\hfill & x_{a_1+1}\mathrm{}x_{a_21}\hfill & \left|\mathrm{}\right|\hfill & x_{a_{l1}+1}\mathrm{}x_{n1}\hfill \\ x_1\mathrm{}x_{a_1}\hfill & |\hfill & x_{a_1+2}\mathrm{}x_{a_2}\hfill & \left|\mathrm{}\right|\hfill & x_{a_{l1}+2}\mathrm{}x_n\hfill \end{array}\right).$$
C) Let $`VGL_2(k)`$ and $`WGL_{nl}(k).`$ Then, the $`2\times 2`$-minors of the conjugate matrix $`VM_{a_1\mathrm{}a_{l1}}W^1`$ generate the same ideal as the $`2\times 2`$-minors of the matrix of $`M_{a_1\mathrm{}a_{l1}}.`$ So, if we subject $`M_{a_1\mathrm{}a_{l1}}`$ to regular $`k`$-linear row and column transformations, the $`2\times 2`$-minors of the resulting $`2\times (nl+1)`$-matrix still generate the vanishing ideal of the same scroll $`S_{a_1\mathrm{}a_{l1}}.`$
D) A $`2\times (nl+1)`$-matrix $`N`$ whose entries are linear forms in $`k[x_0,\mathrm{},x_r]`$ is said to be $`1`$-generic, if no conjugate of $`N`$ has a zero entry. Observe that the property of being $`1`$-generic is preserved under conjugation. Moreover, if $`N^{}`$ is obtained by deleting some columns from the $`1`$-generic matrix $`N,`$ then $`N^{}`$ is again $`1`$-generic.
Finally, let $`N`$ be a $`1`$-generic $`2\times (nl+1)`$-matrix whose entries are linear forms in $`k[x_0,\mathrm{},x_r].`$ Let $`y_0,\mathrm{},y_m`$ be a basis of the $`k`$-vector space $`Lk[x_0,\mathrm{},x_r]_1`$ generated by the entries of $`N.`$ Then, $`m>nl+1,`$ thus $`h:=mn+l>1.`$
Moreover, there are integers $`b_1,\mathrm{},b_{h1}`$ such that with $`b_0=1`$ and $`b_h=m`$ we have $`b_ib_{i1}>1`$ for $`i=1,\mathrm{},h,`$ and such that $`N`$ is conjugate to the $`2\times (ml+1)`$-matrix
$$N^{}:=\left(\begin{array}{ccccc}y_0\mathrm{}y_{b_11}\hfill & |\hfill & y_{b_1+1}\mathrm{}y_{b_21}\hfill & \left|\mathrm{}\right|\hfill & y_{b_{h1}+1}\mathrm{}y_{m1}\hfill \\ y_1\mathrm{}y_{b_1}\hfill & |\hfill & y_{b_1+2}\mathrm{}y_{b_2}\hfill & \left|\mathrm{}\right|\hfill & y_{b_{h1}+2}\mathrm{}y_m\hfill \end{array}\right).$$
So, by parts B) and C) the $`2\times 2`$-minors of $`N`$ generate the vanishing ideal of a rational normal $`h`$-fold scroll in $`_k^m=\mathrm{Proj}(k[y_0,\mathrm{},y_m]).`$
###### Remark 7.2.
A) Let $`s`$ and let $`\stackrel{~}{X}_k^s=\mathrm{Proj}(R),R=k[x_0,\mathrm{},x_s],`$ be a cone over a rational normal scroll $`S_{a_1\mathrm{}a_{l1}}_k^n`$ with $`n\{1,2,\mathrm{},s\}.`$ According to the previous remark we may assume that the vanishing ideal of $`\stackrel{~}{X}`$ in $`R`$ is generated by the $`2\times 2`$-minors of the $`2\times (nl+1)`$-matrix $`M_{a_1\mathrm{}a_{l1}}.`$ In this case (and with the convention that $`_k^1=\mathrm{}`$ and $`dim\mathrm{}=1`$), the vertex $`\mathrm{Sing}(\stackrel{~}{X})`$ of $`\stackrel{~}{X}`$ is given by $`_k^{sn1}=\mathrm{Proj}(R/(x_0,\mathrm{},x_n)R)`$ and so $`dim\mathrm{Sing}(\stackrel{~}{X})=sn1`$ and $`dim\stackrel{~}{X}=l+sn.`$
B) According to 7.1 C) the $`2\times 2`$-minors of any matrix obtained from $`M_{a_1\mathrm{}a_{l1}}`$ by $`k`$-linear row and column operations generate the vanishing ideal of $`\stackrel{~}{X}`$ in $`R.`$
C) Let $`N`$ be a $`1`$-generic $`2\times (nl+1)`$-matrix whose entries are linear forms in $`k[x_0,\mathrm{},x_n].`$ Let $`y_0,\mathrm{},y_m`$ be a basis of the $`k`$-space spanned by the entries of $`N`$ and let $`h:=mn+l.`$ Then, by parts A) and B) and by Remark 7.1 C), the $`2\times 2`$-minors of $`N`$ generate the vanishing ideal of a cone $`Y_k^s`$ over a rational normal $`h`$-fold scroll $`Z_k^m.`$ In particular $`dimY=h+sm=l+sn`$ and $`dim\mathrm{Sing}(Y)=sm1.`$
We now prove the result which shall be crucial in the rest of this chapter.
###### Theorem 7.3.
Let $`\stackrel{~}{X}_k^{r+1}`$ be a (cone over a) rational normal scroll and let $`\varrho :_k^{r+1}\{p\}_k^r`$ be a linear projection from a point $`p_k^{r+1}\stackrel{~}{X}.`$ Then, there is a (cone over a) rational normal scroll $`Y_k^r`$ such that $`Y\varrho (\stackrel{~}{X}),dimY=dim\stackrel{~}{X}+1`$ and $`dim\mathrm{Sing}(\stackrel{~}{X})dim\mathrm{Sing}(Y)dim\mathrm{Sing}(\stackrel{~}{X})+3.`$
###### Proof.
According to Remark 7.2 A) we may assume that the vanishing ideal of $`\stackrel{~}{X}`$ in $`S^{}=k[x_0,\mathrm{},x_{r+1}]`$ is generated by the $`2\times 2`$-minors of the $`2\times (nl+1)`$-matrix
$$M:=\left(\begin{array}{ccccc}x_0\mathrm{}x_{a_11}\hfill & |\hfill & x_{a_1+1}\mathrm{}x_{a_21}\hfill & \left|\mathrm{}\right|\hfill & x_{a_{l1}+1}\mathrm{}x_{n1}\hfill \\ x_1\mathrm{}x_{a_1}\hfill & |\hfill & x_{a_1+2}\mathrm{}x_{a_2}\hfill & \left|\mathrm{}\right|\hfill & x_{a_{l1}+2}\mathrm{}x_n\hfill \end{array}\right)$$
with appropriate integers $`n,l,a_1,\mathrm{},a_{l1}`$ such that, with $`a_0=1,a_l=n,`$ we have $`a_ia_{i1}>1`$ for $`i=1,\mathrm{},l,nr+1,dim\stackrel{~}{X}=l+r+1n`$ and $`dim\mathrm{Sing}(\stackrel{~}{X})=rn.`$
Let $`p:=(c_0:c_1:\mathrm{}:c_{r+1}).`$ As $`p\stackrel{~}{X}`$ there are two different indices $`i,j\{0,1,\mathrm{},n1\}\{a_1,a_2,\mathrm{},a_{l1}\}`$ such that
$$\delta :=det\left(\begin{array}{cc}c_i& c_j\\ c_{i+1}& c_{j+1}\end{array}\right)0.$$
Without loss of generality we may assume that $`c_{i+1}0.`$ Define
$$y_\alpha :=\{\begin{array}{cc}x_\alpha ,\hfill & \text{if}\alpha =i+1,\hfill \\ x_\alpha \frac{c_\alpha }{c_{i+1}}x_{i+1},\hfill & \text{if}\alpha \{0,\mathrm{},r+1\}\{i+1\}.\hfill \end{array}$$
Then $`S^{}=k[y_0,y_1,\mathrm{},y_{r+1}]`$ and with respect to the coordinates $`y_0,\mathrm{},y_{r+1}`$ the point $`p`$ may be written as $`(0:\mathrm{}:0:1:0:\mathrm{}:0)`$ with the entry โ1โ in the $`(i+1)`$-th position. Therefore we may assume that the projection $`\varrho `$ is induced by the inclusion map
$$S^{\prime \prime }:=k[y_0,\mathrm{},y_i,y_{i+2},\mathrm{},y_{r+1}]S^{}.$$
We now express the indeterminates which occur in the matrix $`M`$ in terms of the variables $`y_\alpha :`$
$$x_\alpha =\{\begin{array}{cc}y_\alpha ,\hfill & \text{for}\alpha =i+1,\hfill \\ y_\alpha +\frac{c_\alpha }{c_{i+1}}y_{i+1},,\hfill & \text{if}\alpha \{0,\mathrm{},r+1\}\{i+1\}.\hfill \end{array}$$
Let us first assume that $`i<j.`$ Then, the $`2\times 2`$-submatrix $`U`$ of $`M`$ which contains $`x_i`$ and $`x_j`$ in its first row takes the form
$$U=\left(\begin{array}{cc}y_i+\frac{c_i}{c_{i+1}}y_{i+1}& y_j+\frac{c_j}{c_{i+1}}y_{i+1}\\ y_{i+1}& y_{j+1}+\frac{c_{j+1}}{c_{i+1}}y_{i+1}\end{array}\right).$$
Now performing sucessively $`k`$-linear row and column operations we finally get the following transformed matrix
$$U^{}=\left(\begin{array}{cc}y_i& y_j\frac{c_i}{c_{i+1}}y_{j+1}\frac{c_{j+1}}{c_{i+1}}y_i\frac{\delta }{c_{i+1}^2}y_{i+1}\\ y_{i+1}& y_{j+1}\end{array}\right).$$
If $`i+1=j,`$ then by performing $`k`$-linear row and column operations, $`U`$ can be brought to the form
$$U^{}=\left(\begin{array}{cc}y_i& \frac{c_i}{c_{i+1}}y_{i+2}\frac{c_{i+2}}{c_{i+1}}y_i\frac{\delta }{c_{i+1}^2}y_{i+1}\\ y_{i+1}& y_{i+2}\end{array}\right).$$
Let $`M^{}`$ be the $`2\times (nl+1)`$-matrix which is obtained if the above row and column operations are performed with the whole matrix $`M.`$ Observe that the submatrix $`U`$ of $`M`$ is transformed into the submatrix $`U^{}`$ of $`M^{},`$ which sits in the same columns as $`U.`$ Now, as $`\frac{\delta }{c_{i+1}^2}0`$ we may add appropriate $`k`$-multiples of the columns of $`U^{}`$ to the columns of $`M^{}`$ to remove the indeterminate $`y_{i+1}`$ from all entries of $`M^{}`$ which do not belong to the two columns of $`U^{}.`$ So, we get a $`2\times (nl+1)`$-matrix $`\stackrel{~}{M},`$ conjugate to $`M.`$ In particular, $`\stackrel{~}{M}`$ is $`1`$-generic (cf Remark 7.1 D)) and the entries of $`\stackrel{~}{M}`$ span the same $`k`$-space as the entries of $`M,`$ namely $`_{t=0}^nkx_t=_{t=0}^nky_t.`$ Now, let $`N`$ be the matrix of size $`2\times (nl1)=2\times (n1(l+1)+1)`$ obtained by deleting the two columns of $`U^{}`$ from $`\stackrel{~}{M}.`$ Then $`N`$ is $`1`$-generic (cf Remark 7.1 D)) and $`y_{i+1}`$ does not appear in $`N.`$ So, the entries of $`N`$ span a subspace
$$L\mathrm{\Sigma }_{t=0,ti+1}^nky_tk[y_0,\mathrm{},y_i,y_{i+1},\mathrm{},y_n]S^{\prime \prime },$$
whose dimension $`m`$ is such that $`nm\{1,2,3,4\}.`$
By Remark 7.2 C) the ideal $`I_2(N)S^{\prime \prime }`$ generated by the $`2\times 2`$-minors of $`N`$ is the vanishing ideal of a cone $`Y_k^r=\mathrm{Proj}(S^{\prime \prime })`$ over a rational normal scroll such that $`dimY=(l+1)+r(n1)=dim\stackrel{~}{X}+1`$ and $`dim\mathrm{Sing}(Y)=rm1=dim\mathrm{Sing}(\stackrel{~}{X})+(nm1).`$ As $`nm1\{0,1,2,3\}`$ we get $`dim\mathrm{Sing}(\stackrel{~}{X})dim\mathrm{Sing}(Y)dim\mathrm{Sing}(\stackrel{~}{X})+3.`$ According to Remark 7.2 B) the ideal $`I_2(\stackrel{~}{M})S^{}`$ generated by the $`2\times 2`$-minors of $`\stackrel{~}{M}`$ is the vanishing ideal of $`\stackrel{~}{X}.`$ Therefore $`I_2(\stackrel{~}{M})S^{\prime \prime }`$ is the vanishing ideal of $`\varrho (\stackrel{~}{X}).`$ As each $`2\times 2`$-minor of $`N`$ is a $`2\times 2`$-minor of $`\stackrel{~}{M},`$ it follows $`I_2(N)I_2(\stackrel{~}{M})S^{\prime \prime }`$ and hence $`Y\varrho (\stackrel{~}{X}).`$
This settles the case $`i<j.`$ If $`j<i`$ we first commute the columns of $`U`$ and then conclude as above. โ
We now apply the previous result to varieties of almost minimal degree. We still keep the convention that $`_k^1=\mathrm{}`$ and $`dim\mathrm{}=1,`$ and begin with a preliminary remark.
###### Remark 7.4.
(cf ) A) Let $`l,n,a_1,\mathrm{},a_{l1}`$ be as in Remark 7.1 and consider the rational $`l`$-fold scroll $`S_{a_1\mathrm{}a_{l1}}.`$ We set $`a_0=1,a_l=n,d_i=a_ia_{i1}1,`$ for $`i=1,\mathrm{},l`$ and define the linear subspaces
$$๐ผ_i=\mathrm{Proj}(S/\mathrm{\Sigma }_{j\{a_{i1}+1,\mathrm{},a_i\}}Sx_j)=_k^{d_i}_k^n,i=1,\mathrm{},l.$$
For each $`i\{1,\mathrm{},l\}`$ we consider the Veronese embedding
$$\nu _i:_k^1๐ผ_i,(s:t)(s^{d_i}:s^{d_i1}t:\mathrm{}:st^{d_i1}:t^{d_i}),$$
so that $`\nu _i(_k^1)๐ผ_i=_k^{d_i}`$ is a rational normal curve. Now, for each $`q_k^1`$ let
$$๐ผ(q)=\nu _1(q),\mathrm{},\nu _l(q)=_k^{l1}_k^n$$
be the projective space spanned by the $`l`$ points $`\nu _1(q),\mathrm{},\nu _l(q).`$ Then $`qq^{}`$ implies $`๐ผ(q)๐ผ(q^{})=\mathrm{}`$ for all $`q,q^{}_k^1.`$ Moreover, $`S_{a_1\mathrm{}a_{l1}}=_{q_k^1}๐ผ(q).`$
B) Keep the above notation. For each $`1il`$ let $`๐_i=k^{d_i+1}k^{n+1}`$ be the affine cone over $`๐ผ_i,`$ fix $`q=(s:t)_k^1`$ and set $`v_i=(s^{d_i},s^{d_i1}t,\mathrm{},t^{d_i}).`$ Moreover, for each $`1il`$ let $`w_i=(w_{i0},\mathrm{},w_{id_i})๐_i`$ such that $`(w_{i0}:\mathrm{}:w_{id_i})๐ผ_i\nu _i(q)`$ is a point on the tangent of the rational normal curve $`\nu _i(_k^1)๐ผ_i`$ in the point $`\nu _i(q).`$ In particular, $`v_i`$ and $`w_i`$ are lineraly independent. Now let $`\pi :k^{n+1}\{0\}_k^n=(k^{n+1})`$ be the canonical projection and let
$$r=\pi (\mathrm{\Sigma }_{i=1}^lr_iv_i)๐ผ(q)=\pi (_{i=1}^lkv_i\{0\})=(_{i=1}^lkv_i),$$
where $`(r_1,\mathrm{},r_l)k^l\{0\}.`$ Then the tangent space to the scroll $`S_{a_1\mathrm{}a_{l1}}`$ in the point $`r`$ is given by
$$T_r(S_{a_1\mathrm{}a_{l1}})=(k(\mathrm{\Sigma }_{i=1}^lr_iw_i)+_{i=1}^lkv_i)=๐ผ(q)\pi (\mathrm{\Sigma }_{i=1}^lr_iw_i).$$
From this we easily deduce that $`T_r(S_{a_1\mathrm{}a_{l1}})T_r^{}(S_{a_1\mathrm{}a_{l1}})=๐ผ(q)`$ for all $`r,r^{}๐ผ(q)`$ with $`rr^{}.`$
C) Now, let $`s`$ such that $`n<s`$ and let $`\stackrel{~}{X}_k^s=\mathrm{Proj}(R),R=k[x_0,\mathrm{},x_n],`$ be a cone over the rational normal $`l`$-fold scroll
$$S_{a_1\mathrm{}a_{l1}}_k^n=\mathrm{Proj}(R/(x_{n+1},\mathrm{},x_s)R=k[x_0,\mathrm{},x_n]).$$
Then, the vertex of $`\stackrel{~}{X}`$ is given by $`\mathrm{Sing}(\stackrel{~}{X})=\mathrm{Proj}(R/(x_0,\mathrm{},x_n))=_k^{sn1}_k^s`$ (cf Remark 7.2 A)). Now, for each $`q_k^1`$ let
$$๐ฝ(q)=๐ผ(q)\mathrm{Sing}(\stackrel{~}{X})=_k^{l+sn1}=_k^{dim\stackrel{~}{X}}_k^s$$
be the linear subspace spanned by $`๐ผ(q)=_k^{l1}_k^s`$ and the vertex $`\mathrm{Sing}(\stackrel{~}{X})`$ of $`\stackrel{~}{X}.`$ Then by part A), $`qq^{}`$ implies that $`๐ฝ(q)๐ฝ(q^{})=\mathrm{Sing}(\stackrel{~}{X})`$ for all $`q,q^{}_k^1`$ and moreover $`\stackrel{~}{X}=_{q_k^1}๐ฝ(q).`$
It also follows easily from part B), that for any $`q_k^1`$ and any $`r๐ผ(q)\mathrm{Sing}(\stackrel{~}{X})`$ the tangent space of $`\stackrel{~}{X}`$ at $`r`$ is given by $`T_r(\stackrel{~}{X})=T_{\stackrel{~}{r}}(S_{a_1\mathrm{}a_{l1}})\mathrm{Sing}(\stackrel{~}{X})=_k^d,`$ where $`\stackrel{~}{r}=(r_0:\mathrm{}:r_n)`$ is the canonical projection of $`r`$ from $`\mathrm{Sing}(\stackrel{~}{X}).`$ As a consequence of the last statement in part B) we thus get for all $`q_k^1`$ and all $`r,r^{}๐ฝ(q)\mathrm{Sing}(\stackrel{~}{X}):`$ If $`rr^{},`$ then $`T_r(\stackrel{~}{X})T_r^{}(\stackrel{~}{X})=๐ฝ(q).`$
###### Theorem 7.5.
Let $`X_k^r`$ be a variety of almost minimal degree which is the projection of a (cone over a) rational normal scroll $`\stackrel{~}{X}_k^{r+1}`$ with $`dim\mathrm{Sing}(\stackrel{~}{X})=:h`$ from a point $`p_k^{r+1}\stackrel{~}{X}.`$ Then
* $`X`$ is contained in a (cone over a) rational normal scroll $`Y_k^r`$ such that $`\mathrm{codim}_Y(X)=1`$ and $`hdim\mathrm{Sing}(Y)h+3.`$
* $`X`$ is of arithmetic depth $`th+5.`$
###### Proof.
(a): This is clear by Theorem 7.3.
(b): Assume first that $`X`$ is not arithmetically Cohen-Macaulay. Then, the non CM-locus $`Z`$ of $`X`$ is a linear subspace $`_k^{t2}`$ of $`_k^r`$ (cf Theorem 5.6 (d), (f)). As $`\mathrm{codim}_Y(X)=1,๐ช_{X,x}`$ is a Cohen-Macaulay ring for each point $`xX\mathrm{Sing}(Y).`$ It follows that $`_k^{t2}=ZX\mathrm{Sing}(Y)`$ and hence $`t2dim\mathrm{Sing}(Y)h+3.`$
Now, let $`X`$ be arithmetically Cohen-Macaulay. In this case we conclude by a geometric argument which in fact also implies in the previous case. Let $`d=dimX.`$ After an appropriate change of coordinates, we may assume that we are in the situation of Remark 7.4 B) and C). So, we may write $`\stackrel{~}{X}=_{q_k^1}๐ฝ(q),`$ where $`๐ฝ(q)=_k^{d1}_k^r`$ is a linear subspace for all $`q_k^1.`$ Let $`U=\{(q,q^{})_k^1\times _k^1|qq^{}\}.`$ Then, according to Remark 7.4 C) we have $`๐ฝ(q)๐ฝ(q^{})=\mathrm{Sing}(\stackrel{~}{X})=_k^h,`$ whenever $`(q,q^{})U.`$ Now, for each pair $`(q,q^{})U`$ consider the linear subspace
$$(q,q^{})=๐ฝ(q)\{p\}๐ฝ(q^{})\{p\}_k^{r+1}.$$
Observe that $`\mathrm{Sing}(\stackrel{~}{X})(q,q^{})`$ and $`dim(q,q^{})h+2`$ for all $`(q,q^{})U.`$ Moreover $`dim(q,q^{})=h+2`$ if and only if $`p๐ฝ(q)๐ฝ(q^{}).`$ Consequently we have $`dim(q,q^{})=h+2`$ if and only if there is a line running through $`p`$ and intersecting $`๐ฝ(q)`$ and $`๐ฝ(q^{}).`$ Clearly such a line is contained in $`(q,q^{})`$ and its intersection points with $`๐ฝ(q)`$ and $`๐ฝ(q^{})`$ are different as $`p๐ฝ(q)๐ฝ(q^{}).`$
Let $`VU`$ be the closed subset of all pairs $`(q,q^{})`$ for which $`dim(q,q^{})=h+2.`$ It follows that the union of all proper secant lines of $`\stackrel{~}{X}`$ which run through $`p`$ is a subset of $`W=_{(q,q^{})V}(q,q^{}).`$ Moreover, it follows from the last statement of Remark 7.4 C), that for each point $`q_k^1`$ there is at most one point $`r(q)๐ฝ(q)\mathrm{Sing}(\stackrel{~}{X})`$ such that there is a tangent line $`l(q)=_k^1`$ of $`\stackrel{~}{X}`$ at $`r(q)`$ running through $`p.`$ Let $`T_k^1`$ be the closed subset of all $`q_k^1`$ for which this happens. Then, all tangents to non-singular points of $`\stackrel{~}{X}`$ through $`p`$ are contained in $`Y=_{qT}l(q).`$ Finally observe that the remaining tangents are the lines running trough $`p`$ and $`\mathrm{Sing}(\stackrel{~}{X}).`$ It follows $`_k^{t1}=\mathrm{Sec}_p(\stackrel{~}{X})WY\mathrm{Sing}(\stackrel{~}{X})\{p\}`$ and hence $`t1\mathrm{max}\{dimW,dimY,h+1\}.`$ As $`\{(q,q^{})|(q,q^{})V\}`$ is a family of linear $`(h+2)`$-subspaces of $`_k^{r+1},`$ it follows $`dimWh+2+dimVh+4.`$ As $`\{l(q)|qT\}`$ is a family of lines we have $`dimY1+dimT2.`$ So we get $`t1h+4,`$ hence $`th+5.`$
###### Corollary 7.6.
Let $`X_k^r`$ be a variety of almost minimal degree which is a projection of a rational normal scroll $`\stackrel{~}{X}_k^{r+1}`$ from a point $`p_k^{r+1}\stackrel{~}{X}.`$ Then $`X`$ is of arithmetic depth $`t4.`$
###### Proof.
Clear from Theorem 7.5. โ
As a final comment of this section let us say something about the exceptional case of projections of the Veronese surface.
###### Remark 7.7.
(The exceptional case) Let us recall that the Veronese surface $`F_k^5`$ is defined by the $`2\times 2`$-minors of the matrix
$$M=\left(\begin{array}{ccc}x_0& x_1& x_2\\ x_1& x_3& x_4\\ x_2& x_4& x_5\end{array}\right).$$
Let $`p_K^5F`$ denote a closed point. Suppose that $`\mathrm{rank}M_p=3,`$ i.e. the case of a generic point and remember that $`detM=0`$ defines the secant variety of $`F.`$ Then the projection of $`F`$ from $`p`$ defines a surface $`X_K^4`$ of almost minimal degree and $`0ptA=1.`$
Recall that $`dim_k(I)_2=0`$ (cf. Corollary 4.4 C)). Therefore the surface $`X`$ is cut out by cubics, i.e. it is not contained in a variety of minimal degree.
## 8. Betti numbers
Our next aim is to study the Betti numbers of $`A`$ if $`X`$ is of almost minimal degree, non-arithmetically Cohen-Macaulay and a projection of a (cone over a) rational normal scroll.
###### Lemma 8.1.
Assume that $`X_k^r`$ is of almost minimal degree and of arithmetic depth $`d=dimX.`$ Let $`B=\mathrm{Hom}_A(K(A),K(A)).`$ Then
$$\mathrm{Tor}_i^S(k,B)\{\begin{array}{cc}k(0)k(1),\hfill & \text{if }i=0,\hfill \\ k^{b_i}(i1),\hfill & \text{if }0<ird,\hfill \\ 0,\hfill & \text{if }rd<i,\hfill \end{array}$$
where $`b_i=(r+1d)\left(\genfrac{}{}{0pt}{}{rd}{i}\right)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right)`$ for $`1ird.`$
###### Proof.
By Theorem 5.3 (a) the $`A`$-module $`B`$ is Cohen-Macaulay and hence of depth $`d+1`$ over $`S.`$ Therefore $`\mathrm{Tor}_i^S(k,B)=0`$ for all $`i>rd.`$ According to Theorem 5.3 (b) there is a short exact sequence of graded $`S`$-modules
(8.1)
$$0ABC0,C(S/(y_{t2},\mathrm{},y_r)S)(1),$$
where $`y_0,\mathrm{},y_r`$ form a generic set of linear forms of $`S.`$ In particular, $`C`$ is of dimension $`t1<d`$ and generated by a single element of degree 1. This already shows that $`\mathrm{Tor}_0^S(k,B)k(0)k(1).`$
Applying cohomology to the above short exact sequence we get an isomorphism $`H^{d+1}(A)H^{d+1}(B)`$ which shows that $`\text{end }H^{d+1}(B)=d`$ (cf Theorem 4.2 (b)). As $`0ptB=d+1`$ it follows $`\mathrm{reg}B=1.`$ Moreover the above exact sequence yields
$$dim_kB_1=dim_kA_1+1=r+2=dim_k(S(0)S(1))_1.$$
So, the graded $`S`$-module $`B`$ must have a minimal free resolution of the form
$$0S^{b_{rd}}(r+d1)\mathrm{}S^{b_i}(i1)\mathrm{}S^{b_1}(2)SS(1)B0$$
with $`b_1,\mathrm{},b_{rd}.`$
As $`B`$ is a Cohen-Macaulay module of dimension $`d+1,`$ regularity 1 and of multiplicity $`\mathrm{deg}X=rd+2`$ (cf Theorem 5.3 (c)) its Hilbert series is given by
$$F(\lambda ,B)=\frac{1+(r+1d)\lambda }{(1\lambda )^{d+1}}.$$
On the use of Betti numbers $`b_i`$ we also may write
$$F(\lambda ,B)=\frac{1}{(1\lambda )^{d+1}}(1+\lambda +\underset{i=0}{\overset{rd}{}}(1)^ib_i\lambda ^{i+1}).$$
Comparing coefficients we obtain
$$b_i=(r+1d)\left(\genfrac{}{}{0pt}{}{rd}{i}\right)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right),i=1,\mathrm{},rd,$$
as required. โ
Next we recall a well known result about the Betti numbers of a variety of minimal degree.
###### Lemma 8.2.
Let $`Y_k^r`$ be a variety of minimal degree with $`dimY=d+1.`$ Let $`U`$ be the homogeneous coordinate ring of $`Y.`$ Then
$$\mathrm{Tor}_i^S(k,U)\{\begin{array}{cc}k,\hfill & \text{if }i=0,\hfill \\ k^{c_i}(i1),\hfill & \text{if }0<i<rd,\hfill \\ 0,\hfill & \text{if }rdi,\hfill \end{array}$$
where $`c_i=i\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right)`$ for $`1i<rd.`$
###### Proof.
This is well known (cf for instance ). In fact the Eagon-Northcott complex provides a minimal free resolution of $`U`$ over $`S.`$
###### Theorem 8.3.
Let $`X_k^r`$ be a variety of almost minimal degree which is the projection of a (cone over a) rational normal scroll $`\stackrel{~}{X}_k^{r+1}`$ from a point $`p_k^{r+1}\stackrel{~}{X}.`$ Assume that $`t:=0ptAd:=dimX.`$ Then
$$\mathrm{Tor}_i^S(k,A)\{\begin{array}{cc}k,\hfill & \text{if}i=0,\hfill \\ k^{u_i}(i1)k^{v_i}(i2),\hfill & \text{if}0<irt+1,\hfill \\ 0,\hfill & \text{if}rt+1<i,\hfill \end{array}$$
where
* + $`u_1=t+\left(\genfrac{}{}{0pt}{}{r+1d}{2}\right)d2,`$
+ $`i\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right)u_i(r+1d)\left(\genfrac{}{}{0pt}{}{rd}{i}\right)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right),\text{if }\mathrm{\hspace{0.33em}1}<i<r2d+t1,`$
+ $`u_i=i\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right),\text{if }r2d+t1i<rd,`$
+ $`u_i=0,\text{if }rdi<rt+1,`$
* + $`\mathrm{max}\{0,\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right)(i+2)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right)\}v_i\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right),\text{if }\mathrm{\hspace{0.33em}1}i<r2d+t2,`$
+ $`v_i=\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right)(i+2)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right),\text{if }r2d+t2i<rd,`$
+ $`v_i=\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right),\text{if }rdirt+1.`$
Moreover, $`v_iu_{i+1}=\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right)(rd+1)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right)+\left(\genfrac{}{}{0pt}{}{rd}{i+2}\right)`$ for all $`1i<rd.`$
###### Proof.
As $`0ptA=t`$ and $`\mathrm{reg}A=2`$ (cf Theorem 4.2) the modules $`\mathrm{Tor}_i^S(k,A)`$ behave as stated in the main equality. So let the numbers $`u_i,v_i`$ be defined according to this main equality with the convention that $`u_i=v_i=0`$ for $`i>rt+1.`$ Moreover let $`b_i`$ and $`c_i`$ be as in Lemma 8.1 resp. 8.2 with the convention that $`b_i=0`$ for $`i>rd`$ and $`c_i=0`$ for $`ird.`$
According to Corollary 7.6 there is a (cone over a) rational normal scroll $`Y_k^r`$ of dimension $`d+1`$ such that $`XY.`$ Let $`JS`$ be the vanishing ideal of $`Y`$ and let $`U:=S/J`$ be the homogeneous coordinate ring of $`Y.`$ The short exact sequence $`0I/JUA0`$ together with Lemma 8.2 implies short exact sequences
(8.2)
$$0k^{c_i}(i1)k^{u_i}(i1)k^{v_i}(i2)\mathrm{Tor}_{i1}^S(k,I/J)0$$
for all $`i1.`$
Keep in mind that $`\mathrm{beg}(I/J)=2,u_1=dim_kI_2=t+\left(\genfrac{}{}{0pt}{}{r+1d}{2}\right)d2`$ (cf Corollary 4.4 (c)) and $`dim_kJ_2=\left(\genfrac{}{}{0pt}{}{rd}{2}\right)`$ (cf Lemma 8.2) so that
$$dim_k(I/J)_2=r2d+t2.$$
Whence, by Greenโs Linear Syzygy Theorem (cf \[11, Theorem 7.1\]) we have
$$\mathrm{Tor}_j^S(k,I/J)_{j+2}=0\text{ for all }jr2d+t2.$$
So, the sequence 8.2 yields that $`u_i=c_i`$ for all $`ir2d+t1.`$ This proves statement (a) in the range $`ir2d+t1.`$ The sequence 8.2 also yields that $`c_iu_i`$ for all $`ir2d+t1.`$
Next we consider the short exact sequence of graded $`S`$-modules 8.1. In particular we have $`\mathrm{Tor}_i^S(k,C)k^{\left(\genfrac{}{}{0pt}{}{rt+2}{i}\right)}(i1)`$ for all $`i_0.`$ So, by the sequence 8.1 and in view of Lemma 8.2 we get exact sequences
(8.3)
$$\begin{array}{c}k^{b_{i+1}}(i2)k^{\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right)}(i2)k^{u_i}(i1)k^{v_i}(i2)\hfill \\ \hfill k^{b_i}(i1)k^{\left(\genfrac{}{}{0pt}{}{rt+2}{i}\right)}(i1)k^{u_{i1}}(i)k^{v_{i1}}(i1)k^{b_{i1}}(i)\end{array}$$
for all $`i2.`$ Now, we read off that $`u_ib_i`$ for all $`i1`$ and statement (a) is proved completely.
The sequence 8.3 also yields that
(8.4)
$$v_i=u_{i+1}b_{i+1}+\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right)\text{for all }i1.$$
Observe that $`c_{i+1}b_{i+1}=(i+2)\left(\genfrac{}{}{0pt}{}{rd}{i+1}\right)`$ for $`1i<rd.`$ If $`1i<r2d+t1,`$ statement (a) gives $`c_{i+1}u_{i+1}b_{i+1}`$ so that
$$c_{i+1}b_{i+1}+\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right)v_i\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right).$$
This proves the first estimate of statement (b).
If $`r2d+t2i<rd,`$ statement (a) yields $`u_{i+1}=c_{i+1},`$ hence $`v_i=c_{i+1}b_{i+1}+\left(\genfrac{}{}{0pt}{}{rt+2}{i+1}\right).`$ This proves the second claim of statement (b). Finally, if $`rdi<rt+1`$ statement (a) and Lemma 8.2 yield that $`u_{i+1}=c_{i+1}=0.`$ Now the last claim of the statement (b) follows by 8.4. โ
## 9. Examples
In this final section we present a few examples which illustrate the previous results. All calculations of the โgraded Betti numbersโ $`u_i`$ and $`v_i`$ (cf Theorem 8.3) have been performed by means of the computer algebra system Singular . As for rational scrolls and their secant varieties we refer to and .
First, we present three examples of $`3`$-folds $`X`$ of almost minimal degree in $`_k^{11}`$, one of them being defined by $`32`$ quadrics, the second by $`32`$ quadrics and $`1`$ cubic, the third by $`32`$ quadrics and $`3`$ cubics. These examples show that, contrary to the number of defining quadrics (cf Corollary 4.4 (c) ), the number of defining cubics may vary if the embedding dimension $`r`$, the dimension $`d`$ and the arithmetic $`0ptt`$ of $`X`$ are fixed. Notice that each smooth variety $`X_k^{11}`$ of almost minimal degree which is not arithmetically Cohen-Macaulay is obtained by projecting a rational scroll $`\stackrel{~}{X}_k^{12}`$ from a point $`p_k^{12}\mathrm{Sec}(\stackrel{~}{X})`$ (cf Theorem 5.6).
We first fix some notation. Let $`l,n,d_1,\mathrm{},d_l`$ such that $`d_1d_2leq\mathrm{}d_l`$ and $`_{i=1}^ld_i=nl+1.`$ Let $`a_i=i1+_{j=1}^id_j,i=1,\mathrm{},l1.`$ Then, we write $`S(d_1,\mathrm{},d_l)`$ for the rational normal scroll $`S_{a_1\mathrm{},a_{l1}}`$ (cf Remark 7.1).
###### Example 9.1.
A) Let $`\stackrel{~}{X}_k^{12}`$ the $`3`$-scroll $`S(2,2,6)`$, thus the smooth variety of degree $`10`$ defined by the $`2\times 2`$ minors of the matrix
$$\left(\begin{array}{cccccccc}x_0& x_1|x_3& x_4|x_6& x_7& x_8& x_9& x_{10}& x_{11}\\ x_1& x_2|x_4& x_5|x_7& x_8& x_9& x_{10}& x_{11}& x_{12}\end{array}\right).$$
Its homogeneous coordinate ring is
$$B=k[(s,t)^2u^5,(s,t)^2v^5,(s,t)^6w]k[s,t,u,v,w].$$
Projecting $`\stackrel{~}{X}`$ from the point
$$p_1=(0:0:0:0:0:0:0:0:0:1:0:0:0)_k^{12}\stackrel{~}{X}$$
we get a non-degenerate variety $`X_k^{11}`$ of dimension $`3`$ and of degree $`10`$, (cf Remark 3.3 A) ). Let $`S`$ denote a polynomial ring in $`12`$ indeterminates, let $`IS`$ be the homogeneous vanishing ideal and let $`A=S/I`$ be the homogeneous coordinate ring of $`X`$. Also, consider the only not necessarily vanishing graded Betti numbers
$$u_i:=dim_k\mathrm{Tor}_i^S(k,A)_{i+1},v_i:=dim_k\mathrm{Tor}_i^S(k,A)_{i+2}$$
of $`X`$. These numbers present themselves as shown below:
| $`i`$ | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 | 10 | 11 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`u_i`$ | 32 | 130 | 234 | 234 | 140 | 48 | 7 | 0 | 0 | 0 | 0 |
| $`v_i`$ | 0 | 20 | 155 | 456 | 728 | 728 | 486 | 220 | 66 | 12 | 1 |
In particular $`t:=0ptA=1`$ so that $`X`$ cannot be of minimal degree and hence $`\mathrm{deg}X=10=113+2`$. Therefore, $`X`$ is of almost minimal degree and of arithmetic $`0pt1`$. In particular the projection map $`\nu :\stackrel{~}{X}X`$ is an isomorphism (cf Theorem 5.7) and so $`X`$ becomes smooth. Observe, that $`I`$ is generated by $`32`$ quadrics.
B) Let $`\stackrel{~}{X}_k^{12}`$ be as in part A) but project $`\stackrel{~}{X}`$ from the point
$$p_2=(0:0:0:0:0:0:0:0:0:0:1:0:0)_k^{12}\stackrel{~}{X}.$$
Again let $`X_k^{11}`$ be the image of $`\stackrel{~}{X}`$ under this projection and define $`S,I,A`$ as in part A). Now the Betti numbers $`u_i,v_i`$ present themselves as follows:
| $`i`$ | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 | 10 | 11 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`u_i`$ | 32 | 131 | 234 | 234 | 140 | 48 | 7 | 0 | 0 | 0 | 0 |
| $`v_i`$ | 1 | 20 | 155 | 456 | 728 | 728 | 486 | 220 | 66 | 12 | 1 |
So, as above, we see that $`X`$ is a smooth variety of almost minimal degree having dimension $`3`$ and arithmetic $`0pt1`$. Observe, that now $`I`$ is minimally generated by $`32`$ quadrics and $`1`$ cubic. So, if the same scroll $`\stackrel{~}{X}=S(2,2,6)_k^{12}`$ is projected from two different points $`p_1,p_2_k^{12}\mathrm{Sec}(\stackrel{~}{X})`$, the homological nature of the projection $`X_k^{11}`$ may differ.
C) Now, consider the scroll $`\stackrel{~}{X}:=S(2,4,4)_k^{12}`$, so that $`\stackrel{~}{X}`$ is the smooth variety of dimension $`3`$ and degree $`10`$ defined by the $`2\times 2`$-minors of the matrix
$$\left(\begin{array}{cccccccc}x_0& x_1|x_3& x_4& x_5& x_6|x_8& x_9& x_{10}& x_{11}\\ x_1& x_2|x_4& x_5& x_6& x_7|x_9& x_{10}& x_{11}& x_{12}\end{array}\right).$$
Its homogeneous coordinate ring is
$$B=k[(s,t)^2u^3,(s,t)^4v,(s,t)^4w]k[s,t,u,v,w].$$
Define $`X_k^{11}`$ as the projection of $`\stackrel{~}{X}`$ from the point $`p_2_k^{12}\stackrel{~}{X}`$ (cf part B) ). In this case, the Betti numbers $`u_i`$ and $`v_i`$ of $`X`$ take the values listed in the following table:
| $`i`$ | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 | 10 | 11 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`u_i`$ | 32 | 133 | 248 | 234 | 140 | 48 | 7 | 0 | 0 | 0 | 0 |
| $`v_i`$ | 3 | 34 | 155 | 456 | 728 | 728 | 486 | 220 | 66 | 12 | 1 |
So, again $`X_k^{11}`$ is a smooth variety of almost minimal degree having dimension $`3`$ and arithmetic $`0pt1`$. But this time, besides $`32`$ quadrics three cubics are needed to generate the homogeneous vanishing ideal $`I`$ of $`X`$.
The previous example where all of arithmetic $`0pt1`$ and of dimension $`3`$. By projecting rational $`3`$-scrolls in $`_k^{12}`$ from appropriate points we also may obtain $`3`$-dimensional varieties $`X_k^{11}`$ of almost minimal degree and of arithmetic $`0pt`$ not equal to 1. We present two examples to illustrate this.
###### Example 9.2.
A) Next consider the $`3`$-scroll $`\stackrel{~}{X}:=S(3,3,4)_k^{12}`$ defined by the $`2\times 2`$-minors of the matrix
$$\left(\begin{array}{cccccccc}x_0& x_1& x_2|x_4& x_5& x_6|x_8& x_9& x_{10}& x_{11}\\ x_1& x_2& x_3|x_5& x_6& x_7|x_9& x_{10}& x_{11}& x_{12}\end{array}\right).$$
$`\stackrel{~}{X}`$ has the homogeneous coordinate ring
$$B=k[(s,t)^3u^2,(s,t)^3v^2,(s,t)^4w]k[s,t,u,v,w].$$
We project $`\stackrel{~}{X}`$ from the point
$$p_3=(0:0:0:0:0:0:1:0:0:0:0:0:0)_k^{12}\stackrel{~}{X}.$$
Like above we get a non-degenerate variety $`X_k^{11}`$ of degree $`10=\mathrm{codim}X+2`$ and Betti numbers:
| $`i`$ | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 | 10 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`u_i`$ | 33 | 142 | 278 | 284 | 155 | 48 | 7 | 0 | 0 | 0 |
| $`v_i`$ | 1 | 9 | 40 | 141 | 266 | 266 | 156 | 55 | 11 | 1 |
So, $`X`$ is of arithmetic $`0pt2.`$
The tangent line of the curve
$$\sigma :k\stackrel{~}{X};s\sigma (s):=(0:0:0:0:s^3:s^2:s:1:0:0:0:0:0)$$
in the point $`\sigma (0)`$ contains $`p_3`$. So, the secant cone $`\mathrm{Sec}_{p_3}(\stackrel{~}{X})`$ โ which must be a line according to Theorem 5.6 (d) โ is just the line $`\mathrm{}`$ which joins $`p_3`$ and $`\sigma (0)`$. The projection of $`\mathrm{}`$ from $`p_3`$ to $`_k^{11}`$ is the point
$$q:=(0:0:0:0:0:0:1:0:0:0:0:0)X.$$
So, in the notation of Theorem 5.6, we have $`\mathrm{Sing}(\nu )=\varrho (\mathrm{}\{p_2\})=\{q\}`$. Now, let $`a:=\frac{u^2}{v^2},b:=\frac{s}{t},c:=\frac{tw}{v^2}`$. An easy calculation shows that there is an isomorphism
$$ฯต:X_{t^3v^2}\stackrel{}{}Y:=\mathrm{Spec}(k[a,ab,b^2,b^3,b^2c,bc,c])$$
with $`ฯต(q)=\underset{ยฏ}{0}`$, where $`X_{t^3v^2}X`$ is the affine open neighborhood of $`q`$ defined by $`t^3v^20`$. It is easy to verify, that $`๐ช_{Y,\underset{ยฏ}{0}}`$ is a $`G`$-ring and hence that $`qX`$ is a $`G`$-point, as predicted by Theorem 5.6.
B) Next, we project the $`3`$-scroll $`\stackrel{~}{X}:=S(2,4,4)_k^{12}`$ of Example 9.1 C) from the point
$$p_4=(0:1:0:0:0:0:0:0:0:0:0:0:0)_k^{12}\stackrel{~}{X}.$$
We get a $`3`$-dimensional variety $`X_k^{11}`$ of degree 10 whose non-vanishing Betti numbers are:
| $`i`$ | 1 | 2 | 3 | 4 | 5 | 6 | 7 | 8 | 9 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`u_i`$ | 34 | 151 | 314 | 364 | 230 | 69 | 7 | 0 | 0 |
| $`v_i`$ | 0 | 0 | 0 | 6 | 35 | 56 | 36 | 10 | 1 |
Now, $`X`$ is of arithmetic $`0pt3=dim(X).`$ For each pair $`(s,t)k^2\{(0,0)\}`$ consider the point
$$\pi (s,t):=(s^2:st:t^2:0:0:0:0:0:0:0:0:0:0)\stackrel{~}{X}.$$
Whenever $`st0`$, the two points $`\pi (s,t)`$ and $`\pi (\sqrt{1}s,\sqrt{1}t)`$ are different and the line joining them contains $`p_4`$ and hence belongs to the secant cone $`\mathrm{Sec}_{p_4}(\stackrel{~}{X})`$. Moreover the tangent line of the curve
$$\tau :k\stackrel{~}{X};s\pi (s,1)$$
in the point $`\tau (0)=\pi (0,1)`$ runs through $`p_4`$ and thus belongs to $`\mathrm{Sec}_{p_4}(\stackrel{~}{X})`$. Altogether this shows (cf Theorem 5.6 (d) ) that $`\mathrm{Sec}_{p_4}(\stackrel{~}{X})`$ coincides with the $`2`$-plane
$$\{(a:b:c:0:0:0:0:0:0:0:0:0:0)|(a:b:c)_k^2\}_k^{12}.$$
Projecting this plane from $`p_4`$ we obtain the line
$$h:=\{(a:c:0:0:0:0:0:0:0:0:0:0)|(a:c)_k^1\}X.$$
So, in the notation of Theorem 5.6 we have $`h=\mathrm{Sing}(\nu )`$. Let $`a:=\frac{t}{s},b:=\frac{s^2v}{u^3},c=\frac{s^2w}{u^3}`$ and let $`X_{s^2u^3}X`$ be the affine open set defined by $`s^2u^20`$. It is easy to verify, that there is an isomorphism $`\phi :X_{s^2u^3}\stackrel{}{}Y:=\mathrm{Spec}(k[a^2,b,ab,c,ac])`$ such that $`P:=(b,ab,c,ac)Y`$ is the generic point of $`\phi (hX_{s^2u^3})`$. An easy calculation shows that $`๐ช_{Y,P}`$ is a $`G`$-ring and hence, that the generic point of $`h`$ in $`X`$ is again a $`G`$-point.
We now present a class of non-normal Del Pezzo varieties. Note that these varieties are in fact arithmetically Gorenstein.
###### Example 9.3.
A) Let $`r4`$ and let $`\stackrel{~}{X}_k^{r+1}`$ be the rational surface scroll $`S(2,r1)`$, hence the variety which is defined by the $`2\times 2`$-minors of the matrix
$$\left(\begin{array}{ccccc}x_0& x_1|x_3& x_4& \mathrm{}& x_r\\ x_1& x_2|x_4& x_5& \mathrm{}& x_{r+1}\end{array}\right).$$
$`\stackrel{~}{X}`$ has the homogeneous coordinate ring
$$B:=k[(s,t)^2u^{r2},(s,t)^{r2}v^2]k[s,t,u,v].$$
Now, let $`\varrho :_k^{r+1}\{p\}_k^r,(x_0:x_1:x_2:\mathrm{}:x_{r+1})(x_0:x_2:x_3:\mathrm{}:x_{r+1})`$ be the projection from the point $`p=(0:1:0:\mathrm{}:0)_k^{r+1}\stackrel{~}{X}`$ and let $`X:=\varrho (\stackrel{~}{X})_k^r.`$ Then $`X`$ is a surface and has the homogeneous coordinate ring
$$A:=k[s^2u^{r2},t^2u^{r2},(s,t)^{r2}v^2]B.$$
As $`B`$ is a birational extension of $`A`$, the morphism $`\nu =\varrho :\stackrel{~}{X}X`$ is birational, so that $`\mathrm{deg}X=\mathrm{deg}\stackrel{~}{X}=r`$ and $`X_k^r`$ is a surface of almost minimal degree. Moreover, as $`\stackrel{~}{X}`$ is smooth, $`\nu =\varrho :\stackrel{~}{X}X`$ is a normalization of $`X`$ and $`\mathrm{Sing}(\nu )=\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})`$ is the non-normal locus of $`X`$.
Similar as in example 9.2 B) we can check that the secant cone of $`\stackrel{~}{X}`$at $`p`$ satisfies
$$\mathrm{Sec}_p(\stackrel{~}{X})=\{(a:b:c:0:\mathrm{}:0)|(a:b:c)_k^2\}$$
and hence is a $`2`$-plane. So, by Theorem 5.6, $`X`$ cannot be of arithmetic $`0pt2=dimX`$ and hence is arithmetically Cohen-Macaulay.
Moreover
$$h:=X\mathrm{Nor}(X)=\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})=\{(a:c:0:\mathrm{}:0)_k^r|(a:c)_k^1\}.$$
So, the non-normal locus $`h`$ of $`X`$ is a line.
Now consider the affine open set $`X_{s^2u^{r2}}X`$ defined by $`s^2u^{r2}0`$ and let $`a:=\frac{s^{r5}v^2t}{u^{r2}}`$ and $`b:=\frac{s^{r4}v^2}{u^{r2}}`$. Then, an easy calculation shows that there is an isomorphism
$$\phi :X_{s^2u^{r2}}\stackrel{}{}Y:=\mathrm{Spec}(k[a,b,\frac{a^2}{b^2}])=\mathrm{Spec}(k[a,b,c]/(cb^2a^2))$$
which maps $`hX_{s^2u^{r2}}`$ to the singular line $`a=b=0`$ of the surface $`Y`$. The pinch point $`\underset{ยฏ}{0}`$ of $`Y`$ can be written as $`\phi (\varrho (\mathrm{}\{p\}))`$, where $`\mathrm{}`$ is the tangent line to $`\stackrel{~}{X}`$ at the point $`(1:0:\mathrm{}:0)`$ which contains $`p`$.
The same arguments apply to the affine open set $`X_{t^2u^{r2}}X`$. This allows to conclude that the open neighborhood $`X_{s^2u^{r2}}X_{t^2u^{r2}}`$ of the singular line $`h`$ of $`X`$ is isomorphic to the blow-up $`\mathrm{Proj}(k[a,b][a^2T,b^2T])`$ of the affine plane $`๐ธ_k^2=\mathrm{Spec}(k[a,b])`$ with respect to the polynomials $`a^2`$ and $`b^2`$.
B) Let $`r5`$ and let $`\stackrel{~}{X}_k^{r+1}`$ be the rational normal 3-scroll $`S(1,1,r3)`$, hence the variety which is defined by the $`2\times 2`$-minors of the matrix
$$\left(\begin{array}{cccc}x_0\left|x_2\right|x_4& x_5& \mathrm{}& x_r\\ x_1\left|x_3\right|x_5& x_6& \mathrm{}& x_{r+1}\end{array}\right).$$
$`\stackrel{~}{X}`$ has the homogeneous coordinate ring
$$B:=k[(s,t)u^{r4},(s,t)v^{r4},(s,t)^{r4}w]k[s,t,u,v,w].$$
Now, let $`\varrho :_k^{r+1}\{p\}_k^r,(x_0:x_1:x_2:\mathrm{}:x_{r+1})(x_0:x_1x_2:x_3:\mathrm{}:x_{r+1})`$ be the projection from the point $`p=(0:1:1:0:\mathrm{}:0)_k^{r+1}\stackrel{~}{X}`$ and let $`X:=\varrho (\stackrel{~}{X})_k^r.`$ Then $`X`$ is of dimension 3 and has the homogeneous coordinate ring
$$A:=k[su^{r4},tu^{r4}sv^{r4},tv^{r4},(s,t)^{r4}w]B.$$
As $`B`$ is a birational extension of $`A`$, the morphism $`\nu =\varrho :\stackrel{~}{X}X`$ is birational, so that $`\mathrm{deg}X=\mathrm{deg}\stackrel{~}{X}=r`$ and $`X_k^r`$ has dimension $`3`$ and is of almost minimal degree. Moreover, as $`\stackrel{~}{X}`$ is smooth, $`\nu =\varrho :\stackrel{~}{X}X`$ is a normalization of $`X`$ and $`\mathrm{Sing}(\nu )=\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})`$ is the non-normal locus of $`X`$.
Similar as in example A) above we can check that the secant cone of $`\stackrel{~}{X}`$at $`p`$ satisfies
$$\mathrm{Sec}_p(\stackrel{~}{X})=\{(a:b:c:d:0:\mathrm{}:0)|(a:b:c:d)_k^3\}$$
and hence is a $`3`$-plane. So, by Theorem 5.6 the variety $`X`$ cannot be of arithmetic $`0pt3=dimX`$ and hence is arithmetically Cohen-Macaulay, that is a non-normal Del Pezzo variety of dimension 3.
Moreover
$$\varrho (\mathrm{Sec}_p(\stackrel{~}{X})\{p\})=\{(a:b:d:0:\mathrm{}:0)_k^r|(a:b:d)_k^2\}.$$
So, the non-normal locus of $`X`$ is a plane, in accordance with Proposition 5.8 and Corollary 6.10.
Finally observe that $`X`$ in 9.3 A) is a divisor on the variety of minimal degree $`\varrho (Z)_k^r`$, where $`Z_k^{r+1}`$ is the variety defined by the $`2\times 2`$ minors of the matrix
$$\left(\begin{array}{cccc}x_3& x_4& \mathrm{}& x_r\\ x_4& x_5& \mathrm{}& x_{r+1}\end{array}\right).$$
In the previous example we have met arithmetically Cohen-Macaulay varieties of almost minimal degree which occur as a subvariety of codimension one on a variety of minimal degree. We now present an example of a normal Del Pezzo variety which does not have this property.
###### Example 9.4.
Let $`X_k^9`$ be the smooth $`6`$-dimensional arithmetically Gorenstein variety of degree $`5`$ defined by the $`4\times 4`$ Pfaffian quadrics $`F_1,F_2,F_3,F_4,F_5`$ of the skew symmetric matrix (cf )
$$M=\left(\begin{array}{cccccc}0& x_0& x_1& x_2& x_3& \\ x_0& 0& x_4& x_5& x_6\\ x_1& x_4& 0& x_7& x_8\\ x_2& x_5& x_7& 0& x_9\\ x_3& x_6& x_8& x_9& 0\end{array}\right).$$
According to the columns of $`M`$ provide a minimal system of generators for the first syzygy module of the homogeneous vanishing ideal $`IS=k[x_0,x_2,\mathrm{},x_9]`$ of $`X`$. Assume now that there is a variety $`W_k^9`$ of minimal degree with $`dimW=7`$ and $`XW`$. So, $`W`$ is arithmetically Cohen-Macaulay and of codimension $`2`$ and by the Theorem of Hilbert-Burch the homogeneous vanishing ideal $`JS`$ of $`W`$ is generated by the three $`2\times 2`$-minors $`G_1,G_2,G_3S_2`$ of a $`2\times 3`$-matrix with linearly independent entries in $`S_1`$ (cf ). So, after an eventual renumbering of the generators $`F_i`$, we may assume that $`G_1,G_2,G_3,F_4,F_5I_2`$ is a minimal system of generators of $`I`$. As $`J`$ admits two independent syzygies
$$\lambda _{i1}G_1+\lambda _{i2}G_2+\lambda _{i3}G_3=0,\lambda _{ij}S_1,i=1,2,j=1,2,3,$$
a minimal system of generators for the first syzygy module of $`I`$ would be given by the matrix of the form
$$N=\left(\begin{array}{ccccc}0& 0& & & \\ 0& 0& & & \\ & & & & \\ & & & & \\ & & & & \end{array}\right)S_1^{5\times 5}.$$
On the other hand there should be a $`k`$-linear transformation which converts $`N`$ into $`M`$ โ a contradiction.
###### Remark 9.5.
A) The variety $`X_k^9`$ of Example 9.4 is normal and Dell Pezzo and hence not a projection of a variety $`\stackrel{~}{X}_k^{10}`$ of minimal degree. The non-existence of the above variety $`W_k^9`$ of minimal degree thus is in accordance with Theorem 7.3. By Remark 7.7 the projection $`X_k^4`$ of the Veronese surface $`F_k^5`$ is not contained in a variety $`Y_k^4`$ of minimal degree either, according to the fact, that $`F`$ is not a scroll. So Remark 7.7 and Example 9.4 illustrate that the hypotheses of Theorem 7.3 cannot be weakened.
B) The examples of this section (with the execption of the last one) are all of relatively big codimension. It turns out, that the structure of varieties of almost minimal degree and small codimension is fairly fixed and cannot vary very much. We study these varieties more extensively in .
### Acknowledgment
We thank the referee for his valuable hints concerning further investigations on birational projections from varieties of minimal and almost minimal degree. In fact, a part of his ideas come up already in an on-going joint investigation of the authors on projective surfaces $`X`$ of degree $`r+1`$ in $`P_k^r.`$ We thank also Euisung Park for his comment concerning Proposition 3.4.
| M. Brodmann | | P. Schenzel |
| --- | --- | --- |
| Institut fรผr Mathematik | | Martin-Luther-Univ. Halle-Wittenberg |
| Universitรคt Zรผrich | | Fachbereich Mathematik und Informatik |
| Winterthurerstrasse 190 | | Von-Seckendorff-Platz 1 |
| CH-8057 Zรผrich, Schwitzerland | | D-06120 Halle (Saale), Germany |
| email: brodmann@math.unizh.ch | | email: schenzel@informatik.uni-halle.de | |
warning/0506/nucl-th0506064.html | ar5iv | text | # The importance of strange mesons in neutron star properties
## I Introduction
In the present work we use the relativistic non-linear Walecka model (NLWM) walecka , at zero temperature ($`T=0`$), with the lowest baryon octet $`\{N,\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }\}`$ in $`\beta `$ equilibrium with the lightest leptons $`\{e^{},\mu ^{}\}`$ and compare the results with the same model plus strange meson fields, $`\sigma ^{}(975)`$ and $`\varphi (1020)`$, which introduce strangeness to the interaction according to alemao and greiner . Strange meson fields, namely the scalar meson field $`\sigma ^{}(975)`$ and the vector meson field $`\varphi (1020)`$, had to be considered in order to reproduce the observed strongly attractive $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction. This formalism applied to compact objects like neutron stars, where the energies are such that allow the appearance of the eight lightest baryons. The motivation for this study lies in our interest to describe the interaction between hadrons taking into account a growing number of effects in order to better describe it. Given the difficulty of making comparisons with experimental data we will, for now limit ourselves to verify if the inclusion of these mesons significantly alters some quantities like pressure and energy density in the equation of state and the bulk properties of compact stars.
## II The Formalism
The lagrangian density of the NLWM with the inclusion of the strange meson sector and leptons for $`\beta `$ equilibrium is:
$$=\underset{B=1}{\overset{8}{}}\overline{\mathrm{\Psi }}_B[\gamma _\mu (i^\mu g_{\omega B}V^\mu g_{\rho B}\stackrel{}{\tau }.\stackrel{}{b}^\mu )(M_Bg_{\sigma B}\sigma )]\mathrm{\Psi }_B$$
$$+\frac{1}{2}\left(_\mu \sigma ^\mu \sigma m_\sigma ^2\sigma ^2\right)\frac{1}{3!}k\sigma ^3\frac{1}{4!}\lambda \sigma ^4\frac{1}{4}\mathrm{\Omega }_{\mu \nu }\mathrm{\Omega }^{\mu \nu }+\frac{1}{2}m_\omega ^2V_\mu V^\mu $$
$$\frac{1}{2}\stackrel{}{B}_{\mu \nu }\stackrel{}{B}^{\mu \nu }+\frac{1}{2}m_\rho ^2\stackrel{}{b}_\mu \stackrel{}{b}^\mu $$
$$+\frac{1}{2}\left(_\mu \sigma ^{}^\mu \sigma ^{}m_\sigma ^{}^2\sigma ^2\right)+\frac{1}{2}m_\varphi ^2\varphi _\mu \varphi ^\mu \frac{1}{4}S_{\mu \nu }S^{\mu \nu }\underset{B}{}g_{\sigma ^{}B}\overline{\mathrm{\Psi }}_B\mathrm{\Psi }_B\sigma ^{}\underset{B}{}g_{\varphi B}\overline{\mathrm{\Psi }}_B\gamma _\mu \mathrm{\Psi }_B\varphi ^\mu $$
$$+\underset{l=1}{\overset{2}{}}\overline{\mathrm{\Psi }}_l\left(i\gamma _\mu ^\mu M_l\right)\mathrm{\Psi }_l,$$
(1)
where $`\mathrm{\Omega }_{\mu \nu }=_\mu V_\nu _\nu V_\mu `$, $`\stackrel{}{B}_{\mu \nu }=_\mu \stackrel{}{b}_\nu _\nu \stackrel{}{b}_\mu g_\rho \left(\stackrel{}{b}_\mu \times \stackrel{}{b}_\nu \right)`$ and $`S_{\mu \nu }=_\mu \varphi _\nu _\nu \varphi _\mu `$, with $`B`$ extending over the eight baryons, $`g_{iB}`$ are the coupling constants of mesons $`i`$$`i=\sigma ,\omega ,\rho `$ with baryon $`B`$, and $`m_i`$ is the mass of meson $`i`$. $`\lambda `$ and $`k`$ are the weighs of the non-linear scalar terms and $`\stackrel{}{\tau }`$ is the isospin operator. At this point it is worth emphasizing that the strange mesons are not supposed to act at low densities, where the strangeness content is zero. Moreover, the non-linear terms are normally corrections added to the main linear contributions and hence the non-linear terms in the strange sector are disregarded in the present work. The constants $`g_{iB}`$ are defined by $`g_{iB}=x_Bg_i`$ where $`x_B=\sqrt{2/3}`$, for hyperons, mos , $`x_B=1`$ for the nucleons, and also $`g_\sigma =8.910`$, $`g_\omega =10.626`$, $`g_\rho =8.208`$, $`g_{\sigma ^{}\mathrm{\Lambda }}=g_{\sigma ^{}\mathrm{\Sigma }}=5.11`$, $`g_{\sigma ^{}\mathrm{\Xi }}=9.38`$, $`g_{\varphi \mathrm{\Lambda }}=g_{\varphi \mathrm{\Sigma }}=4.31,g_{\varphi \mathrm{\Xi }}=8.62`$, $`k=\mathrm{6.426\hspace{0.33em}10}^4,\lambda =5.530`$ according to gle1 and gle2 . The strange mesons interact with hyperons only ($`g_{\sigma ^{}p}=g_{\sigma ^{}n}=g_{\varphi p}=g_{\varphi n}=0`$). The masses of baryons of the octect are: $`M_N=938MeV`$ (nucleons), $`M_\mathrm{\Lambda }=1116MeV`$, $`M_\mathrm{\Sigma }=1193MeV`$, $`M_\mathrm{\Xi }=1318MeV`$ and the meson masses are: $`m_\sigma =512MeV`$, $`m_\omega =738MeV`$, $`m_\rho =770MeV`$, $`m_\sigma ^{}=975MeV`$, $`m_\varphi =1020MeV`$. In order to account for the $`\beta `$ equilibrium in the star the leptons are also included in the lagrangian density of eq. (1) as a non-interacting Fermi gas. The masses of the leptons are $`M_e^{}=0.511MeV`$ and $`M_\mu ^{}=105.66MeV`$.
Applying the Euler-Lagrange equations to (1) and using the mean-field approximation ($`\sigma \sigma =\sigma _0,V_\mu V_\mu =\delta _{\mu 0}V_0`$ and $`\stackrel{}{b}_\mu \stackrel{}{b}_\mu =\delta _{\mu 0}b_\mu ^0\delta _{\mu 0}b_\mu `$), we obtain:
$$\sigma _0=\frac{k}{2m_\sigma ^2}\sigma _0^2\frac{\lambda }{6m_\sigma ^2}\sigma _0^3+\underset{B}{}\frac{g_\sigma }{m_\sigma ^2}x_B\rho _{\sigma B},$$
(2)
$`V_0={\displaystyle \underset{B}{}}{\displaystyle \frac{g_\omega }{m_\omega ^2}}x_B\rho _B,`$ (3)
$`b_0={\displaystyle \underset{B}{}}{\displaystyle \frac{g_\rho }{m_\rho ^2}}x_B\tau _3\rho _B,`$ (4)
$`\sigma _0^{}={\displaystyle \underset{B}{}}{\displaystyle \frac{g_\sigma ^{}}{m_\sigma ^{}^2}}x_B\rho _B,`$ (5)
$`\varphi _0={\displaystyle \underset{B}{}}{\displaystyle \frac{g_\varphi }{m_\varphi ^2}}x_B\rho _B,`$ (6)
where
$$\rho _{\sigma B}=\frac{M_B^{}}{\pi ^2}_0^{K_{F_B}}\frac{p^2dp}{\sqrt{p^2+M_B^{}}},$$
(7)
$$\rho _B=\frac{1}{3\pi ^2}K_{F_B}^3,M_B^{}=M_Bg_{\sigma B}\sigma _0g_{\sigma ^{}B}\sigma _0^{},$$
(8)
and the 0 subscripts added to the fields mean that a mean field approximation, where the meson fields were considered as classical fields was performed.
Through the energy-momentum tensor, we obtain:
$$\epsilon _a=\frac{1}{\pi ^2}\left(\underset{i=B,l}{}_0^{K_{F_i}}p^2๐p\sqrt{p^2+M_i^2}\right)+\frac{m_\omega ^2}{2}V_0^2+\frac{m_\rho ^2}{2}b_0^2+\frac{m_\sigma ^2}{2}\sigma _0^2+\frac{k}{6}\sigma _0^3+\frac{\lambda }{24}\sigma _0^4+\frac{m_\sigma ^{}^2}{2}\sigma _0^2+\frac{m_\varphi ^2}{2}\varphi _0^2,$$
(9)
$$P_a=\frac{1}{3\pi ^2}\left(\underset{i=B,l}{}_0^{K_{F_i}}\frac{p^4dp}{\sqrt{p^2+M_i^2}}\right)+\frac{m_\omega ^2}{2}V_0^2+\frac{m_\rho ^2}{2}b_0^2\frac{m_\sigma ^2}{2}\sigma _0^2\frac{k}{6}\sigma _0^3\frac{\lambda }{24}\sigma _0^4\frac{m_\sigma ^{}^2}{2}\sigma _0^2+\frac{m_\varphi ^2}{2}\varphi _0^2.$$
(10)
In a neutron star, charge neutrality and baryon number must be conserved quantities. Moreover, the conditions of chemical equilibrium hold. In terms of the chemical potentials of the constituent particles, these conditions read:
$`\mu _n=\mu _p+\mu _e^{},\mu _\mu ^{}=\mu _e^{},`$
$`\mu _{\mathrm{\Sigma }^0}=\mu _{\mathrm{\Xi }^0}=\mu _\mathrm{\Lambda }=\mu _n,`$
$`\mu _\mathrm{\Sigma }^{}=\mu _\mathrm{\Xi }^{}=\mu _n+\mu _e^{},`$
$`\mu _{\mathrm{\Sigma }^+}=\mu _p=\mu _n\mu _e^{}.`$ (11)
## III Results
The inclusion of hyperons and strange mesons alters the equations of state and the particle fractions, as can be seen from figures 1 and 2.
From fig. 1 we notice that the inclusion of the hyperons softens the equations of state in comparison with the EOS obtained only with nucleons and leptons. The inclusion of the strange mesons hardens these equations a little at higher energy densities. This indicates that the influence of the strange mesons is significant at higher densities, what can be easily seen in fig. 2, where we notice a difference in the fractions of heavier hyperons, at densities above $`5\rho _0`$, where $`\rho _0`$ is the saturation density of the nuclear matter.
Neutron star profiles can be obtained by solving the Tolman-Oppenheimer-Volkoff (TOV) equations tov , resulting from the exact solution of Einsteinโs general relativity equations in the Schwarzschild metric for spherically symmetric, static stars. Applying the equation of states (9) and (10) in TOV equations results in the star properties shown in table I and figure 3. In table I the profiles of the stars with the maximum gravitational mass and with the maximum radius are shown for two possible EOS: without the strange mesons and with them. In these cases the crust of the stars were not included.
The observed values for the mass of the neutron stars lie between 1.2 to 1.8 $`M_{sun}`$. Our results are in the expected range. From table I and figure 3, one can see that the differences in the star properties with and without strange meson are not very relevant. Nevertheless, the constitution of the stars at large densities are somewhat different. At about four times the saturation density (see figure 2) the inclusion of the strange mesons start playing its role in the constitution of the stars. At this high energy a phase transition to a deconfined phase of quarks or to a system with kaon condensates can already take place. These two possibilities are certainly more important to the properties of neutron stars than a system containing strange mesons. The influence of the inclusion of the strange mesons in protoneutron stars with temperatures around 30 to 40 $`MeV`$ and the their importance when trapped neutrinos are included are under investigation.
## ACKNOWLEDGMENTS
This work was partially supported by CNPq (Brazil). |
warning/0506/nucl-th0506015.html | ar5iv | text | # MKPH-T-05-05 Incoherent pion photoproduction on the deuteron with polarization observables I: Formal expressions
## I Introduction
Photoproduction of pions on light nuclei is an important topic in medium energy nuclear physics. It is motivated by different and complementary aspects. On the one hand one wants to study the elementary reaction on the neutron to which otherwise one has no access. On the other hand one is interested in the influence of a nuclear environment on the elementary production amplitude, and last but not least, one hopes to obtain information on nuclear structure.
Besides the study of unpolarized total and differential cross sections, polarization observables provide very often further insight into details of the underlying reaction mechanisms and possible structure effects. In this case, such observables will serve as additional critical tests or check points for theoretical models. The considerable progress in experimental techniques for studying polarization phenomena has brought into focus also the question, what role polarization effects play in pion photoproduction on nuclei. Of particular interest is photoproduction of pions on the deuteron in view of its simple structure. Indeed, it has been studied quite extensively over the past 50 years (see DaA03a and references therein). While in earlier work mainly total and semi-exclusive differential cross sections of incoherent pion production have been studied, polarization observables were considered more recently, both in experiment LEGS ; A2 as well as in theory. For example, the spin asymmetry of the total cross section with respect to circular photon polarization, which determines the Gerasimov-Drell-Hearn sum rule, was investigated theoretically in LeP96 ; DaA03b ; ArF04 and target asymmetries were considered in LoS00 .
Subsequently, various polarization asymmetries of the semi-exclusive differential cross section $`\stackrel{}{d}(\stackrel{}{\gamma },\pi )NN`$ were studied theoretically in a series of papers Dar04a ; Dar05a ; Dar05b ; Dar05c ; DaS05 . Unfortunately, many of the results presented there are based on incorrect expressions for polarization observables, because the formal expressions for them were taken in analogy from the corresponding expressions of deuteron photodisintegration Are88 . This is in principle possible, since the spin degrees of freedom are the same in both reactions, provided one takes care to check where certain symmetry properties of the reaction amplitude have been used in the derivation of the polarization observables in photodisintegration, because they are not identical in both reactions. This caveat refers in particular to those observables which are related to linearly polarized photons. It appears that this fact was not taken into account so that the results in Dar05a ; Dar05b for them cannot be trusted. But also the results for circularly polarized photons are incorrect, namely the claim in Dar04a , that all of them vanish identically, is wrong. Moreover, this statement is in contradiction to Dar05b , where a non-vanishing differential spin asymmetry for circularly polarized photons is reported, because this asymmetry is proportional to the beam-target asymmetry $`T_{10}^c`$ for circularly polarized photons and a vector polarized deuteron, which means that the latter does not vanish. Thus, it is obvious that the importance of polarization effects requires a more careful and thorough treatment as done in Dar04a ; Dar05a ; Dar05b ; Dar05c ; DaS05 .
With the present work we want to provide a solid basis for the formal expressions of the various polarization observables which determine the differential cross section for incoherent pion production on the deuteron with polarized photons and/or polarized deuterons by deriving the general form of the differential cross section including all possible polarization asymmetries. It complements the work of Blaazer et al. BlB94 , who have formally derived all possible polarization observables for coherent pion photoproduction on the deuteron.
## II Kinematics
As a starting point, we will first consider the kinematics of the photoproduction reaction
$$\gamma (k,\stackrel{}{\epsilon }_\mu )+d(p_d)\pi (q)+N_1(p_1)+N_2(p_2),$$
(1)
where we have defined the notation of the four-momenta of the participating particles. The circular polarization vector of the photon is denoted by $`\stackrel{}{\epsilon }_\mu `$ with $`\mu =\pm 1`$. The following formal developments will not depend on the reference frame, laboratory or center-of-momentum (c.m.) frame. However, in view of our explicit application FiA05 in which the reaction is evaluated in the laboratory frame, we will refer sometimes to this frame for definiteness. We choose as independent variables for the description of the final state the outgoing pion momentum $`\stackrel{}{q}=(q,\theta _q,\varphi _q)`$ and the spherical angles $`\mathrm{\Omega }_p=(\theta _p,\varphi _p)`$ of the relative momentum $`\stackrel{}{p}=(\stackrel{}{p}_1\stackrel{}{p}_2)/2=(p,\mathrm{\Omega }_p)`$ of the two outgoing nucleons. Together with the incoming photon energy $`\omega =k_0`$, the momenta of the outgoing nucleons are fixed, i.e. $`\stackrel{}{p}_{1/2}=(\stackrel{}{k}+\stackrel{}{p}_d\stackrel{}{q})/2\pm \stackrel{}{p}`$. The coordinate system is chosen to have a right-handed orientation with $`z`$-axis along the photon momentum $`\stackrel{}{k}`$. We distinguish in general three planes: (i) the photon plane spanned by the photon momentum and the direction of maximal linear photon polarization, which defines the direction of the $`x`$-axis, (ii) the pion plane, spanned by the photon and pion momenta, which intersects the photon plane along the $`z`$-axis with an angle $`\varphi _q`$, and (iii) the nucleon plane spanned by the momenta of the two outgoing nucleons intersecting the pion plane along the total momentum of the two nucleons. This is illustrated in Fig. 1 for the laboratory frame. In case that the linear photon polarization vanishes, one can choose $`\varphi _q=0`$ and then photon and pion planes coincide.
## III The $`T`$-matrix
All observables are determined by the $`T`$-matrix elements of the electromagnetic pion production current $`\stackrel{}{J}_{\gamma \pi }`$ between the initial deuteron and the final $`\pi NN`$ states. In a general frame, it is given by
$$T_{sm_s,\mu m_d}=^{()}\stackrel{}{p}_1\stackrel{}{p}_2sm_s,\stackrel{}{q}|\stackrel{}{\epsilon }_\mu \stackrel{}{J}_{\gamma \pi }(0)|\stackrel{}{p}_d\mathrm{\hspace{0.17em}1}m_d,$$
(2)
where $`s`$ and $`m_s`$ denote the total spin and its projection on the relative momentum $`\stackrel{}{p}`$ of the outgoing two nucleons, and $`m_d`$ correspondingly the deuteron spin projection on the $`z`$-axis as quantization axis. Furthermore, transverse gauge has been chosen. The knowledge of the specific form of $`\stackrel{}{J}_{\gamma \pi }`$ is not needed for the following formal considerations.
The general form of the $`T`$-matrix after separation of the overall c.m.-motion is given by
$`T_{sm_s\mu m_d}(q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)`$ $`=`$ $`^{()}\stackrel{}{p}sm_s,\stackrel{}{q}|J_{\gamma \pi ,\mu }(\stackrel{}{k})|1m_d`$ (3)
$`=`$ $`\sqrt{2\pi }{\displaystyle \underset{L}{}}i^L\widehat{L}^{()}\stackrel{}{p}sm_s,\stackrel{}{q}|๐ช_\mu ^{\mu L}|1m_d`$
with $`\mu =\pm 1`$ and transverse electric and magnetic multipoles
$`๐ช_M^{\mu L}`$ $`=`$ $`E_M^L+\mu M_M^L.`$ (4)
Furthermore, we use throughout the notation $`\widehat{L}=\sqrt{2L+1}`$. It is convenient to introduce a partial wave decomposition of the final states by
$`{}_{}{}^{()}\stackrel{}{p}sm_s|`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}{\displaystyle \underset{l_pj_pm_p}{}}\widehat{l}_p(l_p0sm_s|j_pm_s)D_{m_s,m_p}^{j_p}(\varphi _p,\theta _p,\varphi _p)^{()}p(l_ps)j_pm_p|,`$ (5)
$`{}_{}{}^{()}\stackrel{}{q}|`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}{\displaystyle \underset{l_qm_q}{}}\widehat{l}_qD_{0,m_q}^{l_q}(\varphi _q,\theta _q,\varphi _q)^{()}ql_qm_q|,`$ (6)
where $`m_p`$ and $`m_q`$ like $`m_d`$ refer to the photon momentum $`\stackrel{}{k}`$ as quantization axis. Here, the rotation matrices $`D_{m^{}m}^j`$ are taken in the convention of Rose Ros57 . Using the multipole decomposition and applying the Wigner-Eckart theorem yields
$`{}_{}{}^{()}p(l_ps)j_pm_p,ql_qm_q|๐ช_M^{\mu L}|1m_d`$ $`=`$ $`{\displaystyle \underset{JM_J}{}}()^{j_pl_q+J}\widehat{J}\left(\begin{array}{ccc}j_p& l_q& J\\ m_p& m_q& M_J\end{array}\right)\left(\begin{array}{ccc}J& L& 1\\ M_J& M& m_d\end{array}\right)`$ (12)
$`pq((l_ps)j_pl_q)J๐ช^{\mu L}1,`$
with the selection rule $`m_p+m_q=M_J=M+m_d`$. Rewriting the angular dependence
$`D_{m_s,m_p}^{j_p}(\varphi _p,\theta _p,\varphi _p)D_{0,m_q}^{l_q}(0,\theta _q,\varphi _q)`$ $`=`$ $`d_{m_s,m_p}^{j_p}(\theta _p)d_{0,m_q}^{l_q}(\theta _q)e^{i((m_pm_s)\varphi _p+m_q\varphi _q)},`$ (13)
and rearranging, using the foregoing selection rule for $`M=\mu `$,
$`(m_pm_s)\varphi _p+m_q\varphi _q`$ $`=`$ $`(m_pm_s)\varphi _{pq}+(\mu +m_dm_s)\varphi _q`$ (14)
with $`\varphi _{pq}=\varphi _p\varphi _q`$, one finds that the $`T`$-matrix can be written as
$`T_{sm_s\mu m_d}(\mathrm{\Omega }_p,\mathrm{\Omega }_q)`$ $`=`$ $`e^{i(\mu +m_dm_s)\varphi _q}t_{sm_s\mu m_d}(\theta _p,\theta _q,\varphi _{pq}),`$ (15)
where the small $`t`$-matrix depends only on $`\theta _p`$, $`\theta _q`$, and the relative azimuthal angle $`\varphi _{pq}`$. Explicitly one has
$`t_{sm_s\mu m_d}(\theta _p,\theta _q,\varphi _{pq})`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2\pi }}}{\displaystyle \underset{Ll_pj_pm_pl_qm_qJJM_J}{}}i^L\widehat{L}\widehat{J}\widehat{l}_q\widehat{l}_p\widehat{j}_p()^{J+l_p+j_ps+m_sl_q}`$ (23)
$`\left(\begin{array}{ccc}l_p& s& j_p\\ 0& m_s& m_s\end{array}\right)\left(\begin{array}{ccc}j_p& l_q& J\\ m_p& m_q& M_J\end{array}\right)\left(\begin{array}{ccc}J& L& 1\\ M_J& \mu & m_d\end{array}\right)`$
$`pq((l_ps)j_pl_q)J๐ช^{\mu L}1d_{m_s,m_p}^{j_p}(\theta _p)d_{0,m_q}^{l_q}(\theta _q)e^{i(m_pm_s)\varphi _{pq}}.`$
Using this explicit form for the small $`t`$-matrix, it is quite straightforward to show that, if parity is conserved, the following symmetry relation holds for the inverted spin projections
$`t_{sm_s\mu m_d}(\theta _p,\theta _q,\varphi _{pq})`$ $`=`$ $`()^{s+m_s+\mu +m_d}t_{sm_s\mu m_d}(\theta _p,\theta _q,\varphi _{pq}).`$ (24)
In the derivation of this relation one has made use of the parity selection rules for the multipole transitions to a final partial wave $`|pq((l_ps)j_pl_q)J`$ with parity $`\pi _{J(l_p,l_q)}=()^{l_p+l_q+1}`$
$`\left\{\begin{array}{ccc}E^L\hfill & \pi _d\pi _{J(l_p,l_q)}()^L=1\hfill & ()^{l_p+l_q+L}=1\hfill \\ M^L\hfill & \pi _d\pi _{J(l_p,l_q)}()^L=1\hfill & ()^{l_p+l_q+L}=1\hfill \end{array}\right\}.`$ (27)
Therefore, invariance under a parity transformation results in the following property of the reduced matrix element
$`()^{l_p+l_q+L}pq((l_ps)j_pl_q)J๐ช^{\mu L}1`$ $`=`$ $`pq((l_ps)j_pl_q)J๐ช^{\mu L}1.`$ (28)
The symmetry property (24) leads to a corresponding relation for the $`T`$-matrix
$`T_{sm_s\mu m_d}(\theta _p,\varphi _p,\theta _q,\varphi _q)`$ $`=`$ $`()^{s+m_s+\mu +m_d}T_{sm_s\mu m_d}(\theta _p,\varphi _p,\theta _q,\varphi _q).`$ (29)
For an uncoupled spin representation, one finds accordingly, using the transformation
$`T_{m_1m_2\mu m_d}(\theta _p,\varphi _p,\theta _q,\varphi _q)`$ $`=`$ $`{\displaystyle \underset{sm_s}{}}({\displaystyle \frac{1}{2}}m_1{\displaystyle \frac{1}{2}}m_2|sm_s)T_{sm_s\mu m_d}(\theta _p,\varphi _p,\theta _q,\varphi _q),`$ (30)
where $`m_j`$ denotes the spin projection of the โjthโ nucleon on the quantization axis, as symmetry relation
$`T_{m_1m_2\mu m_d}(\theta _p,\varphi _p,\theta _q,\varphi _q)`$ $`=`$ $`()^{1+m_1+m_2+\mu +m_d}T_{m_1m_2\mu m_d}(\theta _p,\varphi _p,\theta _q,\varphi _q).`$ (31)
The small $`t`$-matrix elements are the basic quantities which determine differential cross section and asymmetries. The latter are given as ratios of bilinear hermitean forms in terms of the $`t`$-matrix elements (see (51) and (52) below).
## IV The differential cross section including polarization asymmetries
The usual starting point is the general expression for the differential cross section
$`{\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`c(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)tr(T^{}T\rho _i),`$ (32)
where $`T`$ denotes the reaction matrix, $`\rho _i`$ the density matrix for the spin degrees of the initial system. The trace refers to all initial and final state spin degrees of freedom comprising incoming photon, target deuteron, and final nucleons. Furthermore, $`c(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)`$ denotes a kinematic factor which comprises the final state phase space and the incoming flux. In an arbitrary frame one has
$$c(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)=\frac{1}{(2\pi )^5}\frac{E_d}{E_d+p_d}\frac{m_N^2}{4\omega \omega _\pi }\frac{p^{}q^2}{W_{NN}}$$
(33)
with
$$p^{}=\frac{1}{2}\sqrt{W_{NN}^24m_N^2},$$
(34)
as the relative momentum of the final two nucleons in their c.m. sytem, and
$$\omega =k_0,E_d=\sqrt{p_d^2+m_d^2},\omega _\pi =\sqrt{q^2+m_\pi ^2},W_{NN}^2=(\omega +E_d\omega _\pi )^2(\stackrel{}{k}+\stackrel{}{p}_d\stackrel{}{q})^2.$$
(35)
The density matrix $`\rho _i`$ in (32) is a direct product of the density matrices $`\rho ^\gamma `$ of the photon and $`\rho ^d`$ of the deuteron
$$\rho _i=\rho ^\gamma \rho ^d.$$
(36)
The photon density matrix has the form
$$\rho _{\mu \mu ^{}}^\gamma =\frac{1}{2}(\delta _{\mu \mu ^{}}+\stackrel{}{P}^\gamma \stackrel{}{\sigma }_{\mu \mu ^{}})$$
(37)
with respect to circular polarization $`\mu =\pm 1`$. Here, $`|\stackrel{}{P}^\gamma |`$ describes the total degree of polarization, $`P_z^\gamma =P_c^\gamma `$ the degree of circular polarization, and $`P_l^\gamma =\sqrt{(P_x^\gamma )^2+(P_y^\gamma )^2}`$ the degree of linear polarization. By a proper rotation around the photon momentum, one can choose the $`x`$-axis in the direction of maximum linear polarization, i.e., $`P_x^\gamma =P_l^\gamma `$ and $`P_y^\gamma =0`$. Then one has explicitly
$$\rho _{\mu \mu ^{}}^\gamma =(1+\mu P_c^\gamma )\delta _{\mu \mu ^{}}P_l^\gamma \delta _{\mu ,\mu ^{}}e^{2i\mu \varphi _q}.$$
(38)
Furthermore, the deuteron density matrix $`\rho ^d`$ can be expressed in terms of irreducible spin operators $`\tau ^{[I]}`$ with respect to the deuteron spin space
$$\rho _{m_dm_{d}^{}{}_{}{}^{}}^d=\frac{1}{3}\underset{IM}{}()^M\widehat{I}1m_d|\tau _M^{[I]}|1m_d^{}P_{IM}^d,$$
(39)
where $`P_{00}^d=1`$, and $`P_{1M}^d`$ and $`P_{2M}^d`$ describe vector and tensor polarization components of the deuteron, respectively. The spin operators are defined by their reduced matrix elements
$$1\tau ^{[I]}1=\sqrt{3}\widehat{I}\text{for}I=0,1,2.$$
(40)
From now on we will assume that the deuteron density matrix is diagonal with respect to an orientation axis $`\stackrel{}{d}`$ having spherical angles $`(\theta _d,\varphi _d)`$ with respect to the coordinate system associated with the photon plane in the lab frame. Then one has with respect to $`\stackrel{}{d}`$ as quantization axis
$$\rho _{mm^{}}^d=p_m\delta _{mm^{}},$$
(41)
where $`p_m`$ denotes the probability for finding a deuteron spin projection $`m`$ on the orientation axis. With respect to this axis one finds from (39) $`P_{IM}^d(\stackrel{}{d})=P_I^d\delta _{M,0}`$, where the orientation parameters $`P_I^d`$ are related to the probabilities $`\{p_m\}`$ by
$`P_I^d`$ $`=`$ $`\sqrt{3}\widehat{I}{\displaystyle \underset{m}{}}()^{1m}\left(\begin{array}{ccc}1& 1& I\\ m& m& 0\end{array}\right)p_m`$ (42)
$`=`$ $`\delta _{I0}+\sqrt{{\displaystyle \frac{3}{2}}}(p_1p_1)\delta _{I1}+{\displaystyle \frac{1}{\sqrt{2}}}(13p_0)\delta _{I2}.`$
The polarization components in the chosen lab frame are obtained from the $`P_I^d`$ by a rotation, transforming the quantization axis along the orientation axis into the direction of the photon momentum, i.e.
$$P_{IM}^d(\stackrel{}{z})=P_I^de^{iM\varphi _d}d_{M0}^I(\theta _d),$$
(43)
where $`d_{mm^{}}^j`$ denotes a small rotation matrix Ros57 . Thus the deuteron density matrix becomes finally
$$\rho _{m_dm_{d}^{}{}_{}{}^{}}^d=\frac{1}{\sqrt{3}}()^{1m_d}\underset{IM}{}\widehat{I}\left(\begin{array}{ccc}1& 1& I\\ m_d^{}& m_d& M\end{array}\right)P_I^de^{iM\varphi _d}d_{M0}^I(\theta _d).$$
(44)
This means, the deuteron target is characterized by four parameters, namely the vector and tensor polarization parameters $`P_1^d`$ and $`P_2^d`$, respectively, and by the orientation angles $`\theta _d`$ and $`\varphi _d`$. If one chooses the c.m. frame as reference frame, one should note that the deuteron density matrix undergoes no change in the transformation from the lab to the c.m. system, since the boost to the c.m. system is collinear with the deuteron quantization axis Rob74 .
The evaluation of the general expression of the differential cross section in (32) can be done analogously to deuteron photodisintegration as described in detail in Are88 . In fact, one can follow the same steps except for the use of the symmetry relation of Eq. (2) in Are88 which is different in case of pion production (see (24)) because of the additional pion degree of freedom in the final state, in particular its pseudovector character. In terms of the small $`t`$-matrices as defined in (15), one finds, inserting the density matrices of photon and deuteron for the general five-fold differential cross section,
$`{\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu ^{}\mu IM}{}}P_I^de^{iM(\varphi _q\varphi _d)}d_{M0}^I(\theta _d)u_{IM}^{\mu ^{}\mu }\left[(1+\mu P_c^\gamma )\delta _{\mu \mu ^{}}P_l^\gamma \delta _{\mu ,\mu ^{}}e^{2i\mu \varphi _q}\right],`$ (45)
where we have introduced the quantities
$`u_{IM}^{\mu ^{}\mu }(q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`c(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p){\displaystyle \frac{\widehat{I}}{\sqrt{3}}}{\displaystyle \underset{m_dm_d^{}}{}}()^{1m_d}\left(\begin{array}{ccc}1& 1& I\\ m_d^{}& m_d& M\end{array}\right)`$ (46)
$`{\displaystyle \underset{sm_s}{}}t_{sm_s\mu ^{}m_d^{}}^{}(q,\theta _q,\theta _p,\varphi _{pq})t_{sm_s\mu m_d}(q,\theta _q,\theta _p,\varphi _{pq}).`$
It is straighforward to prove that they behave under complex conjugation as
$$u_{IM}^{\mu ^{}\mu }(q,\theta _q,\theta _p,\varphi _{pq})^{}=()^Mu_{IM}^{\mu \mu ^{}}(q,\theta _q,\theta _p,\varphi _{pq}).$$
(47)
Furthermore, with the help of the symmetry in (24) one finds
$$u_{IM}^{\mu ^{}\mu }(q,\theta _q,\theta _p,\varphi _{pq})=()^{I+M+\mu ^{}+\mu }u_{IM}^{\mu ^{}\mu }(q,\theta _q,\theta _p,\varphi _{pq}),$$
(48)
which yields in combination with (47)
$$u_{IM}^{\mu ^{}\mu }(q,\theta _q,\theta _p,\varphi _{pq})=()^{I+\mu ^{}+\mu }u_{IM}^{\mu \mu ^{}}(q,\theta _q,\theta _p,\varphi _{pq})^{}.$$
(49)
This relation is quite useful for a further simplification of the semi-exclusive differential cross section later on.
Separating the polarization parameters of photon ($`P_l^\gamma `$ and $`P_c^\gamma `$) and deuteron ($`P_I^d`$), it is then straightforward to show that the differential cross section can be brought into the form
$`{\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}P_I^d{\displaystyle \underset{M=I}{\overset{I}{}}}e^{iM\varphi _{qd}}d_{M0}^I(\theta _d)[v_{IM}^1+v_{IM}^1`$ (50)
$`+P_c^\gamma (v_{IM}^1v_{IM}^1)+P_l^\gamma (w_{IM}^1e^{2i\varphi _q}+w_{IM}^1e^{2i\varphi _q})],`$
with $`\varphi _{qd}=\varphi _q\varphi _d`$, where we have introduced for convenience the quantities
$`v_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`u_{IM}^{\mu \mu }(q,\theta _q,\theta _p,\varphi _{pq}),`$ (51)
$`w_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`u_{IM}^{\mu \mu }(q,\theta _q,\theta _p,\varphi _{pq}).`$ (52)
According to (47) and (49), they have the following properties under complex conjugation
$`v/w_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq})^{}`$ $`=`$ $`()^Mv/w_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq}),`$ (53)
$`v_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq})^{}`$ $`=`$ $`()^Iv_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq}),`$ (54)
$`w_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq})^{}`$ $`=`$ $`()^Iw_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq}).`$ (55)
From Eq. (53) follows that $`v_{I0}^\mu `$ and $`w_{I0}^\mu `$ are real. The sum over $`M`$ in (50) can be rearranged with the help of the relation (53) and $`d_{M0}^I(\theta _d)=()^Md_{M0}^I(\theta _d)`$
$`{\displaystyle \underset{M=I}{\overset{I}{}}}e^{iM\varphi _{qd}}d_{M0}^I(\theta _d)(v_{IM}^1\pm v_{IM}^1)`$ $`=`$ $`{\displaystyle \underset{M=0}{\overset{I}{}}}{\displaystyle \frac{d_{M0}^I(\theta _d)}{1+\delta _{M0}}}\left(e^{iM\varphi _{qd}}(v_{IM}^1\pm v_{IM}^1)+e^{iM\varphi _{qd}}()^M(v_{IM}^1\pm v_{IM}^1)\right)`$ (56)
$`=`$ $`{\displaystyle \underset{M=0}{\overset{I}{}}}{\displaystyle \frac{d_{M0}^I(\theta _d)}{1+\delta _{M0}}}\left(e^{iM\varphi _{qd}}(v_{IM}^1\pm v_{IM}^1)+\text{c.c.}\right),`$
and furthermore with $`\psi _M=M\varphi _{qd}2\varphi _q`$
$`{\displaystyle \underset{M=I}{\overset{I}{}}}e^{iM\varphi _{qd}}d_{M0}^I(\theta _d)(w_{IM}^1e^{2i\varphi _q}+w_{IM}^1e^{2i\varphi _q})`$ $`=`$ $`{\displaystyle \underset{M=I}{\overset{I}{}}}d_{M0}^I(\theta _d)\left(e^{i\psi _M}w_{IM}^1+e^{i\psi _M}()^Mw_{IM}^1\right)`$ (57)
$`=`$ $`{\displaystyle \underset{M=I}{\overset{I}{}}}d_{M0}^I(\theta _d)\left(e^{i\psi _M}w_{IM}^1+\text{c.c.}\right).`$
This then yields for the differential cross section
$`{\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`{\displaystyle \underset{I}{}}P_I^d\{{\displaystyle \underset{M=0}{\overset{I}{}}}{\displaystyle \frac{1}{1+\delta _{M0}}}d_{M0}^I(\theta _d)\mathrm{}e[e^{iM\varphi _{qd}}(v_{IM}^++P_c^\gamma v_{IM}^{})]`$ (58)
$`+P_l^\gamma {\displaystyle \underset{M=I}{\overset{I}{}}}d_{M0}^I(\theta _d)\mathrm{}e[e^{i\psi _M}w_{IM}^1]\},`$
where we have defined
$$v_{IM}^\pm =v_{IM}^1\pm v_{IM}^1.$$
(59)
Now, introducing various beam, target and beam-target asymmetries by
$`\tau _{IM}^{0/c}(q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`{\displaystyle \frac{1}{1+\delta _{M0}}}\mathrm{}ev_{IM}^\pm (q,\theta _q,\theta _p,\varphi _{pq}),M0,`$ (60)
$`\sigma _{IM}^{0/c}(q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`\mathrm{}mv_{IM}^\pm (q,\theta _q,\theta _p,\varphi _{pq}),M>0,`$ (61)
$`\tau _{IM}^l(q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`\mathrm{}ew_{IM}^1(q,\theta _q,\theta _p,\varphi _{pq}),`$ (62)
$`\sigma _{IM}^l(q,\theta _q,\theta _p,\varphi _{pq})`$ $`=`$ $`\mathrm{}mw_{IM}^1(q,\theta _q,\theta _p,\varphi _{pq}),M0,`$ (63)
where we took into account that $`v_{I0}^\mu `$ and $`w_{I0}^\mu `$ are real, one obtains as final expression for the general five-fold differential cross section with beam and target polarization
$`{\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`{\displaystyle \underset{I}{}}P_I^d\{{\displaystyle \underset{M=0}{\overset{I}{}}}d_{M0}^I(\theta _d)[\tau _{IM}^0\mathrm{cos}(M\varphi _{qd})+\sigma _{IM}^0\mathrm{sin}(M\varphi _{qd})`$ (64)
$`+P_c^\gamma (\tau _{IM}^c\mathrm{cos}(M\varphi _{qd})+\sigma _{IM}^c\mathrm{sin}(M\varphi _{qd}))]`$
$`+P_l^\gamma {\displaystyle \underset{M=I}{\overset{I}{}}}d_{M0}^I(\theta _d)[\tau _{IM}^l\mathrm{cos}\psi _M+\sigma _{IM}^l\mathrm{sin}\psi _M]\}.`$
This constitutes our central result.
We will now turn to the semi-exclusive reaction $`\stackrel{}{d}(\stackrel{}{\gamma },\pi )NN`$ where only the produced pion is detected, which means integration of the five-fold differential cross section $`d^5\sigma /dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p`$ over $`\mathrm{\Omega }_p`$. The resulting cross section will then be governed by the integrated asymmetries $`๐\mathrm{\Omega }_p\tau _{IM}^\alpha `$ and $`๐\mathrm{\Omega }_p\sigma _{IM}^\alpha `$ ($`\alpha \{0,c,l\}`$), of which quite a few will vanish, either $`๐\mathrm{\Omega }_p\tau _{IM}^\alpha `$ or $`๐\mathrm{\Omega }_p\sigma _{IM}^\alpha `$. To show this, we first introduce the quantities
$`W_{IM}(q,\theta _q)`$ $`=`$ $`{\displaystyle ๐\mathrm{\Omega }_pw_{IM}^1(q,\theta _q,\theta _p,\varphi _{pq})}`$ (65)
$`=`$ $`{\displaystyle \frac{\widehat{I}}{\sqrt{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{m_dm_d^{}}{}()^{1m_d}\left(\begin{array}{ccc}1& 1& I\\ m_d^{}& m_d& M\end{array}\right)}`$
$`{\displaystyle \underset{sm_s}{}}t_{sm_s1m_d^{}}^{}(q,\theta _q,\theta _p,\varphi _{pq})t_{sm_s1m_d}(q,\theta _q,\theta _p,\varphi _{pq}),`$
$`V_{IM}^\pm (q,\theta _q)`$ $`=`$ $`V_{IM}^1(q,\theta _q)\pm V_{IM}^1(q,\theta _q),`$ (66)
with
$`V_{IM}^\mu (q,\theta _q)`$ $`=`$ $`{\displaystyle ๐\mathrm{\Omega }_pv_{IM}^\mu (q,\theta _q,\theta _p,\varphi _{pq})}`$ (67)
$`=`$ $`{\displaystyle \frac{\widehat{I}}{\sqrt{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{m_dm_d^{}}{}()^{1m_d}\left(\begin{array}{ccc}1& 1& I\\ m_d^{}& m_d& M\end{array}\right)}`$
$`{\displaystyle \underset{sm_s}{}}t_{sm_s\mu m_d^{}}^{}(q,\theta _q,\theta _p,\varphi _{pq})t_{sm_s\mu m_d}(q,\theta _q,\theta _p,\varphi _{pq}).`$
Using now the property (54), one finds with the help of
$$_0^{2\pi }๐\varphi _pf(\varphi _{pq})=_0^{2\pi }๐\varphi _pf(\varphi _{pq})$$
(68)
for a periodic function $`f(\varphi _{pq}+2\pi )=f(\varphi _{pq})`$, the relation
$$V_{IM}^1(q,\theta _q)=๐\mathrm{\Omega }_pv_{IM}^1(q,\theta _q,\theta _p,\varphi _{pq})=()^I๐\mathrm{\Omega }_pv_{IM}^1(q,\theta _q,\theta _p,\varphi _{pq})^{}=()^IV_{IM}^1(q,\theta _q)^{},$$
(69)
and thus
$$V_{IM}^\pm (q,\theta _q)=V_{IM}^1(q,\theta _q)\pm ()^IV_{IM}^1(q,\theta _q)^{}.$$
(70)
Correspondingly, using (55) one obtains
$$W_{IM}(q,\theta _q)^{}=()^I๐\mathrm{\Omega }_pw_{IM}^1(q,\theta _q,\theta _p,\varphi _{pq})=()^IW_{IM}(q,\theta _q).$$
(71)
From the two foregoing equations we can conclude that $`V_{IM}^+`$ and $`W_{IM}`$ are real for $`I=0`$ and 2 and imaginary for $`I=1`$, whereas $`V_{IM}^{}`$ is imaginary for $`I=0`$ and 2 and real for $`I=1`$. Therefore, according to (60) through (63) the following integrated asymmetries vanish
$`{\displaystyle ๐\mathrm{\Omega }_p\tau _{IM}^\alpha }`$ $`=`$ $`0\text{ for }\left\{\begin{array}{cc}\alpha \{0,l\},\hfill & \text{and}I=1\hfill \\ \alpha \{c\},\hfill & \text{and}I=0,2\hfill \end{array}\right\},`$ (74)
$`{\displaystyle ๐\mathrm{\Omega }_p\sigma _{IM}^\alpha }`$ $`=`$ $`0\text{ for }\left\{\begin{array}{cc}\alpha \{0,l\},\hfill & \text{and}I=0,2\hfill \\ \alpha \{c\},\hfill & \text{and}I=1\hfill \end{array}\right\}.`$ (77)
Instead of using these results for deriving from (64) the three-fold semi-exclusive differential cross section, we prefer to start from the expression in (58), and obtain
$`{\displaystyle \frac{d^3\sigma }{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \underset{I}{}}P_I^d\left\{{\displaystyle \underset{M=0}{\overset{I}{}}}{\displaystyle \frac{1}{1+\delta _{M0}}}d_{M0}^I(\theta _d)\mathrm{}e[e^{iM\varphi _{qd}}(V_{IM}^++P_c^\gamma V_{IM}^{})]+P_l^\gamma {\displaystyle \underset{M=I}{\overset{I}{}}}d_{M0}^I(\theta _d)\mathrm{}e[e^{i\psi _M}W_{IM}]\right\}.`$ (78)
This expression can be simplified using the fact that $`i^{\delta _{I1}}W_{IM}`$, $`i^{\delta _{I1}}V_{IM}^+`$ and $`i^{1\delta _{I1}}V_{IM}^{}`$ are real according to (70) and (71). The latter two quantities can be written as
$`i^{\delta _{I1}}V_{IM}^+`$ $`=`$ $`2\mathrm{}e(i^{\delta _{I1}}V_{IM}^1),`$ (79)
$`i^{1\delta _{I1}}V_{IM}^{}`$ $`=`$ $`2\mathrm{}e(i^{1\delta _{I1}}V_{IM}^1)=2\mathrm{}m(i^{\delta _{I1}}V_{IM}^1).`$ (80)
Using now
$`\mathrm{}e[e^{iM\varphi _{qd}}V_{IM}^+]`$ $`=`$ $`\mathrm{}e[e^{i(M\varphi _{qd}\delta _{I1}\pi /2)}i^{\delta _{I1}}V_{IM}^+]=2\mathrm{}e(i^{\delta _{I1}}V_{IM}^1)\mathrm{cos}[M\varphi _{qd}\delta _{I1}\pi /2],`$ (81)
$`\mathrm{}e[e^{iM\varphi _{qd}}V_{IM}^{}]`$ $`=`$ $`\mathrm{}e[{\displaystyle \frac{1}{i}}e^{i(M\varphi _{qd}+\delta _{I1}\pi /2)}i^{1\delta _{I1}}V_{IM}^{}]=2\mathrm{}m(i^{\delta _{I1}}V_{IM}^1)\mathrm{sin}[M\varphi _{qd}+\delta _{I1}\pi /2],`$ (82)
$`\mathrm{}e[e^{i\psi _M}W_{IM}]`$ $`=`$ $`\mathrm{}e[e^{i(\psi _M\delta _{I1}\pi /2)}i^{\delta _{I1}}W_{IM}]=i^{\delta _{I1}}W_{IM}\mathrm{cos}[\psi _M\delta _{I1}\pi /2],`$ (83)
we find as final form for the three-fold semi-exclusive differential cross section
$`{\displaystyle \frac{d^3\sigma }{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}[1+P_l^\gamma \{\stackrel{~}{\mathrm{\Sigma }}^l\mathrm{cos}2\varphi _q+{\displaystyle \underset{I=1}{\overset{2}{}}}P_I^d{\displaystyle \underset{M=I}{\overset{I}{}}}\stackrel{~}{T}_{IM}^l\mathrm{cos}[\psi _M\delta _{I1}\pi /2]d_{M0}^I(\theta _d)\}`$ (84)
$`+{\displaystyle \underset{I=1}{\overset{2}{}}}P_I^d{\displaystyle \underset{M=0}{\overset{I}{}}}(\stackrel{~}{T}_{IM}^0\mathrm{cos}[M\varphi _{qd}\delta _{I1}\pi /2]+P_c^\gamma \stackrel{~}{T}_{IM}^c\mathrm{sin}[M\varphi _{qd}+\delta _{I1}\pi /2])d_{M0}^I(\theta _d)].`$
Here the unpolarized cross section and the asymmetries are given by
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`V_{00}^1(q,\theta _q),`$ (85)
$`\stackrel{~}{\mathrm{\Sigma }}^l(q,\theta _q){\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`W_{00}(q,\theta _q),`$ (86)
$`\stackrel{~}{T}_{IM}^0(q,\theta _q){\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`(2\delta _{M0})\mathrm{}e[i^{\delta _{I1}}V_{IM}^1(q,\theta _q)],\text{for }0MI,`$ (87)
$`\stackrel{~}{T}_{IM}^c(q,\theta _q){\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`(2\delta _{M0})\mathrm{}m[i^{\delta _{I1}}V_{IM}^1(q,\theta _q)],\text{for }0MI,`$ (88)
$`\stackrel{~}{T}_{IM}^l(q,\theta _q){\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`i^{\delta _{I1}}W_{IM}(q,\theta _q),\text{for }IMI.`$ (89)
Because $`V_{I0}^1`$ is real according to (53), the asymmetries $`\stackrel{~}{T}_{10}`$ and $`\stackrel{~}{T}_{20}^c`$ vanish identically. We would like to point out that in forward and backward pion emission, i.e. for $`\theta _q=0`$ and $`\pi `$, the following asymmetries have to vanish
$$\stackrel{~}{\mathrm{\Sigma }}^l=0,\stackrel{~}{T}_{IM}^{0,c}=0\text{for}M0,\text{and}T_{IM}^l=0\text{for}M2,$$
(90)
because in that case the differential cross section cannot depend on $`\varphi _q`$, since at $`\theta _q=0`$ or $`\pi `$ the azimuthal angle $`\varphi _q`$ is undefined or arbitrary. This feature can also be shown by straightforward evaluation of $`V_{IM}^\mu `$ and $`W_{IM}`$ using the explicit representation of the $`t`$-matrix in (23). One finds
$$V_{IM}^\mu (q,\theta _q=0/\pi ,\theta _p,\varphi _{pq})=0\text{for }M0\text{and}W_{IM}(q,\theta _q=0/\pi ,\theta _p,\varphi _{pq})=0\text{for }M2.$$
(91)
The authors of DaS05 were not aware of this general kinematic property because they evaluate the asymmetries numerically for $`\theta _q=0`$ and $`\pi `$ and find that the obtained values are of the order of $`10^3`$. They conclude in the case of $`T_{11}`$ that it vanishes there but point out that $`\mathrm{\Sigma }^l`$ does not vanish. For completeness and also in view of the numerous errors in Dar04a ; Dar05a ; Dar05b ; DaS05 , we list in the appendix A the explicit expressions of the asymmetries in terms of the $`t`$-matrix elements.
In case that only the direction of the outgoing pion is measured and not its momentum, the corresponding differential cross section $`d^2\sigma /d\mathrm{\Omega }_q`$ is given by an expression formally analogous to (84) where only the above asymmetries are integrated over the pion momentum, i.e., by the replacements
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $``$ $`{\displaystyle \frac{d^2\sigma _0}{d\mathrm{\Omega }_q}}={\displaystyle _{q_{min}(\theta _q)}^{q_{max}(\theta _q)}}๐q{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}},`$ (92)
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{\mathrm{\Sigma }}^l(q,\theta _q)`$ $``$ $`{\displaystyle \frac{d^2\sigma _0}{d\mathrm{\Omega }_q}}\mathrm{\Sigma }^l(\theta _q)={\displaystyle _{q_{min}(\theta _q)}^{q_{max}(\theta _q)}}๐q{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{\mathrm{\Sigma }}^l(q,\theta _q),`$ (93)
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{T}_{IM}^\alpha (q,\theta _q)`$ $``$ $`{\displaystyle \frac{d^2\sigma _0}{d\mathrm{\Omega }_q}}T_{IM}^\alpha (\theta _q)={\displaystyle _{q_{min}(\theta _q)}^{q_{max}(\theta _q)}}๐q{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{T}_{IM}^\alpha (q,\theta _q),\alpha \{0,l,c\}.`$ (94)
The upper and lower integration limits are given by
$`q_{max}(\theta _q)`$ $`=`$ $`{\displaystyle \frac{1}{2b}}\left(a\omega \mathrm{cos}\theta _q+E_{\gamma d}\sqrt{a^24bm_\pi ^2}\right),`$ (95)
$`q_{min}(\theta _q)`$ $`=`$ $`\mathrm{max}\{0,{\displaystyle \frac{1}{2b}}\left(a\omega \mathrm{cos}\theta _qE_{\gamma d}\sqrt{a^24bm_\pi ^2}\right)\},`$ (96)
where
$`a`$ $`=`$ $`W_{\gamma d}^2+m_\pi ^24m_N^2,`$ (97)
$`b`$ $`=`$ $`W_{\gamma d}^2+\omega ^2\mathrm{sin}^2\theta _q,`$ (98)
$`W_{\gamma d}^2`$ $`=`$ $`m_d(m_d+2\omega ),`$ (99)
$`E_{\gamma d}`$ $`=`$ $`m_d+\omega .`$ (100)
The general total cross section is obtained from (84) by integrating over $`q`$ and $`\mathrm{\Omega }_q`$ resulting in
$$\sigma (P_l^\gamma ,P_c^\gamma ,P_1^d,P_2^d)=\sigma _0\left[1+P_2^d\overline{T}_{20}^{\mathrm{\hspace{0.17em}0}}\frac{1}{2}(3\mathrm{cos}^2\theta _d1)+P_c^\gamma P_1^d\overline{T}_{10}^c\mathrm{cos}\theta _d+P_l^\gamma P_2^d\overline{T}_{22}^l\mathrm{cos}(2\varphi _d)\frac{\sqrt{6}}{4}\mathrm{sin}^2\theta _d\right],$$
(101)
where the unpolarized total cross section and the corresponding asymmetries are given by
$`\sigma _0`$ $`=`$ $`{\displaystyle ๐\mathrm{\Omega }_q_{q_{min}(\theta _q)}^{q_{max}(\theta _q)}๐q\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}},`$ (102)
$`\sigma _0\overline{T}_{IM}^\alpha `$ $`=`$ $`{\displaystyle ๐\mathrm{\Omega }_q_{q_{min}(\theta _q)}^{q_{max}(\theta _q)}๐q\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}\stackrel{~}{T}_{IM}^\alpha },`$ (103)
with $`\alpha \{0,l,c\}`$.
Finally, we would like to point out that for coherent photoproduction of $`\pi ^0`$ on the deuteron formally the same expression as in (84) holds with unpolarized differential cross section and asymmetries $`\mathrm{\Sigma }^l(\theta _q)`$, $`T_{IM}(\theta _q)`$, and $`T_{IM}^{c/l}(\theta _q)`$, which are defined in analogy to (85) through (89) with the replacements
$`V_{IM}^1`$ $``$ $`c(\omega ,\mathrm{\Omega }_q){\displaystyle \frac{\widehat{I}}{\sqrt{3}}}{\displaystyle \underset{m_dm_d^{}}{}}()^{1m_d}\left(\begin{array}{ccc}1& 1& I\\ m_d^{}& m_d& M\end{array}\right){\displaystyle \underset{m_d^{\prime \prime }}{}}t_{m_d^{\prime \prime }1m_d^{}}^{}(\theta _q)t_{m_d^{\prime \prime }1m_d}(\theta _q),`$ (104)
$`W_{IM}`$ $``$ $`c(\omega ,\mathrm{\Omega }_q){\displaystyle \frac{\widehat{I}}{\sqrt{3}}}{\displaystyle \underset{m_dm_d^{}}{}}()^{1m_d}\left(\begin{array}{ccc}1& 1& I\\ m_d^{}& m_d& M\end{array}\right){\displaystyle \underset{m_d^{\prime \prime }}{}}t_{m_d^{\prime \prime }1m_d^{}}^{}(\theta _q)t_{m_d^{\prime \prime }1m_d}(\theta _q).`$ (105)
Here, $`c(\omega ,\mathrm{\Omega }_q)`$ denotes a kinematic factor. A complete listing of all polarization observables including recoil polarization of the final deuteron can be found in BlB94 .
## V Conclusions
In this work we have derived formal expressions for the differential cross section of incoherent pion photoproduction on the deuteron including various polarization asymmetries with respect to polarized photons and deuterons. Obviously, these expressions are generally valid for pseudoscalar meson production. We did not consider polarization analysis of the final state, i.e. spin analysis of one or both outgoing nucleons. In this case one has to evaluate instead of (32)
$`P_\alpha (j){\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`c(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)tr(T^{}\sigma _\alpha (j)T\rho _i),`$ (106)
for the polarization of the โjthโ outgoing nucleon, or
$`P_{\alpha _1\alpha _2}{\displaystyle \frac{d^5\sigma }{dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p}}`$ $`=`$ $`c(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)tr(T^{}\sigma _{\alpha _1}(1)\sigma _{\alpha _2}(2)T\rho _i),`$ (107)
for the polarization of both outgoing nucleons. For the evaluation of these expressions one can proceed straightforwardly as has been done in Are88 . In a subsequent paper FiA05 , we will investigate the influence of $`NN`$\- and $`\pi N`$-rescattering on the various asymmetries of the semi-exclusive differential cross section of incoherent pion photoproduction on the deuteron.
###### Acknowledgements.
We would like to thank Michael Schwamb for interesting discussions and a careful reading of the manuscript. This work was supported by the Deutsche Forschungsgemeinschaft (SFB 443). *
## Appendix A Explicit expressions for the various polarization asymmetries
We list here the explicit hermitean, bilinear forms in terms of the $`t`$-matrix elements for cross section and the various asymmetries:
1. The semi-exclusive differential cross section
$$\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}=\frac{1}{3}๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{sm_sm_d}{}|t_{sm_s1m_d}|^2.$$
(A1)
2. The photon asymmetry for linearly polarized photons and unpolarized deuterons
$$\stackrel{~}{\mathrm{\Sigma }}^l\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}=\frac{1}{3}๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{sm_sm_d}{}t_{sm_s1m_d}^{}t_{sm_s1m_d}.$$
(A2)
3. The target asymmetry for vector polarized deuterons and unpolarized photons
$$\stackrel{~}{T}_{11}^0\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}=\sqrt{\frac{2}{3}}๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}m\underset{sm_s}{}(t_{sm_s11}^{}t_{sm_s10}+t_{sm_s10}^{}t_{sm_s11}).$$
(A3)
4. The target asymmetries for tensor polarized deuterons and unpolarized photons
$`\stackrel{~}{T}_{20}^0{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{1}{3\sqrt{2}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{sm_s}{}(|t_{sm_s11}|^2+|t_{sm_s11}|^22|t_{sm_s10}|^2)},`$ (A4)
$`\stackrel{~}{T}_{21}^0{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}e\underset{sm_s}{}(t_{sm_s11}^{}t_{sm_s10}t_{sm_s10}^{}t_{sm_s11})},`$ (A5)
$`\stackrel{~}{T}_{22}^0{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}e\underset{sm_s}{}t_{sm_s11}^{}t_{sm_s11}}.`$ (A6)
5. The beam-target asymmetries for circularly polarized photons and vector polarized deuterons
$`\stackrel{~}{T}_{10}^c{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{sm_s}{}(|t_{sm_s11}|^2|t_{sm_s11}|^2)},`$ (A7)
$`\stackrel{~}{T}_{11}^c{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}e\underset{sm_s}{}(t_{sm_s11}^{}t_{sm_s10}+t_{sm_s10}^{}t_{sm_s11})}.`$ (A8)
6. The beam-target asymmetries for circularly polarized photons and tensor polarized deuterons
$`\stackrel{~}{T}_{21}^c{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}m\underset{sm_s}{}(t_{sm_s10}^{}t_{sm_s11}t_{sm_s11}^{}t_{sm_s10})},`$ (A9)
$`\stackrel{~}{T}_{22}^c{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}m\underset{sm_s}{}t_{sm_s11}^{}t_{sm_s11}}.`$ (A10)
7. The beam-target asymmetries for linearly polarized photons and vector polarized deuterons
$`\stackrel{~}{T}_{10}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}m\underset{sm_s}{}(t_{sm_s11}^{}t_{sm_s11})},`$ (A11)
$`\stackrel{~}{T}_{11}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}m\underset{sm_s}{}(t_{sm_s11}^{}t_{sm_s10})},`$ (A12)
$`\stackrel{~}{T}_{11}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}m\underset{sm_s}{}(t_{sm_s11}^{}t_{sm_s10})}.`$ (A13)
8. The beam-target asymmetries for linearly polarized photons and tensor polarized deuterons
$`\stackrel{~}{T}_{20}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{3}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}e\underset{sm_s}{}(t_{sm_s10}^{}t_{sm_s10}t_{sm_s11}^{}t_{sm_s11})},`$ (A14)
$`\stackrel{~}{T}_{21}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}e\underset{sm_s}{}(t_{sm_s10}^{}t_{sm_s11})},`$ (A15)
$`\stackrel{~}{T}_{21}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\mathrm{}e\underset{sm_s}{}(t_{sm_s10}^{}t_{sm_s11})},`$ (A16)
$`\stackrel{~}{T}_{22}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{sm_s}{}t_{sm_s11}^{}t_{sm_s11}},`$ (A17)
$`\stackrel{~}{T}_{22}^l{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}{\displaystyle ๐\mathrm{\Omega }_pc(\omega ,q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)\underset{sm_s}{}t_{sm_s11}^{}t_{sm_s11}}.`$ (A18) |
warning/0506/hep-ph0506029.html | ar5iv | text | # General soft terms from Supergravity including D-terms
## 1 Introduction
A classical way of communicating supersymmetry breaking to the visible sector is through gravitational interactions. In supergravity$`^\mathrm{?}`$, the hidden sector scalar potential is assumed to have a minimum, preferably generated dynamically, leading to a vacuum expectation value (vev) for at least one of the auxiliary fields. Tree level gravitational interactions then communicate this breaking to the visible sector generating soft terms in the global limit. General expressions for these soft terms can then be derived in terms of these auxiliary fields as has been pointed out long ago by $`^\mathrm{?}`$. By their very nature, such general expressions can be applied to study the soft terms in several classes of models such as supergravity lagrangians derived from superstring theories$`^\mathrm{?}`$.
While the existing expressions have been extremely useful, they could be considered in a certain way as incomplete as they have been concentrating solely on the $`F`$ type supersymmetry breaking terms. It is well known that there could be $`D`$ type contributions too $`^\mathrm{?}`$, that can arise for example in models based on anomalous $`U(1)`$ symmetries $`^\mathrm{?}`$. Furthermore, in effective lagrangians from the Type II orientifolds with intersecting D-branes, one can expect such D-term contributions to be naturally present. There is also a second motivation. Recently, influenced by the multivacuum structure of string theory$`^\mathrm{?}`$ as possible new view of cosmological constant problem, a new proposal has been put forward by the authors of Ref.$`^\mathrm{?}`$. Here, it is proposed that the fermionic superpartners stay close to the weak scale, whereas the scalar superpartners can be present at scales as high as $`10^9`$ GeV. It would be very difficult to achieve this kind of splitting between superpartners in supergravity models with only $`F`$ type supersymmetry breaking. Given these motivations, we present here the results obtained in $`^\mathrm{?}`$ for the general expressions for soft terms in presence of non-zero D-term contributions and study a few applications for them. Particularly, we sketch a model where such large D-terms can be utilised in realising split supersymmetry and address the issue of moduli stabilisation, of particular relevance for any scenario of supersymmetry breaking.
## 2 General Expressions Including D-breaking
Let us now proceed to generalise the analysis in the literature by including abelian gauge groups $`_AU(1)_A`$ and the corresponding D-type contributions to the SUSY breaking. The scalar potential now takes the form:
$$V=e^G(G^MG_M3)+\frac{1}{2}\underset{A}{}g_A^2D_A^2,$$
(1)
where the auxiliary $`F`$ fields are given by $`G_M=\frac{G}{z^M}`$ and $`z`$ represents the scalar part of a chiral superfield. The index $`M`$ runs over all the chiral superfields present, matter as well as hidden sector and/or moduli fields. At the minimum, the hidden sector auxiliary fields attain a vev breaking supersymmetry spontaneously. The D-terms are given by
$$D_A=z^IX_I^A\frac{K}{z^I}+\xi _A=\overline{z}^{\overline{I}}X_I^A\frac{K}{\overline{z}^{\overline{I}}}+\xi _A;\xi _A\eta _A^\alpha _\alpha K,$$
(2)
where $`X_I^A`$ represents the $`U(1)_A`$ charges of the field $`\varphi ^I`$ and $`\xi _A`$ denotes the Fayet-Iliopoulos terms for the abelian $`U(1)`$ factors. Note that the equality between the first two terms is a straightforward consequence of the gauge invariance of the Kรคhler potential. We consider the Fayet-Iliopoulos terms to be moduli dependent and we will not explicitly discuss here the various possible mechanisms of moduli stabilisation. The conditions of the cancellation of the cosmological constant and the requirement of existence of a minimum gives
$`<e^G(G^MG_M3)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2>=0,`$ (3)
$`<e^G(G^M_KG_M+G_K)+{\displaystyle \underset{A}{}}g_A^2D_A(_KD_A{\displaystyle \frac{1}{2}}G_KD_A)>=0,`$ (4)
where $``$ denotes the covariant derivative on the Kรคhler manifold. The scalar soft spectrum is defined as :
$`m_{I\overline{J}}^2=<_I_{\overline{J}}V>=<_I_{\overline{J}}V>,`$ (5)
$`m_{IJ}^2=<_I_JV>=<_I_JV>,`$ (6)
$`A_{IJK}=<_I_J_KV>.`$ (7)
We will further distinguish the visible sector (matter) fields from those of hidden sector fields $`T^\alpha `$ (and later on flavon fields), by requiring $`G^i=0,\mathrm{\Phi }^i=0`$, with $`\mathrm{\Phi }`$ representing the scalar part of a matter field. Using these we recover in the absence of D-term contributions the standard form$`^{\mathrm{?},\mathrm{?}}`$ for the soft scalar masses. Given the form of D-terms above (2), we have in the vacuum, after setting the matter fields vevs to zero
$`_jD_A=\overline{v}_{\overline{\beta }}X_{\overline{\beta }}^AK_{\overline{j}\beta }+\eta _A^{\overline{\alpha }}K_{j\overline{\alpha }}=0,_i_jD_A=0,`$
$`_i_{\overline{j}}D_A`$ $`=`$ $`K_{i\overline{j}}X_i^A+(\overline{v}^{\overline{l}}X_A^{\overline{l}}_{\overline{l}}+\eta _A^{\overline{\alpha }}_{\overline{\alpha }})K_{i\overline{j}},_i_j_lD_A=0`$ (8)
The equations for the soft terms are now given by $`^\mathrm{?}`$:
$`m_{i\overline{j}}^2`$ $`=`$ $`m_{3/2}^2\left(G_{i\overline{j}}R_{i\overline{j}\alpha \overline{\beta }}G^\alpha G^{\overline{\beta }}\right)+{\displaystyle \underset{A}{}}g_A^2D_A\left(X_i^A+\overline{v}_{\overline{l}}X_{\overline{l}}^A_{\overline{l}}+\eta _A^{\overline{\alpha }}_{\overline{\alpha }}{\displaystyle \frac{1}{2}}D_A\right)G_{i\overline{j}},`$ (9)
$`m_{ij}^2`$ $`=`$ $`m_{3/2}^2\left(2_iG_j+G^\alpha _i_jG_\alpha \right){\displaystyle \frac{1}{2}}{\displaystyle \underset{A}{}}g_A^2D_A^2(_iG_j+{\displaystyle \frac{1}{2}}g_A^2_i_jf_A),`$ (10)
$`A_{ijk}`$ $`=`$ $`m_{3/2}^2\left(3_i_jG_k+G^\alpha _i_j_kG_\alpha \right){\displaystyle \frac{1}{2}}_i_jG_k{\displaystyle \underset{A}{}}g_A^2D_A^2,`$ (11)
where we have identified the gravitino mass $`<e^G>=m_{3/2}^2`$ and $`f_A`$ is the gauge kinetic function. While these expressions are given for the tree level potential, higher order corrections can play a significant role, depending on the specifics of the model of supersymmetry breaking. In models with small tree-level contributions, the dominant set of corrections are of anomaly mediated type$`^\mathrm{?}`$ which are proportional to the gravitino mass $`m_{3/2}`$. A detailed analysis including various parameterisations will be presented in $`^\mathrm{?}`$. The $`\mu `$ term and the tree level gaugino mass terms are then given by
$$\mu _{ij}=m_{3/2}_iG_j,M_{1/2}^A=\frac{1}{2}(Ref_A)^1m_{3/2}f_{A\alpha }G^\alpha .$$
(12)
It is clear from the above analysis that in the F-limit where D term contributions are absent, the soft terms all typically of the same order of magnitude without large hierarchies within themselves. These expressions have been used to parameterise soft terms from superstring theories as well as supergravity$`^\mathrm{?}`$. Of course, if the gravitino mass is very larg $`m_{3/2}>>TeV`$, possible higher derivative operators can change the pattern displayed above completely. A consistent supergravity analysis in such a case, however, becomes considerably more involved.
### 2.1 Implications of large D-terms on the soft parameters
Let us now systematically see the impact of the D-terms on each of the soft parameters. As has been noted earlier, as long as they are of the $`๐ช(m_{3/2}^2)`$, they do not have strong impact. Let us now consider the limit $`m_{3/2}^2\stackrel{<}{}D_A\stackrel{<}{}m_{3/2}M_P`$.
(i). Sfermion Mass Terms: The most dominant contribution to the sfermion masses from the $`D`$-terms are the ones which are linear in $`D`$ which for $`m_{3/2}`$ TeV push the scalar masses to intermediate energy scale. Note that these terms depend on the charges of the fields under the additional $`U(1)`$ gauge group, thus putting a constraint that these charges to be of definite sign. If all the three generations of the sfermions have the same charges under the $`U(1)`$ groups, this term would also be universal.
(ii). Higgs mass terms and the $`B\mu `$: The Higgs masses follow almost the same requirements as the soft masses. Usually, their charges are linked with the Giudice-Masiero mechanism$`^\mathrm{?}`$. The $`B\mu `$ term is however special. Unlike the Higgs mass terms, it does not receive large contributions from D-terms, whose contributions can be utmost of $`๐ช(m_{3/2}^2)`$. If the splitting between the Higgs masses and the $`B\mu `$ is too large, it could lead to unphysical regions in $`\mathrm{tan}\beta `$. To remedy this, alternative schemes have to be devised, an explicit example being discussed in the next subsection.
(iii). A-terms: Even if the D-terms are large, the A-terms are typically proportional to $`๐ช(m_{3/2})`$. No large enhancement is present. This is expected as A-terms break R-symmetries.
### 2.2 A model for Split supersymmetry
The requirement of Split supersymmetry type soft spectra are as follows :
(i) Scalar soft terms : $`m_{\stackrel{~}{f}}^2๐ช(10^910^{15})`$ GeV, $`(\stackrel{~}{f}=Q,u^c,d^c,L,e^c)`$
(ii). Higgs mass parameters $`m_{H_1}^2m_{H_2}^2B_\mu ๐ช(10^910^{15})`$ GeV, with one of them fine tuned to be around the electroweak scale.
(iii). The gaugino masses and the $`\mu `$ term are around the weak scale.
From the discussion in the previous section, it was obvious that it is just not sufficient to choose the $`U(1)`$ charges of the sfermions to be positive to realise the split spectrum<sup>a</sup><sup>a</sup>asee also Ref.$`^\mathrm{?}`$. Since $`B_\mu `$ term does not have large D-term contributions, we need to disentangle the $`\mu `$ and the $`B_\mu `$ term by introducing (at least ) one new field $`X`$ and allowing a term of the type $`XH_1H_2`$ in the superpotential. If the auxiliary field $`F_XD`$, whereas $`X<<m_{3/2}`$, then the tuning of Higgs parameters is technically possible. The minimal field content realising this is as follows. The model contains an additional $`U(1)`$ group, with two additional fields $`X`$ and $`\varphi `$ with charges $`+2`$ and $`1`$. The $`\varphi `$ field can act as a flavon field attaining a large vev close to the fundamental scale. The superpotential and the relevant terms in the Kahler potential, obtained by expanding in powers of the matter fields the full supergravity , are
$`W=W_0+W_{SSM}+\lambda _1XH_1H_2+\lambda _2X\varphi ^2+\mathrm{},`$
$`K=K_0+{\displaystyle \underset{i\overline{j}}{}}Z_{i\overline{j}}\varphi ^i\overline{\varphi }^{\overline{j}}+Z^{}(\varphi ^{})^2H_1H_2+\mathrm{}.`$ (13)
In (13), $`W_0`$ is a holomorphic function of moduli fields , $`K_0`$ is the Kahler potential for moduli, $`Z_{i\overline{j}}`$ and $`Z^{}`$ are generically also moduli dependent and the dots denote higher order terms in an expansion in matter fields. The main features of the model are already captured by performing a vacuum analysis at the global supersymmetry level. In this case, the scalar potential is given by
$$V=\lambda _2^2(|\varphi |^4+4|X|^2|\varphi |^2)+\frac{1}{2}g^2(2|X|^2|\varphi |^2+\xi )^2+\mathrm{},$$
(14)
For $`\xi >0`$, the stable extremum of the above and the auxiliary fields are given by:
$$\varphi =\frac{g^2}{2\lambda _2^2+g^2}\xi ,X=0,F_\varphi =0,F_X=\frac{\lambda _2g^2}{2\lambda _2^2+g^2}\xi ,D=\frac{2\lambda _2^2}{2\lambda _2^2+g^2}\xi .$$
(15)
From the above it is clear that $`F_Xg^2D`$ and moreover are of order of the FI term $`\xi `$. This is sufficient to enable the $`B`$ term to receive large contributions through the term $`G^X_{H_1}_{H_2}G_X`$ in the eq.(10). As long as $`\xi `$ is close to an intermediate scale $`m_{3/2}^2<<\xi m_{3/2}M_P`$, this model seems to replicate the hierarchical split spectrum, if one fixes the gravitino mass at 1 TeV. However, in typical string models, the FI term is of the $`๐ช(M_{Pl}^2/16\pi ^2)`$, which would give a too large contribution to the vacuum energy. The correct order of magnitude could be achieved by incorporating the above model into a higher dimensional theory. For illustration lets us consider a 5D theory compactified over $`S^1/Z_2`$. The Standard Model and the $`X,\varphi `$ fields live on a 3D brane, whereas the gauge fields of the $`U(1)`$ are allowed to propagate in the bulk. We will use Scherk-Schwarz mechanism to break supersymmetry.
The various scales in the problem are $`R=tM_5^1,RM_5^3=M_P^2`$, where $`tReT`$, the modulus field. After canonically normalizing the various fields by $`\widehat{\varphi }_i=\sqrt{t/3}\varphi _i`$ and at the global supersymmetry level, the potential retains the form (14) with $`\xi M_5^2=M_P^2/t`$. The four dimensional gauge coupling is given by $`g^2=1/t=1/(RM_5)`$, whereas the gravitino mass is given by $`m_{3/2}=\omega /R`$, where $`\omega `$ is a number of order one. The D-term contribution to the vacuum energy is then of the form $`V_Dg^2M_5^4m_{3/2}^2M_P^2`$, in the right order as required by the cancellation of the vacuum energy and realisation of the split spectrum. If the no-scale structure is broken by the dynamics, the gauginos attain their masses through anomaly mediation and thus we choose the gravitino mass to be of the order of 100 TeV. In the opposite case, new sources of gaugino masses have to be invoked, like Dirac-type masses or higher dimensional operators if the gravitino is much heavier. The $`\mu `$ term can be generated by Giudice-Masiero mechanism and is $`\mu (v/M_5)^2m_{3/2}`$. So, this model replicates the spectrum of the split supersymmetry at the weak scale using large D-terms of the intermediate scale and a 100 TeV massive gravitino.
### 2.3 Moduli stabilisation problem
As transparent in (2), the FI terms are field (moduli) dependent. If no additional dynamics is present, the moduli fields will always exbihit a runaway behaviour and the FI terms disappear. We resume here the issue of moduli stabilisation with realisation of large D-term contributions to soft terms discussed in $`^\mathrm{?}`$in a context similar to, but having some new features compared the one discussed some time ago in $`^\mathrm{?}`$. We would like to stress that the analysis performed in $`^\mathrm{?}`$ and summarized here is also relevant for the issue of the uplift of the energy density in the context of KKLT type moduli stabilisation $`^\mathrm{?}`$. The gauge group consists of the Standard Model supplemented by a confining hidden sector group and an anomalous $`U(1)_X`$. For simplicity we discuss the case of an supersymetric SU(2) gauge group with one quark flavor $`Q^a`$ and anti-quark $`\stackrel{~}{Q}^a`$ where $`a=1,2`$ is an index in the fundamental representation of the $`SU(2)`$ gauge group. The hidden sector consists of a stack of two magnetised D9 branes in the type I string with kinetic function $`f=S+kT`$, where S is the dilaton (super)field, T a volume (Kahler) modulus and $`k`$ is an positive or negative integer determined by the magnetic fluxes in two compact torii. The dynamical scale of the hidden sector gauge group and the effective superpotential are
$$\mathrm{\Lambda }=M_Pe^{8\pi ^2(S+kT)/5},W=W_0+\frac{\mathrm{\Lambda }^5}{M}+\lambda \phi M.$$
(16)
In order to stabilise the modulus S we invoke the three-form NS-NS and RR fluxes. The low energy dynamics is described by $`M=Q^a\stackrel{~}{Q}^a`$, the composite โmesonโ field. $`W_0`$ depends on the modulus $`S`$ and eventually other (complex structure) moduli of the theory and stabilise them $`S=S_0`$ by giving them a very large mass. If the other relevant mass scales, the FI term and the dynamical scale $`\mathrm{\Lambda }`$ have much lower values, we can safely integrate out these fields, by keeping the T modulus in the low energy dynamics. Minimisation with respect to $`T`$ stabilises also the Kahler modulus. Due to the anomalous nature of the $`U(1)_X`$, there are mixed anomalies with the hidden sector gauge group which translate into a chiral nature of the quark and anti-quark field, such that the sum of their charges, equal to the $`M`$ meson charge, is different from zero and, in our example, equal to $`+1`$. $`\phi `$ is a field of charge $`1`$ which originally participate in the Yukawa coupling $`\lambda \phi Q^a\stackrel{~}{Q}^a`$, which plays the role of meson mass after the spontaneous symmetry breaking of the $`U(1)_X`$. Along the $`SU(2)`$ flat directions, the D-term scalar potential is
$$V_D=\frac{g_X^2}{2}\left[(M^{}M)^{1/2}|\phi |^2+k\mu ^2\right]^2,$$
(17)
where $`\mu `$ is a mass scale determined by the T-modulus vev. The new feature of (17) is that $`k`$ and consequently the FI term can have both signs, whereas in the effective heterotic string framework worked out in $`^\mathrm{?}`$, the FI term had only one possible sign. In the limit $`\mathrm{\Lambda }<<\mu `$, the vacuum structure and the pattern of supersymmetry breaking in the two cases of $`k`$ positive and negative are vastly different :
i) $`k>0`$. In this case the vacuum can be determined as in $`^\mathrm{?}`$. We find, to the lowest order in the parameter $`ฯต(\mathrm{\Lambda }/k^{1/2}\mu )^{5/2}`$, a hierarchically small scale of supersymmetry breaking
$`|\phi |^2=k\mu ^2,M=\lambda ^{1/2}\mathrm{\Lambda }^2(\mathrm{\Lambda }/k^{1/2}\mu )^{1/2},`$
$`D_X={\displaystyle \frac{\lambda \mathrm{\Lambda }^5}{(k^{1/2}\mu )^3}},F_\phi =\mathrm{\Lambda }^2({\displaystyle \frac{\lambda \mathrm{\Lambda }}{k^{1/2}\mu }})^{1/2},F^{\overline{M}}=K^{M\overline{M}}_MW={\displaystyle \frac{\mathrm{\Lambda }^5}{k\mu ^2M_P^2}}.`$ (18)
ii) $`k<0`$. In this case we find, to the lowest order in the parameter $`ฯต^{}(\mathrm{\Lambda }^2/|k|\mu ^2)^5`$ , a large scale of supersymmetry breaking (for the complete expressions, see $`^\mathrm{?}`$)
$$|\phi |\frac{\mathrm{\Lambda }^5}{k^2\mu ^4},M|k|\mu ^2,D_Xk\mu ^2,F_\phi k\mu ^2,F^{\overline{M}}\frac{\mathrm{\Lambda }^5}{k\mu ^2M_P^2}.$$
(19)
Interestingly enough, this second case generate a large scale for supersymmetry breaking with large D-term contributions.
### 2.4 Acknowledgments
We would like to thank the Moriond organizers for providing a stimulating atmosphere. We wish to thank S. Lavignac and Carlos Savoy for useful discussions. This work is supported in part by the CNRS PICS no. 2530 and 3059, INTAS grant 03-51-6346, the RTN grants MRTN-CT-2004-503369, MRTN-CT-2004-005104 and by a European Union Excellence Grant, MEXT-CT-2003-509661. SKV is also supported by Indo-French Centre for Promotion of Advanced Research (CEFIPRA) project No: 2904-2 โBrane World Phenomenologyโ.
## References |
warning/0506/hep-th0506131.html | ar5iv | text | # Phase Transitions Patterns in Relativistic and Nonrelativistic Multi-Scalar-Field Models
## I INTRODUCTION
The study of symmetry breaking (SB) and symmetry restoration (SR) mechanisms have proved to be extremely useful in the analysis of phenomena related to phase transitions in almost all branches of physics. Some topics of current interest which make extensive use of SB/SR mechanisms are topological defects formation in cosmology, the Higgs-Kibble mechanism in the standard model of elementary particles and the Bose-Einstein condensation (BEC) in condensed matter physics. An almost general rule that arises from those studies is that a symmetry which is broken at zero temperature should get restored as the temperature increases. Examples range from the traditional ferromagnet to the more up to date chiral symmetry breaking/restoration in QCD, with the transition pattern being the simplest one of going from the broken phase to the symmetric one as temperature goes from below to above some critical value and vice-versa.
However, a counter-intuitive example may happen in multi-field models, as first noticed by Weinberg weinberg who considered an $`O(N_\varphi )\times O(N_\psi )`$ invariant relativistic model with two types of scalar fields (with $`N_\varphi `$ and $`N_\psi `$ components) and different types of self and crossed interactions. Using the one-loop approximation he has shown that it is possible for the crossed coupling constant to be negative, while the model is still bounded from below, leading, for some parameter values, to an enhanced symmetry breaking effect at high temperatures. This would predict that a symmetry which is broken at $`T=0`$ may not get restored at high temperatures, a phenomenon known as symmetry non restoration (SNR), or, in the opposite case, a symmetry that is unbroken at $`T=0`$ would become broken at high temperatures, thus characterizing inverse symmetry breaking (ISB). Here, one could argue that SNR/ISB are perhaps just artifacts of the simple one-loop perturbative approximation and that the consideration of higher order terms and effects like the temperature dependence of the couplings could change the situation. To answer this question the model has been re-investigated by many other authors using a variety of different methods with most results giving further support to the idea (see, e.g., Ref. borut for a short review of SNR/ISB). For example, the SNR/ISB phenomena were studied using the Wilson Renormalization Group roos and the explicit running of the (temperature dependent) coupling constants has been taken into account, showing that in fact the strength of all couplings increase in approximately the same way with the temperature. This analysis shows that once the couplings are set, at some (temperature) scale, such as to make SNR/ISB possible, the situation cannot be reversed at higher temperatures. Two of the present authors have also treated the problem nonperturbatively taking full account of the cumbersome two-loop contributions MR1 . The results obtained in Ref. MR1 were shown to be in good agreement with those obtained using the renormalization group approach of Ref. roos and, therefore, also support the possibility of SNR/ISB occurring in relativistic multi-scalar field models even at extremely high temperatures, where standard perturbation theory would break down.
The mechanisms of SNR/ISB have found a variety of applications. For instance, in cosmology, where they have been implemented in realistic models, their consequences have been explored in connection with high temperature phase transitions in the early Universe, with applications covering problems involving CP violation and baryogenesis, topological defect formation, inflation, etc moha ; rio . For example, the Kibble-Higgs sector of a $`SU(5)`$ grand unified theory can be mimicked by considering the case $`N_\varphi =90`$ and $`N_\psi =24`$ and has been used to treat the monopole problem moha ; rio ; lozano . Setting $`N_\varphi =N_\psi =1`$ the model becomes invariant under the discrete transformation $`Z_2\times Z_2`$. The latter version has been used in connection with the domain wall problem domainwall . Most applications are listed in Ref. borut which gives an introduction to the subject discussing other contexts in which SNR/ISB can take place in connection with cosmology and condensed matter physics. These interesting results from finite temperature quantum field theory raise important questions regarding their possible manifestation in condensed matter systems which can be described by means of nonrelativistic scalar field theories in the framework of the phenomenology of Ginzburg-Landau potentials, like, for example, in homogeneous Bose gases bec . As far as these systems are concerned, we are unaware of any applications or studies of analogue SNR/ISB phenomena in the context of nonrelativistic scalar field models.
In the context of condensed matter physics more exotic transitions are well known to be possible and similar phenomena to ISB/SNR have been observed in a large variety of materials. One of the best known examples is the symmetry pattern observed in potassium sodium tartrate tetrahydrate, $`KNa(C_4H_4O_6).4H_2O`$, most commonly known as the Rochelle salt, which goes salt , as the temperature increases, from a more symmetric orthorhombic crystalline structure to a less symmetric monoclinic structure at $`T255K`$. It then returns to be orthorhombic phase at $`T297K`$, till it melts at $`T348K`$. It thus exhibits an intermediary inverse symmetry breaking like phenomenon through a reentrant phase. Other materials which arose great interest recently due to their potential applications include, for example, the liquid crystals smectic and spin glass materials spinglass , which exhibit analogous phenomena of having less symmetric phases at intermediary temperature ranges, known as nematic to smectic phases (ferro and antiferro electric and magnetic like phases), and compounds known as the manganites, e.g. $`(Pr,Ca,Sr)MnO_3`$, which can exhibit ferromagnetic like reentrant phases above the Curie (critical) temperature manganites . Actually, in the condensed matter literature we can find many other examples of physical materials exhibiting analogue phenomena of SNR/ISB. This same trend of the emergence of reentrant phases also seems to include low dimensional systems lowd . A discussion on these inverse like symmetry breaking phenomena in condensed matter systems has been recently summarized in Ref. inverse . Another motivation for the present work is the growing interest in investigating parallels between symmetry breaking in particle physics (Cosmology) and condensed matter physics (the Laboratory) as discussed by Rivers rivers in a recent review related to the COSLAB programme. One of the most exciting aspects of such investigation is due to the fact that condensed matter allows for experiments which can, in principle, test models and/or methods used in Cosmology.
Here, our aim is to analyze a nonrelativistic model composed of two different types of multi component fields. To investigate SNR/ISB we consider a model possessing an $`U(1)\times U(1)`$ global symmetry, that is analogous to the $`O(N_\varphi )\times O(N_\psi )`$ relativistic model studied in weinberg ; roos ; MR1 , for $`N_\varphi =N_\psi =2`$, including both one and two-body interactions in the potential. Further, by disregarding the bosonic internal degrees of freedom, the model is considered as representing a system of hard core spheres. In the analysis that follows in the next sections we do not claim that this simplified model described in terms of scalar fields with local interactions will be simulating the phases behavior of any of the condensed matter system cited in the previous paragraph, but just that it suffices, as a toy model, to show the generality of the possibility of emergence of reentrant behavior in some simple condensed matter systems which can be modeled by coupled multi-scalar field models. The chosen non relativistic model is also simple enough to show the differences and analogies regarding the phenomena of SNR/ISB which occurs on its relativistic counterpart.
We will show that, like in the relativistic case, SNR/ISB can take place when thermal effects on the couplings are neglected. We then consider these thermal effects by computing the first one-loop contributions to the couplings finding that, contrary to the relativistic case, SNR/ISB cannot persist indefinitely at higher temperatures when all symmetries are restored. In summary, the possible phase transition patterns seem to be completely different for the relativistic and nonrelativistic cases when the important thermal effects on the couplings are taken into account. This paper is divided as follows. In Sec. II we review the original relativistic prototype model. In Sec. III we present a similar nonrelativistic model of hard core spheres with quadratic and quartic interactions. We show how SNR/ISB cannot occur for such a system when the temperature effects on the couplings are considered, but they can only manifest through reentrant like phases, with symmetry restoration always happening at high enough temperatures. Our conclusions and final remarks are presented in Sec. IV. An appendix is included to show some technical details of the calculations.
## II THE EMERGENCE OF SNR/ISB PHENOMENA IN THE RELATIVISTIC MODEL
At finite temperature the relativistic multi-scalar field theory was first studied by Weinberg weinberg who found evidence of SNR/ISB taking place at finite temperatures. On his work, he considered a prototype model composed of two types of scalar fields, $`\varphi `$ and $`\psi `$ with $`N_\varphi `$ and $`N_\psi `$ components, respectively, which is invariant under the $`O(N_\varphi )\times O(N_\psi )`$ transformation. Such a model has a lagrangian density which can then be written as
$$(\varphi ,\psi )=\frac{1}{2}(_\mu \varphi )^2\frac{m_\varphi ^2}{2}\varphi ^2\frac{\lambda _\varphi }{4!}(\varphi ^2)^2+\frac{1}{2}(_\mu \psi )^2\frac{m_\psi ^2}{2}\psi ^2\frac{\lambda _\psi }{4!}(\psi ^2)^2\frac{\lambda }{4}\varphi ^2\psi ^2.$$
(1)
The self-coupling constants $`\lambda _\varphi `$ and $`\lambda _\psi `$ and the cross coupling $`\lambda `$ in Eq. (1) are traditionally considered as all positive. However, it is still possible to consider $`\lambda `$ negative in (1) provided the potential is kept bounded from below. It is easily seen in this case that the boundness condition for the model (1) requires that the couplings satisfy
$$\lambda _\varphi >0,\lambda _\psi >0,\lambda _\varphi \lambda _\psi >9\lambda ^2.$$
(2)
The fact that the cross coupling, $`\lambda `$, is allowed to be negative has interesting consequences as is seen from the one-loop thermal mass evaluation. As usual, the temperature effects on the zero temperature mass parameters $`m_i^2`$ (where $`i=\varphi `$ or $`\psi `$) can be computed from the (thermal) self-energy corrections $`\mathrm{\Sigma }_i(T)`$ from which the thermal masses, $`M_i^2(T)=m_i^2(0)+\mathrm{\Sigma }_i(T)`$ are obtained. The thermal masses have been first calculated with the one-loop approximation weinberg which, using the usual rules of finite temperature quantum field theory (see e.g. jackiw ; weinberg ) and in the high temperature approximation, $`m_\varphi /T,m_\psi /T1`$, leads to the results
$$M_\varphi ^2(T)m_\varphi ^2+\frac{T^2}{12}\left[\lambda _\varphi \frac{1}{2}\left(\frac{N_\varphi +2}{3}\right)+\lambda \frac{N_\psi }{2}\right],$$
(3)
and
$$M_\psi ^2(T)m_\psi ^2+\frac{T^2}{12}\left[\lambda _\psi \frac{1}{2}\left(\frac{N_\psi +2}{3}\right)+\lambda \frac{N_\varphi }{2}\right],$$
(4)
where we kept only the leading order relevant thermal contributions in the high temperature expansion of $`\mathrm{\Sigma }_i(T)`$, which will be enough for the analysis that follows. Note also that the zero temperature quantum corrections to both masses and coupling constants are divergent quantities and so require renormalization. This is done the standard way by adding the appropriate counterterms of renormalization in (1) (see also MR1 ). We are only interested in the thermal quantities (that are finite) since the zero temperature quantum corrections to masses and couplings can be regarded as negligible as compared to the finite temperature contributions. In Eqs. (3), (4) as well as in the relations below, the mass parameters $`m_\varphi `$ and $`m_\psi `$ and couplings $`\lambda _\varphi ,\lambda _\varphi `$ and $`\lambda `$ are just to be interpreted here as the renormalized quantities instead of the bare ones. It is obvious, from the potential term in the lagrangian density (1), that if one of the mass parameters $`m_i^2`$ is negative the $`O(N_i)`$ symmetry related to that sector is broken at $`T=0`$: $`O(N_i)O(N_i1)`$. Thermal effects tend to restore that symmetry at a certain critical temperature, upon using Eqs. (3) and (4), given by
$$T_{c,i}=\left\{12m_i^2\left[\lambda _i\frac{1}{2}\left(\frac{N_i+2}{3}\right)+\lambda \frac{N_j}{2}\right]^1\right\}^{1/2}.$$
(5)
However, if $`m_i^2<0`$, Eq. (5) shows that for a negative cross-coupling constant, $`\lambda <0`$, and for $`|\lambda |>\lambda _i(N_i+2)/(3N_j)`$, $`T_{c,i}`$ cannot be real. In other words, the broken symmetry is never restored (SNR). At the same time if $`m_i^2>0`$ (unbroken $`O(N_i)`$ symmetry at $`T=0`$), but $`\lambda <0`$ and $`|\lambda |>\lambda _i(N_i+2)/(3N_j)`$ then from Eqs. (3), (4) and (5), we can predict that, as the temperature is increased, the symmetry will be broken at $`T_c`$, instead of being restored (ISB). For example, let us suppose that $`\lambda <0`$ and
$$|\lambda |>\frac{\lambda _\varphi }{N_\psi }\left(\frac{N_\varphi +2}{3}\right).$$
(6)
In this case the boundness condition assures that $`|\lambda |<\lambda _\psi (N_\psi +2)/(3N_\varphi )`$. Then, if $`m_\psi ^2<0`$ one has broken $`O(N_\varphi )`$ symmetry at $`T=0`$, but $`M_\psi ^2(T)`$ will eventually become positive at the corresponding $`T_{c,\psi }`$, given by Eq. (5), restoring the symmetry. If $`m_\psi ^2>0`$, then $`M_\psi ^2(T)>0`$ for all values of $`T`$ and the model is always symmetric under $`O(N_\psi )`$. On the other hand, if $`m_\varphi ^2<0`$, our choice of parameters predicts that the $`O(N_\varphi )`$ symmetry is broken at $`T=0`$ and that it does not get restored at high temperatures, a clear manifestation of SNR. At the same time, if $`m_\varphi ^2>0`$, the $`O(N_\varphi )`$ symmetry, which is unbroken at $`T=0`$, becomes broken at a $`TT_{c,\varphi }`$, which is a manifestation of ISB. Obviously, which field will suffer SNR or ISB depends on our initial arbitrary choice of parameter values. Note that when $`\lambda =0`$ the theory decouples and SNR/ISB cannot take place. In this case one observes the usual SR which happens in the simple $`O(N)`$ scalar model.
An issue that arises, concerning the results discussed above, is that the coupling constants are scale dependent in accordance with the renormalization group equations. Therefore, at high temperatures not only the masses get dressed by thermal corrections but also the coupling constants, so we must answer whether the intriguing phase transitions patterns discussed above, for $`\lambda <0`$, can hold in terms of the equivalent running coupling constants. This issue was analyzed by Roos roos , who used the Wilson Renormalization Group (WRG) to evaluate the $`\lambda _i(T)`$ and $`\lambda (T)`$. His calculations revealed that the strength of all couplings increase, at high $`T`$, in a way which excludes the possibility of SR in cases where SNR/ISB happen. He also showed that the running of coupling constants with temperature as predicted by the one-loop approximation, as adopted in the present work, is robust up to very large scales. In addition to that, the two-loop nonperturbative calculations performed in Ref. MR1 also support, from a qualitative point of view, Weinbergโs one-loop results.
As one notices from the equations which describe the thermal masses, Eqs. (3) and (4), the appearance of SNR/ISB is directly related to the relation among the different couplings when $`\lambda <0`$. It is then useful to define the quantity
$$\mathrm{\Delta }_i=\lambda _i\frac{1}{2}\left(\frac{N_i+2}{3}\right)+\lambda \frac{N_j}{2},$$
(7)
which takes those relations into account. Then, in terms of temperature independent couplings, the critical temperature, Eq. (5), can be written as
$$T_{c,i}=\left(\frac{12M_i^2}{\mathrm{\Delta }_i}\right)^{1/2}.$$
(8)
One can easily see that SNR/ISB may occur when one<sup>1</sup><sup>1</sup>1One of the main results of Ref. MR1 states that SNR/ISB can occur in both sectors, for some parameter values, a situation which is not allowed at the one-loop level. of the $`\mathrm{\Delta }_i`$ is negative.
Let us now check the robustness of SNR/ISB when the effective, temperature dependent couplings are considered. The thermal effects on all the three couplings, at the one-loop order, are considered in terms of the corrections to the four-point 1PI Greenโs functions. All diagrams at the one-loop level contributing to the effective couplings $`\lambda _\varphi (T)`$, $`\lambda _\psi (T)`$ and $`\lambda (T)`$ are shown in Figs. 1, 2 and 3, respectively. These diagrams, with zero external momenta <sup>2</sup><sup>2</sup>2Recall that those are contributions to the effective potential, which generates all 1PI Greenโs function with zero external momenta., are easily computed at finite temperature (see for instance Refs. fendley ; roos ). Using again the high temperature approximation at leading order, one obtains
$$\lambda _\varphi (T)\lambda _\varphi +\frac{3}{8\pi ^2}\mathrm{ln}\left(\frac{T}{M_0}\right)\left[\frac{1}{2}\left(\frac{N_\varphi +8}{9}\right)\lambda _\varphi ^2+\frac{N_\psi }{2}\lambda ^2\right],$$
(9)
$$\lambda _\psi (T)\lambda _\psi +\frac{3}{8\pi ^2}\mathrm{ln}\left(\frac{T}{M_0}\right)\left[\frac{1}{2}\left(\frac{N_\psi +8}{9}\right)\lambda _\psi ^2+\frac{N_\varphi }{2}\lambda ^2\right],$$
(10)
and
$$\lambda (T)\lambda +\frac{\lambda }{8\pi ^2}\mathrm{ln}\left(\frac{T}{M_0}\right)\left[\frac{1}{2}\left(\frac{N_\varphi +2}{3}\right)\lambda _\varphi +\frac{1}{2}\left(\frac{N_\psi +2}{3}\right)\lambda _\psi \right]+\frac{\lambda ^2}{4\pi ^2}\mathrm{ln}\left(\frac{T}{M_0}\right),$$
(11)
where $`M_0`$ is a regularization scale. In writing the above equations we are once again assuming that the tree-level couplings in Eqs. (9), (10) and (11) are the renormalized ones and we also are only showing the relevant high temperature corrections. The same expressions were also obtained by Roos in roos (note however that different normalizations for the tree-level potential as well as $`M_0=T_0,N_\varphi =N_\psi =1`$ were used in that reference). In roos , the numerical solution of the one-loop Wilson renormalization group equations was also compared to the usual flow equations for the constants obtained from the one-loop beta-functions and shown to agree well with each other up to very high scales. The flow equations referring to the perturbative effective coupling constants Eqs. (9), (10) and (11) are expressed in term of the dimensionless scale $`T/M_0`$ fendley as
$$\frac{d\lambda _\varphi (t)}{dt}=\frac{3}{8\pi ^2}\left[\frac{1}{2}\left(\frac{N_\varphi +8}{9}\right)\lambda _\varphi ^2(t)+\frac{N_\psi }{2}\lambda ^2(t)\right],$$
(12)
$$\frac{d\lambda _\psi (t)}{dt}=\frac{3}{8\pi ^2}\left[\frac{1}{2}\left(\frac{N_\psi +8}{9}\right)\lambda _\psi ^2(t)+\frac{N_\varphi }{2}\lambda ^2(t)\right],$$
(13)
and
$$\frac{d\lambda (t)}{dt}=\frac{\lambda (t)}{8\pi ^2}\left[\frac{1}{2}\left(\frac{N_\varphi +2}{3}\right)\lambda _\varphi (t)+\frac{1}{2}\left(\frac{N_\psi +2}{3}\right)\lambda _\psi (t)\right]+\frac{\lambda ^2(t)}{4\pi ^2},$$
(14)
where $`t=\mathrm{ln}(T/M_0)`$ was used. The solutions of the flow equations (12), (13) and (14), with initial conditions given by the renormalized tree-level coupling constants, can also be easily seen to be equivalent to the solutions for the set of linear coupled equations,
$`\lambda _\varphi (T)`$ $`=`$ $`\lambda _\varphi +{\displaystyle \frac{3}{8\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{N_\varphi +8}{9}}\right)\lambda _\varphi \lambda _\varphi (T)+{\displaystyle \frac{N_\psi }{2}}\lambda \lambda (T)\right],`$
$`\lambda _\psi (T)`$ $`=`$ $`\lambda _\psi +{\displaystyle \frac{3}{8\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{N_\psi +8}{9}}\right)\lambda _\psi \lambda _\psi (T)+{\displaystyle \frac{N_\varphi }{2}}\lambda \lambda (T)\right],`$
$`\lambda (T)`$ $`=`$ $`\lambda +{\displaystyle \frac{\lambda }{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{N_\varphi +2}{3}}\right)\lambda _\varphi (T)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{N_\psi +2}{3}}\right)\lambda _\psi (T)\right]`$ (15)
$`+`$ $`{\displaystyle \frac{\lambda (T)}{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)\left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{N_\varphi +2}{3}}\right)\lambda _\varphi +{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{N_\psi +2}{3}}\right)\lambda _\psi \right]+{\displaystyle \frac{1}{4\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)\lambda \lambda (T).`$
The results obtained from the flow equations given above, or equivalently from the solutions of coupled set of equations (15), are standard ways of nonperturbatively resumming the leading order corrections (in this case the leading log temperature dependent corrections) to the coupling constants. For instance, Eq. (15) is exactly the analogous procedure used for the one-field case for summing all ladder (1-loop or bubble) contributions to the effective coupling constant. For the multi-field case, the perturbative approximation for (15) is again given by Eqs. (9), (10) and (11) at the one-loop level. In our case, these equations are useful to test how robust is the phenomena of SNR/ISB and will be used below in our analysis. Later, in the next section for the nonrelativistic limit of Eq. (1), we will also construct the analogous of these nonperturbative equations for the temperature dependent effective couplings.
In terms of the effective temperature dependent couplings, $`\lambda _\varphi (T)`$, $`\lambda _\psi (T)`$ and $`\lambda (T)`$ the quantity analogous to Eq. (7) becomes
$$\mathrm{\Delta }_i(T)=\lambda _i(T)\frac{1}{2}\left(\frac{N_i+2}{3}\right)+\lambda (T)\frac{N_j}{2},$$
(16)
or, more explicitly, using Eqs. (9), (10) and (11),
$`\mathrm{\Delta }_i(T)`$ $`=`$ $`\lambda _i\left({\displaystyle \frac{N_i+2}{6}}\right)+\lambda {\displaystyle \frac{N_j}{2}}+\lambda _i^2{\displaystyle \frac{(N_i+8)(N_i+2)}{288\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)+\lambda N_j{\displaystyle \frac{(N_i+2)\lambda _i+(N_j+2)\lambda _j}{96\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right)`$ (17)
$`+`$ $`\lambda ^2{\displaystyle \frac{N_j(N_i+6)}{32\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{T}{M_0}}\right).`$
It is clear from the expressions for the effective couplings, Eqs. (9), (10), (11) and Eq. (17), that for perturbative values for the tree-level coupling parameters the predicted results for SNR/ISB are very stable even for very large temperatures (in units of the regularization scale $`M_0`$) which is due to the slow logarithmic change with the temperature. As an illustration, consider for example the tree-level coupling parameters that satisfy the boundness condition Eq. (2), $`\lambda _\varphi =7\times 10^5`$, $`\lambda _\psi =5\times 10^4`$ and $`\lambda =6\times 10^5`$ and $`N_\varphi =N_\psi =2`$. For these values of parameters Eq. (6) is satisfied and the one-loop equations for the effective masses predict ISB or SNR, along the $`\varphi `$ direction, for $`m_\varphi ^2>0`$ or $`m_\varphi ^2<0`$, respectively. Fig. 4 shows that the boundness condition also holds true for the effective temperature dependent couplings. Fig. 5 shows the quantity $`\mathrm{\Delta }_\varphi (T)`$, defined by Eq. (16), which remains negative for the whole range of temperatures considered, thus predicting SNR/ISB along the $`\varphi `$ direction, in accordance with the WRG results roos . At the same time, $`\mathrm{\Delta }_\psi (T)`$, shown in Fig. 6, remains always positive.
Note that the apparent almost constancy in a wide range of temperatures seem from the Figs. 4, 5 and 6 is only a consequence of the effective couplings be only logarithmically dependent on $`T`$ and the very small values for the tree-level couplings that we have considered. Had we taken larger values for the tree-level couplings, obviously would lead to a much larger variation with increasing temperature.
Given the results shown above for the relativistic case, we can conclude, therefore, that the inclusion of thermal effects on the couplings does not exclude the possibility of SNR/ISB occurring at high temperatures. We recall that although the results were obtained with the one-loop approximation this feature does not seem to be an artifact of perturbation theory as confirmed by the results produced by nonperturbative methods, such as the Wilson Renormalization Group procedure used in Ref. roos , as well as the optimized perturbation theory used in Ref. MR1 , where not only thermal corrections to the couplings are accounted for but also to the masses (like in the Schwinger-Dyson or gap equations for the masses).
## III SEARCHING FOR SNR/ISB PATTERNS IN THE nonrelativistic CASE
We now turn our attention to the analysis of similar SNR/ISB phenomena displayed by the relativistic model, given by the lagrangian density, Eq. (1), in the case of its nonrelativistic counterpart. Let us first recall some fundamental differences between relativistic and nonrelativistic theories that will be important in our analysis. Firstly, the obvious reduction from Lorentz to Galilean invariance. Secondly, it should be noted that in the nonrelativistic description particle number is conserved and so, only complex fields are allowed. This second point will be particularly important to us since, for the processes entering in the effective couplings shown in Figs. 1, 2 and 3, only those that do not change particle number (the elastic processes) will be allowed (e.g. this selects the processes (a) shown in Fig. 3 but not the (b) and (c), inelastic, ones). Another important difference between relativistic and nonrelativistic models concerns the structure of the respective propagators. While the relativistic propagator allows for both forward and backward particle propagation (which is associated to particles and anti-particles, respectively), the nonrelativistic propagator of scalar theories at $`T=0`$ only has forward propagation (see e.g. the discussion in Ref. oren ). Note however that the structure of the propagators (or two-point Greenโs function) in a thermal bath includes both backward and forward propagation mahan , which can be interpreted in terms of excitations to and from the thermal bath (or, equivalently, emission and absorption of particle to and from the thermal bath KB ).
We should also say that, alternatively to the derivation of the nonrelativistic analog of (1), we could as well consider the relevant equations leading e.g. to the derived effective couplings in the previous section and take the appropriate low-energy limit for those equations. However it is more practical, and indeed it is the procedure usually adopted in atomic and low energy nuclear physics, to start directly from the nonrelativistic Hamiltonian or Lagrangian densities. This is an one step procedure leading, say, to the Feynman rules that can be applied to any other quantity that we may be interested in computing, without having first to compute the corresponding relativistic expressions and then working out the corresponding nonrelativistic. So, let us now initially consider the nonrelativistic limit of the lagrangian density given by Eq. (1). This can be obtained by first expressing the fields $`\varphi `$ and $`\psi `$ in terms of (complex) nonrelativistic fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ as oren ; davis ; zee
$$\varphi (\stackrel{}{x},t)=\frac{1}{\sqrt{2m_\varphi }}\left[\mathrm{exp}(im_\varphi t)\mathrm{\Phi }(\stackrel{}{x},t)+\mathrm{exp}(im_\varphi t)\mathrm{\Phi }^{}(\stackrel{}{x},t)\right],$$
(18)
and
$$\psi (\stackrel{}{x},t)=\frac{1}{\sqrt{2m_\psi }}\left[\mathrm{exp}(im_\psi t)\mathrm{\Psi }(\stackrel{}{x},t)+\mathrm{exp}(im_\psi t)\mathrm{\Psi }^{}(\stackrel{}{x},t)\right],$$
(19)
where it is assumed that the fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ oscillate in time much more slowly than $`\mathrm{exp}(im_\varphi t)`$ and $`\mathrm{exp}(im_\psi t)`$, respectively. By substituting (18) and (19) in (1) and taking the nonrelativistic limit of large masses, the oscillatory terms with frequencies $`m_\varphi `$ and $`m_\psi `$ can be dropped. The resulting lagrangian density in terms of $`\mathrm{\Phi }`$, $`\mathrm{\Psi }`$ and complex conjugate fields becomes
$`(\mathrm{\Phi }^{},\mathrm{\Phi },\mathrm{\Psi }^{},\mathrm{\Psi })`$ $`=`$ $`{\displaystyle \frac{1}{2m_\varphi }}\left[im_\varphi (_t\mathrm{\Phi }^{})\mathrm{\Phi }+im_\varphi \mathrm{\Phi }^{}(_t\mathrm{\Phi })|\mathrm{\Phi }|^2+\left|_t\mathrm{\Phi }\right|^2\right]{\displaystyle \frac{\lambda _\varphi }{16m_\varphi ^2}}(\mathrm{\Phi }^{}\mathrm{\Phi })^2`$ (20)
$`+`$ $`{\displaystyle \frac{1}{2m_\varphi }}\left[im_\psi (_t\mathrm{\Psi }^{})\mathrm{\Psi }+im_\psi \mathrm{\Psi }^{}(_t\mathrm{\Psi })|\mathrm{\Psi }|^2+\left|_t\mathrm{\Psi }\right|^2\right]{\displaystyle \frac{\lambda _\psi }{16m_\psi ^2}}(\mathrm{\Psi }^{}\mathrm{\Psi })^2`$
$``$ $`{\displaystyle \frac{\lambda }{4m_\varphi m_\psi }}(\mathrm{\Phi }^{}\mathrm{\Phi })(\mathrm{\Psi }^{}\mathrm{\Psi }),`$
where we have assumed for simplicity, in the derivation of the last term in (20), the cross-fields interaction term, that $`m_\varphi m_\psi `$ <sup>3</sup><sup>3</sup>3Note that for equal masses there is the possibility of an additional symmetric interaction term of the form $`[(\mathrm{\Phi }\mathrm{\Psi }^{})^2+(\mathrm{\Phi }^{}\mathrm{\Psi })^2]`$ in (20), however this term will not be relevant for our analysis and conclusions since it can be absorbed in a redefinition of the cross-coupling constant $`\lambda `$, especially when we work with densities, or averages of the fields, like in an effective potential calculation.. By further considering
$`\left|_t\mathrm{\Phi }\right|^22m_\varphi \mathrm{Im}(\mathrm{\Phi }_t\mathrm{\Phi }^{}),`$
$`\left|_t\mathrm{\Psi }\right|^22m_\psi \mathrm{Im}(\mathrm{\Psi }_t\mathrm{\Phi }^{}),`$ (21)
we can omit the terms with two time derivatives in Eq. (20). So the Lorentz invariance in (20) is lost and the nonrelativistic analogue of (1) is obtained. The interaction terms in Eq. (20) are the same as those obtained by approximating the usual nonrelativistic two-body interaction potentials by hard core (delta) potentials, e.g.,
$`{\displaystyle d^3x\mathrm{\Phi }^{}(๐ฑ,t)\mathrm{\Phi }(๐ฑ,t)V_\mathrm{\Phi }(๐ฑ๐ฑ^{})\mathrm{\Phi }^{}(๐ฑ^{},t)\mathrm{\Phi }(๐ฑ^{},t)}g_\mathrm{\Phi }\left[\mathrm{\Phi }^{}(๐ฑ,t)\mathrm{\Phi }(๐ฑ,t)\right]^2,`$
$`{\displaystyle d^3x\mathrm{\Psi }^{}(๐ฑ,t)\mathrm{\Psi }(๐ฑ,t)V_\mathrm{\Psi }(๐ฑ๐ฑ^{})\mathrm{\Psi }^{}(๐ฑ^{},t)\mathrm{\Psi }(๐ฑ^{},t)}g_\mathrm{\Psi }\left[\mathrm{\Psi }^{}(๐ฑ,t)\mathrm{\Psi }(๐ฑ,t)\right]^2,`$
$`{\displaystyle d^3x\mathrm{\Phi }^{}(๐ฑ,t)\mathrm{\Phi }(๐ฑ,t)V_{\mathrm{\Phi }\mathrm{\Psi }}(๐ฑ๐ฑ^{})\mathrm{\Psi }^{}(๐ฑ^{},t)\mathrm{\Psi }(๐ฑ^{},t)}g_{\mathrm{\Phi }\mathrm{\Psi }}\left[\mathrm{\Phi }^{}(๐ฑ,t)\mathrm{\Phi }(๐ฑ,t)\right]\left[\mathrm{\Psi }^{}(๐ฑ,t)\mathrm{\Psi }(๐ฑ,t)\right].`$ (22)
The approximation of the two-body potential interactions like in (22) is also commonly adopted in the description of cold dilute atomic systems, where only binary type interactions at low energy are relevant. In that case, the local coupling parameters $`g_\mathrm{\Phi },g_\mathrm{\Psi }`$ and $`g_{\mathrm{\Phi }\mathrm{\Psi }}`$ are also associated to the s-wave scattering lengths $`a_i`$ bec , e.g., $`g_i=2\pi a_i/m_i`$. For nonrelativistic systems in general, besides the two-body interaction terms like (22) (in the hard core approximation) we can also include additional one-body like interaction terms, e.g., $`\kappa _\mathrm{\Phi }\mathrm{\Phi }^{}\mathrm{\Phi }`$, etc. This is the case when we submit the system to an external potential (for example a magnetic field). It can also represent an internal energy term (like the internal molecular energy relative to free atoms in which case the fields in the lagrangian would be related to molecular dimers). In models of superconductivity a constant one-body like interaction term represents the opening of an explicit gap of energy in the system. In the grand-canonical formulation $`\kappa _i`$ can represent chemical potentials included in the action formulation, so that one can also describe density effects (in addition to those from the temperature). In order to retain the symmetry breaking analogies to the previous relativistic model (1), and since our intention here is to keep the analysis as general as possible, we shall also consider additional one-body interaction terms for $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ while their precise interpretation is left as open and will depend on the particular system Eq. (20) is intented to represent.
With the considerations assumed above, we therefore take the following nonrelativistic lagrangian model, that is analogue to the relativistic model Eq. (1),
$`(\mathrm{\Phi }^{},\mathrm{\Phi },\mathrm{\Psi }^{},\mathrm{\Psi })`$ $`=`$ $`\mathrm{\Phi }^{}\left(i_t+{\displaystyle \frac{1}{2m_\mathrm{\Phi }}}^2\right)\mathrm{\Phi }\kappa _\mathrm{\Phi }\mathrm{\Phi }^{}\mathrm{\Phi }{\displaystyle \frac{g_\mathrm{\Phi }}{3!}}(\mathrm{\Phi }^{}\mathrm{\Phi })^2`$ (23)
$`+`$ $`\mathrm{\Psi }^{}\left(i_t+{\displaystyle \frac{1}{2m_\mathrm{\Psi }}}^2\right)\mathrm{\Psi }\kappa _\mathrm{\Psi }\mathrm{\Psi }^{}\mathrm{\Psi }{\displaystyle \frac{g_\mathrm{\Psi }}{3!}}(\mathrm{\Psi }^{}\mathrm{\Psi })^2`$
$``$ $`g(\mathrm{\Phi }^{}\mathrm{\Phi })(\mathrm{\Psi }^{}\mathrm{\Psi }),`$
where we have expressed the derivative terms in their more common form (by doing an integration by parts in the action context). The numerical factors and signs in the one and two-body potential terms in (23) have been chosen in such a way so that the potential in (23) are analogous the one considered in (1). The coupling constants shown in (23) are related to those in (20) by $`g_i=3\lambda _i/(8m_i^2)`$ (with $`i=\mathrm{\Phi },\mathrm{\Psi }`$) and $`g=\lambda /(4m_\mathrm{\Phi }m_\mathrm{\Psi })`$ while $`m_\mathrm{\Phi }`$ and $`m_\mathrm{\Psi }`$ represent the (atomic) masses. In addition, notice that for the nonrelativistic limit which leads to Eq. (23) to be valid, one must keep $`Tm_i`$. Since for nonrelativistic systems in general, the masses $`m_i`$ are of order of typical atomic masses, $`m_i๐ช(1100)\mathrm{GeV}`$, and the typical temperatures in condensed matter systems are at most of order of a few eV, this condition will always hold for the ranges of temperature we will be interested in below.
For multi-component fields, Eq. (23) is the nonrelativistic multi-scalar model with symmetry $`U(N_\mathrm{\Phi })\times U(N_\mathrm{\Psi })`$ that is the analogue of the original relativistic model Eq. (1). For simplicity, in the following we assume the simplest version of (23) where $`N_\mathrm{\Phi }=N_\mathrm{\Psi }=1`$, corresponding to an $`U(1)\times U(1)`$ symmetric model. In this case, by writing the complex fields in terms of real components,
$`\mathrm{\Phi }={\displaystyle \frac{1}{\sqrt{2}}}\left(\varphi _1+i\varphi _2\right),\mathrm{\Psi }={\displaystyle \frac{1}{\sqrt{2}}}\left(\psi _1+i\psi _2\right),`$ (24)
we see that the lagrangian model (23) falls in the same class of universality as that of Eq. (1) for the case of the $`O(2)\times O(2)`$ symmetry. The extension to higher symmetries can be done starting from (23) but the case of simplest symmetry involving the coupling of complex scalar fields will already be sufficient for our study (physically, this system may, for example, describe the coupling of Bose atoms or molecules in an atomic dilute gas system).
Just like in the relativistic case, in the current application we consider the $`\kappa _i`$, that appears in Eq. (23), as simple temperature independent parameters for which thermal corrections arise from the evaluation of the corresponding field self-energies. Now, to make contact with the analogous potential used in the prototype relativistic models for SNR/ISB, we take the overall potential as being repulsive, bounded from below. This requirement imposes a constraint condition analogous to the one found in the relativistic case, $`g_\mathrm{\Psi }>0`$, $`g_\mathrm{\Phi }>0`$ and $`g_\mathrm{\Psi }g_\mathrm{\Phi }>9g^2`$. At a given temperature, the phase structure of the model is then given by the sign of $`\kappa _i(T)=\kappa _i+\mathrm{\Sigma }_i`$, where $`\mathrm{\Sigma }_i`$ is the field temperature dependent self-energy. In the broken phase $`\kappa _i(T)<0`$, while in the symmetric phase $`\kappa _i(T)>0`$. The phase transition occurs at $`\kappa _i(T=T_c^i)=0`$. At the one-loop level the diagrams contributing to the field self-energies are the same as those in the relativistic case. The main difference is that the momentum integrals in the loops are now given in terms of the nonrelativistic propagators for $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$,
$$D_i(\omega _n,๐ช)=\frac{1}{i\omega _n+\omega _i(๐ช)},$$
(25)
where $`\omega _n=2\pi nT`$ ($`n=0,\pm 1,\pm 2,\mathrm{}`$) are the bosonic Matsubara frequencies and $`\omega _i(๐ช)=๐ช^2/(2m_i)+\kappa _i`$. At the one-loop level we then obtain (see appendix)
$$\mathrm{\Sigma }_i(T)=T\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}\frac{d^3q}{(2\pi )^3}\left[\frac{2g_i/3}{i\omega _n+\omega _i(๐ช)}+\frac{g}{i\omega _n+\omega _j(๐ช)}\right],$$
(26)
The sum in Eq. (26) can be easily performed and the resulting momentum integrals lead to well known Bose integrals (see e.g. Ref. pathria and also the appendix for the derivations of these equations). We then obtain the results
$$\kappa _\mathrm{\Phi }(T)=\kappa _\mathrm{\Phi }+\frac{2g_\mathrm{\Phi }}{3}\left(\frac{m_\mathrm{\Phi }T}{2\pi }\right)^{3/2}\mathrm{Li}_{3/2}[\mathrm{exp}(\kappa _\mathrm{\Phi }/T)]+g\left(\frac{m_\mathrm{\Psi }T}{2\pi }\right)^{3/2}\mathrm{Li}_{3/2}[\mathrm{exp}(\kappa _\mathrm{\Psi }/T)],$$
(27)
and
$$\kappa _\mathrm{\Psi }(T)=\kappa _\mathrm{\Psi }+\frac{2g_\mathrm{\Psi }}{3}\left(\frac{m_\mathrm{\Psi }T}{2\pi }\right)^{3/2}\mathrm{Li}_{3/2}[\mathrm{exp}(\kappa _\mathrm{\Psi }/T)]+g\left(\frac{m_\mathrm{\Phi }T}{2\pi }\right)^{3/2}\mathrm{Li}_{3/2}[\mathrm{exp}(\kappa _\mathrm{\Phi }/T)],$$
(28)
where $`\mathrm{Li}_n(z_i)`$ is the polylogarithmic function. Like in the relativistic effective mass terms, in Eqs. (27) and (28) we have once again limited to showing the temperature dependent corrections coming from $`\mathrm{\Sigma }_i`$, omitting the divergent (zero-point energy terms) contributions to the effective one-body terms (so the tree-level parameters in Eqs. (27) and (28) are assumed to be already the renormalized ones). Let us now consider the high temperature approximation, which in the nonrelativistic case is valid as long as $`\kappa _iTm_i`$. Then, one is allowed to consider the approximation for the polylogarithmic functions in Eqs. (27) and (28), $`\mathrm{Li}_{3/2}[\mathrm{exp}(\kappa _i/T)]\zeta (3/2)`$, where $`\zeta (x)`$ is the Riemann zeta function and $`\zeta (3/2)2.6124`$. One can then write the two equations (27) and (28) in the high temperature approximation more compactly, as
$$\kappa _i(T)\kappa _i+\left(\frac{T}{2\pi }\right)^{3/2}\zeta (3/2)\mathrm{\Delta }_i^{\mathrm{NR}},$$
(29)
where we defined the quantity $`\mathrm{\Delta }_i^{\mathrm{NR}}`$ analogous to that of the relativistic case,
$$\mathrm{\Delta }_i^{\mathrm{NR}}=\frac{2}{3}g_im_i^{3/2}+gm_j^{3/2},$$
(30)
in terms of which we obtain the critical temperature for symmetry restoration/breaking analogous to the relativistic expression,
$$T_{c,i}^{\mathrm{NR}}=2\pi \left[\frac{\kappa _i}{\mathrm{\Delta }_i^{\mathrm{NR}}\zeta (3/2)}\right]^{2/3}.$$
(31)
Eq. (31) shows that there are three interesting cases which depend on the sign and magnitude of the cross coupling $`g`$. Taking $`\kappa _i<0`$ and $`g>0`$ one observes a shift in the critical temperatures indicating that the transition occurs at lower temperatures compared to the decoupled case ($`g=0`$). If $`g<0`$ but $`|g|<(2/3)g_i(m_i/m_j)^{3/2}`$ then the transition occurs at higher temperatures. Despite these quantitative differences symmetry restoration does take place in both cases. Now consider for instance the case, with $`\kappa _\mathrm{\Phi }<0`$, where $`g<0`$ but $`|g|>(2/3)g_\mathrm{\Phi }(m_\mathrm{\Phi }/m_\mathrm{\Psi })^{3/2}`$ (in this case, and assuming $`m_\mathrm{\Phi }m_\mathrm{\Psi }`$, the boundness condition assures that $`|g|<(2/3)g_\mathrm{\Psi }`$). Under these conditions, for the $`\mathrm{\Phi }`$ field we have a similar situation as that for the corresponding relativistic case studied in Sec. II, where, Eq. (31) does not give a finite, positive real quantity. This is a manifestation of ISB (for $`\kappa _\mathrm{\Phi }<0`$, or SNR, for $`\kappa _\mathrm{\Phi }>0`$) within our two-field complex nonrelativistic model being analogous to what is seen in the relativistic case. At the same time the field $`\mathrm{\Psi }`$ suffers the expected phase transition at a higher $`T_c`$ compared to the $`g=0`$ case. As in the relativistic case, which field will suffer SNR/ISB depends on our initial choice of parameters in the tree level potential, so it is model dependent. The SNR/ISB result as seen above in this nonrelativistic model is however misleading, as we next show by considering the same phenomenon in terms of the effective, temperature dependent, coupling parameters.
Physically, the possibility of SNR/ISB occurring at high temperatures as predicted by the naive perturbative approximation, Eq. (29), is much harder to be accepted in the nonrelativistic case (condensed matter) than in the relativistic case (cosmology), which already seems to indicate that the results here should be different when the simple perturbative calculations are improved.
Following the analogy with the relativistic model calculation done in Sec. II, we next evaluate the temperature effects on the nonrelativistic couplings.
The diagrams contributing to the effective self-couplings $`g_\mathrm{\Phi }(T)`$ and $`g_\mathrm{\Psi }(T)`$, at the one-loop level, are those shown in Figs. 1 and 2, respectively, except that the s-channel ones with internal propagators for fields different from those in the external legs (the second one-loop diagrams in Figs. 1 and 2 made of vertices nonconserving particles) are absent. For the effective cross-coupling $`g(T)`$, as discussed at the beginning of this section, the diagrams contributing at the one-loop level are the ones shown in Fig. 3(a), corresponding to the particle number preserving processes. Using (25) for the nonrelativistic field propagators, the explicit expressions for the effective couplings at one-loop order are found to be (see also the appendix)
$`g_\mathrm{\Phi }(T)=g_\mathrm{\Phi }`$ $``$ $`{\displaystyle \frac{g_\mathrm{\Phi }^2}{3}}{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\frac{1}{2\omega _\mathrm{\Phi }(๐ช)}\left\{1+2n(\omega _\mathrm{\Phi })+8\beta \omega _\mathrm{\Phi }n(\omega _\mathrm{\Phi })\left[1+n(\omega _\mathrm{\Phi })\right]\right\}}`$ (32)
$``$ $`3g^2{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\beta n(\omega _\mathrm{\Psi })\left[1+n(\omega _\mathrm{\Psi })\right]},`$
$`g_\mathrm{\Psi }(T)=g_\mathrm{\Psi }`$ $``$ $`{\displaystyle \frac{g_\mathrm{\Psi }^2}{3}}{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\frac{1}{2\omega _\mathrm{\Psi }(๐ช)}\left\{1+2n(\omega _\mathrm{\Psi })+8\beta \omega _\mathrm{\Psi }n(\omega _\mathrm{\Psi })\left[1+n(\omega _\mathrm{\Psi })\right]\right\}}`$ (33)
$``$ $`3g^2{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\beta n(\omega _\mathrm{\Phi })\left[1+n(\omega _\mathrm{\Phi })\right]},`$
and
$`g(T)=g`$ $``$ $`{\displaystyle \frac{2g}{3}}{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\beta \left\{g_\mathrm{\Phi }n(\omega _\mathrm{\Phi })\left[1+n(\omega _\mathrm{\Phi })\right]+g_\mathrm{\Psi }n(\omega _\mathrm{\Psi })\left[1+n(\omega _\mathrm{\Psi })\right]\right\}}`$ (34)
$``$ $`g^2{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\left\{\frac{\omega _\mathrm{\Psi }}{\omega _\mathrm{\Psi }^2\omega _\mathrm{\Phi }^2}\left[1+2n(\omega _\mathrm{\Phi })\right]\frac{\omega _\mathrm{\Phi }}{\omega _\mathrm{\Psi }^2\omega _\mathrm{\Phi }^2}\left[1+2n(\omega _\mathrm{\Psi })\right]\right\}}.`$
The numerical factors in Eqs. (32), (33) and (34) are due to the symmetries of the diagrams and normalizations chosen in Eq. (23). $`n(\omega )`$ in the above equations is the Bose-Einstein distribution.
The zero temperature contributions in Eqs. (32) - (34) are divergent and require proper renormalization. This is mostly simply done by performing the momentum integrals in $`d=3ฯต`$ dimensions and the resulting integrals are all found to be finite in dimensional regularization (when taking $`ฯต0`$ at the end). The momentum integrals for the finite temperature contributions in Eqs. (32) - (34) can again be performed with the help of the Bose integrals given in the appendix. We can further simplify the equations by considering parameters such as $`m_\mathrm{\Phi }m_\mathrm{\Psi }=m`$, $`\kappa _\mathrm{\Phi }\kappa _\mathrm{\Psi }=\kappa `$ and temperatures satisfying $`\kappa Tm`$ in which case the zero temperature corrections to the couplings are negligible compared to the finite temperature ones and can safely be neglected. At leading order, in $`T`$, we then obtain the results
$$g_\mathrm{\Phi }(T)g_\mathrm{\Phi }\frac{mT}{12\pi }\sqrt{\frac{2m}{\kappa }}\left(5g_\mathrm{\Phi }^2+9g^2\right)+๐ช(\kappa /T),$$
(35)
$$g_\mathrm{\Psi }(T)g_\mathrm{\Psi }\frac{mT}{12\pi }\sqrt{\frac{2m}{\kappa }}\left(5g_\mathrm{\Psi }^2+9g^2\right)+๐ช(\kappa /T),$$
(36)
and
$$g(T)g\frac{mT}{4\pi }\sqrt{\frac{2m}{\kappa }}g\left(g+\frac{2g_\mathrm{\Phi }}{3}+\frac{2g_\mathrm{\Psi }}{3}\right)+๐ช(\kappa /T).$$
(37)
Note from Eqs. (35), (36) and (37) that the effective couplings in the nonrelativistic theory have a much stronger dependence with the temperature than in those in the equivalent relativistic theory, Eqs. (9), (10) and (11). We therefore expect to see larger deviations at high temperatures for the effective couplings as compared with the same case in the relativistic problem (by high temperature we mean here temperatures larger than the typical one-body potential coefficients in (23), but much less than the particle masses, see above). It is also evident from the analysis of higher loop corrections to the effective couplings in the nonrelativistic model that all bubble like corrections contribute with the same power in temperature as the one-loop terms, which can easily be checked by simple power-counting in the momentum. A side effect of this is the breakdown, at high temperatures, of the simple one-loop perturbation theory applied here. Another symptom is the apparent running of the effective self-couplings, shown above, to negative values for sufficiently high temperatures given by $`TT_{\mathrm{break}}\mathrm{min}(12\pi \sqrt{\kappa /(2m^3)}g_\mathrm{\Phi }/(5g_\mathrm{\Phi }^2+9g^2),\mathrm{\hspace{0.33em}12}\pi \sqrt{\kappa /(2m^3)}g_\mathrm{\Psi }/(5g_\mathrm{\Psi }^2+9g^2))`$. Nevertheless, it is easy to check that (for the parameters adopted below) the results obtained by just plugging Eqs. (35), (36) and (37) above into Eqs. (29) already show a drastic qualitative difference between this simple improved approximation and the naive perturbative evaluation given by Eq. (29). It looks that for the nonrelativistic case SNR/ISB is a mere artifact of perturbation theory. Intuitively, this is already a rather satisfactory result which, as we will show below, will be confirmed by a nonperturbative resummation of the bubble like corrections. The resumming of all leading order bubble corrections to the couplings can again be done by solving the set of homogeneous linear equations for $`g_\mathrm{\Phi }(T),g_\mathrm{\Psi }(T)`$ and $`g(T)`$,
$`g_\mathrm{\Phi }(T)`$ $`=`$ $`g_\mathrm{\Phi }g_\mathrm{\Phi }(T){\displaystyle \frac{g_\mathrm{\Phi }}{3}}I_1(\beta \kappa _\mathrm{\Phi })3g(T)gI_2(\beta \kappa _\mathrm{\Psi }),`$
$`g_\mathrm{\Psi }(T)`$ $`=`$ $`g_\mathrm{\Psi }g_\mathrm{\Psi }(T){\displaystyle \frac{g_\mathrm{\Psi }}{3}}I_1(\beta \kappa _\mathrm{\Psi })3g(T)gI_2(\beta \kappa _\mathrm{\Phi }),`$
$`g(T)`$ $`=`$ $`gg(T){\displaystyle \frac{2g_\mathrm{\Phi }}{6}}I_2(\beta \kappa _\mathrm{\Phi })g(T){\displaystyle \frac{2g_\mathrm{\Psi }}{6}}I_2(\beta \kappa _\mathrm{\Psi })g_\mathrm{\Phi }(T){\displaystyle \frac{2g}{6}}I_2(\beta \kappa _\mathrm{\Phi })g_\mathrm{\Psi }(T){\displaystyle \frac{2g}{6}}I_2(\beta \kappa _\mathrm{\Psi }),`$ (38)
$``$ $`g(T)gI_3(\beta \kappa _\mathrm{\Phi },\beta \kappa _\mathrm{\Psi }),`$
where we have defined the functions
$`I_1(\beta \kappa _i)`$ $`=`$ $`{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\frac{1}{2\omega _i(๐ช)}\left\{1+2n(\omega _i)+8\beta \omega _in(\omega _i)\left[1+n(\omega _i)\right]\right\}},`$
$`I_2(\beta \kappa _i)`$ $`=`$ $`{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\beta n(\omega _i)\left[1+n(\omega _i)\right]},`$
$`I_3(\beta \kappa _i,\beta \kappa _j)`$ $`=`$ $`{\displaystyle \frac{d^3๐ช}{(2\pi )^3}\left\{\frac{\omega _j}{\omega _j^2\omega _i^2}\left[1+2n(\omega _i)\right]\frac{\omega _i}{\omega _j^2\omega _i^2}\left[1+2n(\omega _j)\right]\right\}}.`$ (39)
One is now in position to investigate how thermal effects on the effective nonrelativistic couplings manifest themselves in phenomena similar to SNR/ISB. First we show in Fig. 7 some representative results for the effective interactions obtained from the solutions of Eq. (38). The tree-level parameters considered here are: $`g_\mathrm{\Phi }=2\times 10^{15}\mathrm{eV}^2,g_\mathrm{\Psi }=10^{16}\mathrm{eV}^2,g=10^{16}\mathrm{eV}^2`$, $`m_\mathrm{\Phi }m_\mathrm{\Psi }=1\mathrm{G}\mathrm{e}\mathrm{V}`$ and $`\kappa _\mathrm{\Phi }=\kappa _\mathrm{\Psi }=1\mathrm{n}\mathrm{e}\mathrm{V}`$. We note that all couplings tend to evolve to zero at very high temperatures as $`T`$ gets closer to $`m`$. So, the apparent instability caused by the fact that $`g_\mathrm{\Psi }(T)`$ as well as $`g_\mathrm{\Phi }(T)`$ could become negative beyond some temperature, $`T_{\mathrm{break}}`$, as suggested by Eqs. (35) and (36), has disappeared completely when the nonperturbative flow of the couplings are considered (using the tree-level parameters given above and from Eqs. (35) and (36), one would get $`T_{\mathrm{break}}0.019`$ eV (or $`220`$ K)).
It is tempting, by looking at Fig. 7, to associate this to a free model at high energies, however we must recall that the nonrelativistic model, Eq. (23), will eventually no longer be valid for such high values of temperature, in which case it should be replaced by the original relativistic model (in any case the model Eq. (23) should of course be regarded as an effective model valid at low energy scales only).
In Fig. 8 we show the equivalent of Eq. (30) for the $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ fields, $`\mathrm{\Delta }_\mathrm{\Phi }^{\mathrm{NR}}(T)`$ and $`\mathrm{\Delta }_\mathrm{\Psi }^{\mathrm{NR}}(T)`$, respectively, given in terms of the temperature dependent bubble resummed couplings for the same tree-level parameters considered above. We note that for the parameters considered $`\mathrm{\Delta }_\mathrm{\Psi }^{\mathrm{NR}}(T)`$ is initially negative and reverse sign at some temperature, indicating that symmetry breaking at high temperatures tend to happen in the $`\mathrm{\Psi }`$ field direction, while the potential in the $`\mathrm{\Phi }`$ direction remains unbroken (actually, the temperature where $`\mathrm{\Delta }_\mathrm{\Psi }^{\mathrm{NR}}(T)`$ crosses zero, is close to the point of a reentrant transition for $`\mathrm{\Psi }`$). $`\mathrm{\Delta }_\mathrm{\Phi }^{\mathrm{NR}}(T)`$ however always remain positive, which then points to no transition in the $`\mathrm{\Phi }`$ direction. These aspects are also clearly seen in the plot, shown in Fig. 9, for the effective one-body term $`\kappa _\mathrm{\Psi }(T)`$, expressed in terms of the effective nonperturbative couplings. In terms of the parameters considered, symmetry breaking is seen to happen at a temperature $`T_{c,\mathrm{\Psi }}^{(SB)}3.4\times 10^4\mathrm{eV}`$ (or $`4`$ K) while the reentrant phase (symmetry restoration) happens at a temperature $`T_{c,\mathrm{\Psi }}^{(SR)}1.4\times 10^2\mathrm{eV}`$ (or $`161`$ K). In between these two temperatures we see a manifestation of an ISB phase. In the $`\mathrm{\Phi }`$ direction there is no symmetry breaking or reentrant phases at any temperature for the parameters considered.
In Fig. 10 we show a phase diagram for the system as a function of the tree-level coupling $`g_\mathrm{\Psi }`$ and the temperature. The thin horizontal line at $`g_\mathrm{\Psi }=10^{16}\mathrm{eV}^2`$ illustrates the reentrant transition (through an inverse symmetry breaking) shown in Fig. 9. All other parameters are the same as considered above. Note that the condition of stability, $`g_\mathrm{\Psi }g_\mathrm{\Phi }>9g^2`$, expressed in terms of the nonperturbative and temperature dependent couplings, is always satisfied at any temperature for the parameters considered for the previous figures (the effective couplings $`g_\mathrm{\Psi }`$ and $`g_\mathrm{\Phi }`$ also remain always positive, as is clear from Fig. 7 and previous discussion).
## IV CONCLUSIONS
We have reviewed how symmetry non-restoration and inverse symmetry breaking may take place, at arbitrarily large temperatures, in multi-field scalar relativistic and nonrelativistic theories. These counter intuitive phenomena appear due to the fact that the crossed interaction can be negative while the models are still bounded from below. We have recalled that, in the relativistic case, SNR/ISB are not a mere artifact of calculational approximations. We have then set to investigate the possible SNR/ISB manifestation and consequences in a nonrelativistic $`U(N_\mathrm{\Phi })\times U(N_\mathrm{\Psi })`$ scalar model of hard core spheres by considering the simplest case, $`N_\mathrm{\Phi }=N_\mathrm{\Psi }=1`$, which may be relevant for condensed matter systems of bosonic atoms or molecules.
Performing a naive perturbative one-loop calculation, which includes only the first thermal contribution to the self-energy, we have shown that, for negative values of the crossed coupling, SNR/ISB can take place like in the relativistic case. However, the manifestation of SNR/ISB in condensed matter systems of hard core spheres seems to be more counter intuitive than in the relativistic case where the model may represent, e.g., the Higgs sector. With this in mind we have investigated the explicit (temperature) running of the nonrelativistic couplings. One first improvement was to evaluate the perturbative one-loop thermal corrections to the couplings which already indicate that SNR/ISB do not seem to happen, at high temperatures, for the nonrelativistic case. Next, we have resummed the bubble like contributions in a nonperturbative fashion. This procedure fixed the instability problem related to the possibility of $`g_i(T)`$ becoming negative as observed in the calculation which considered only the simplest one-loop corrections to the couplings. Our nonperturbative calculation also showed that the phase transitions happening in the nonrelativistic case includes a continuous SB/SR pattern characterized, at intermediate temperatures, by a reentrant, continuous transition. Therefore, we can state as our major result that, contrary to the relativistic case, SNR/ISB does not seem to occur in the nonrelativistic model of hard core spheres. Instead, reentrant like phenomena become possible, as our results have indicated.
Finally, it would be interesting to investigate SNR/ISB in connection with the Bose-Einstein condensation problem and the present work gives some of the ideas and tools needed for this task that we hope to pursue and report in a future publication.
###### Acknowledgements.
The authors were partially supported by Conselho Nacional de Desenvolvimento Cientรญfico e Tecnolรณgico (CNPq-Brazil).
## Appendix A Temperature dependent corrections for the nonrelativistic model
Here we give the main steps used in the evaluation of the thermal masses and couplings for the nonrelativistic model. Similar derivations for the relativistic model can be found e.g. in Refs. roos ; fendley ; MR1 .
The effective one-body terms $`\kappa _i(T)`$ and the effective couplings (two-body terms) for the nonrelativistic model are most easily to obtain directly from a computation of the one-loop effective potential for the model lagrangian (23). The effective one and two-body terms will then be identified with the appropriate derivatives of the effective potential.
As usual in the computation of the one-loop potential, we start by decomposing the fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ in (23) in terms of (constant) background fields (which, without loss of generality, can be taken as real fields) $`\varphi _0`$ and $`\psi _0`$, respectively, and fluctuations $`\varphi `$ and $`\psi `$, which in terms of real components, become
$`\mathrm{\Phi }={\displaystyle \frac{1}{\sqrt{2}}}\left(\varphi _0+\varphi _1+i\varphi _2\right),`$ (40)
$`\mathrm{\Psi }={\displaystyle \frac{1}{\sqrt{2}}}\left(\psi _0+\psi _1+i\psi _2\right).`$ (41)
When substituting Eqs. (40) and (41) in (23) we only need keep the quadratic terms in the fluctuation fields for the computation of the one-loop potential for the background fields $`\varphi _0`$ and $`\psi _0`$. We then obtain the (Euclidean) lagrangian density in terms of $`\varphi _0`$ and $`\psi _0`$,
$`_E`$ $`=`$ $`{\displaystyle \frac{\kappa _\mathrm{\Phi }}{2}}\varphi _0^2+{\displaystyle \frac{g_\mathrm{\Phi }}{4!}}\varphi _0^4+{\displaystyle \frac{\kappa _\mathrm{\Psi }}{2}}\psi _0^2+{\displaystyle \frac{g_\mathrm{\Psi }}{4!}}\psi _0^4+{\displaystyle \frac{g}{4}}\varphi _0^2\psi _0^2`$ (42)
$`+`$ $`{\displaystyle \frac{1}{2}}\chi \widehat{M}\chi +\mathrm{cubic}\mathrm{and}\mathrm{quartic}\mathrm{interaction}\mathrm{terms},`$
where we have defined the vector $`\chi =(\varphi _1,\varphi _2,\psi _1,\psi _2)`$ and $`\widehat{M}`$ is the matrix operator for the quadratic terms in the fluctuations,
$$\widehat{M}=\left(\begin{array}{cccc}\frac{^2}{2m_\mathrm{\Phi }}+\kappa _\mathrm{\Phi }+\frac{g_\mathrm{\Phi }}{2}\varphi _0^2+\frac{g}{2}\psi _0^2& i_\tau & g\varphi _0\psi _0& 0\\ i_\tau & \frac{^2}{2m_\mathrm{\Phi }}+\kappa _\mathrm{\Phi }+\frac{g_\mathrm{\Phi }}{6}\varphi _0^2+\frac{g}{2}\psi _0^2& 0& 0\\ g\varphi _0\psi _0& 0& \frac{^2}{2m_\mathrm{\Psi }}+\kappa _\mathrm{\Psi }+\frac{g_\mathrm{\Psi }}{2}\psi _0^2+\frac{g}{2}\varphi _0^2& i_\tau \\ 0& 0& i_\tau & \frac{^2}{2m_\mathrm{\Psi }}+\kappa _\mathrm{\Psi }+\frac{g_\mathrm{\Psi }}{6}\psi _0^2+\frac{g}{2}\varphi _0^2\end{array}\right).$$
(43)
The partial time derivative in (43) is over Euclidean time: $`_\tau =/\tau `$, $`\tau =it`$. By performing the functional integration in the quadratic fluctuations $`\chi `$, the one-loop effective potential $`V_{\mathrm{eff}}(\varphi _0,\psi _0)`$ obtained from Eq. (42) is given by
$`V_{\mathrm{eff}}(\varphi _0,\psi _0)={\displaystyle \frac{\kappa _\mathrm{\Phi }}{2}}\varphi _0^2+{\displaystyle \frac{g_\mathrm{\Phi }}{4!}}\varphi _0^4+{\displaystyle \frac{\kappa _\mathrm{\Psi }}{2}}\psi _0^2+{\displaystyle \frac{g_\mathrm{\Psi }}{4!}}\psi _0^4+{\displaystyle \frac{g}{4}}\varphi _0^2\psi _0^2+{\displaystyle \frac{1}{2}}\mathrm{ln}det\widehat{M},`$ (44)
where the last term on the rhs of (44) comes from the functional integral over the components of $`\chi `$,
$$\frac{1}{2}\mathrm{ln}det\widehat{M}=\frac{1}{\beta V}D\varphi _1D\varphi _2D\psi _1D\psi _2\mathrm{exp}\left[_0^\beta ๐\tau d^3x\left(\frac{1}{2}\chi \widehat{M}\chi \right)\right],$$
(45)
and $`V`$ is the volume of space. Expressing Eqs. (45) and (43) in the space-time momentum Fourier transform form, we then also obtain
$`{\displaystyle \frac{1}{2}}\mathrm{ln}det\widehat{M}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^3๐ช}{(2\pi )^3}}\mathrm{ln}\{[\omega _n^2+E_\mathrm{\Phi }^2(๐ช)][\omega _n^2+E_\mathrm{\Psi }^2(๐ช)]`$ (46)
$``$ $`g^2\varphi _0^2\psi _0^2[\omega _\mathrm{\Phi }(๐ช)+{\displaystyle \frac{g_\mathrm{\Phi }}{6}}\varphi _0^2+{\displaystyle \frac{g}{2}}\psi _0^2][\omega _\mathrm{\Psi }(๐ช)+{\displaystyle \frac{g_\mathrm{\Psi }}{6}}\psi _0^2+{\displaystyle \frac{g}{2}}\varphi _0^2]\},`$
with
$$E_\mathrm{\Phi }(๐ช)=\sqrt{\left[\omega _\mathrm{\Phi }(๐ช)+\frac{g_\mathrm{\Phi }}{2}\varphi _0^2+\frac{g}{2}\psi _0^2\right]\left[\omega _\mathrm{\Phi }(๐ช)+\frac{g_\mathrm{\Phi }}{6}\varphi _0^2+\frac{g}{2}\psi _0^2\right]},$$
(47)
and
$$E_\mathrm{\Psi }(๐ช)=\sqrt{\left[\omega _\mathrm{\Psi }(๐ช)+\frac{g_\mathrm{\Psi }}{2}\psi _0^2+\frac{g}{2}\varphi _0^2\right]\left[\omega _\mathrm{\Psi }(๐ช)+\frac{g_\mathrm{\Psi }}{6}\psi _0^2+\frac{g}{2}\varphi _0^2\right]}.$$
(48)
From $`V_{\mathrm{eff}}(\varphi _0,\psi _0)`$, Eq. (44), we now define the effective one-body terms as
$$\kappa _\mathrm{\Phi }(T)=\frac{^2V_{\mathrm{eff}}(\varphi _0,\psi _0)}{\varphi _0^2}|_{\varphi _0=0,\psi _0=0},$$
(49)
$$\kappa _\mathrm{\Psi }(T)=\frac{^2V_{\mathrm{eff}}(\varphi _0,\psi _0)}{\psi _0^2}|_{\varphi _0=0,\psi _0=0},$$
(50)
which then gives
$$\kappa _\mathrm{\Phi }(T)=\kappa _\mathrm{\Phi }+\frac{2g_\mathrm{\Phi }}{3}\frac{1}{\beta }\underset{n}{}\frac{d^3๐ช}{(2\pi )^3}\frac{\omega _\mathrm{\Phi }(๐ช)}{\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)}+g\frac{1}{\beta }\underset{n}{}\frac{d^3๐ช}{(2\pi )^3}\frac{\omega _\mathrm{\Psi }(๐ช)}{\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)},$$
(51)
and
$$\kappa _\mathrm{\Psi }(T)=\kappa _\mathrm{\Psi }+\frac{2g_\mathrm{\Psi }}{3}\frac{1}{\beta }\underset{n}{}\frac{d^3๐ช}{(2\pi )^3}\frac{\omega _\mathrm{\Psi }(๐ช)}{\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)}+g\frac{1}{\beta }\underset{n}{}\frac{d^3๐ช}{(2\pi )^3}\frac{\omega _\mathrm{\Phi }(๐ช)}{\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)},$$
(52)
The Eqs. (51) and (52) can also easily be expressed in terms of the free nonrelativistic propagators $`D_\mathrm{\Phi }(\omega _n,๐ช)`$ and $`D_\mathrm{\Psi }(\omega _n,๐ช)`$. The sum over the Matsubara frequencies in (51) and (52) are easily performed by using the identity jackiw ,
$$\frac{1}{\beta }\underset{n}{}\frac{1}{\omega _n^2+\omega ^2}=\frac{1}{\omega }\left[\frac{1}{2}+n(\omega )\right],$$
(53)
where $`n(\omega )`$ is the Bose-Einstein distribution,
$$n(\omega )=\frac{1}{e^{\beta \omega }1}.$$
(54)
The momentum integrals in Eqs. (51) and (52) can be expressed in terms of standard Bose integrals as follows (see for example pathria ). Consider the integral (where $`\omega =๐ช^2/(2m)+\kappa `$ and $`\eta =\beta \kappa `$)
$`{\displaystyle \frac{d^3q}{(2\pi )^3}n(\omega )}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}\left({\displaystyle \frac{2m}{\beta }}\right)^{3/2}{\displaystyle _0^{\mathrm{}}}๐x{\displaystyle \frac{x^2}{e^{x^2+\eta }1}}`$ (55)
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}\left({\displaystyle \frac{2m}{\beta }}\right)^{3/2}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}e^{l\eta }{\displaystyle _0^{\mathrm{}}}๐xx^2e^{lx^2}`$
$`=`$ $`\left({\displaystyle \frac{m}{2\pi \beta }}\right)^{3/2}\mathrm{Li}_{3/2}\left(e^{\beta \kappa }\right),`$
where we used the definition for the polylogarithmic function,
$$\mathrm{Li}_\alpha (z)=\underset{l=1}{\overset{\mathrm{}}{}}\frac{z^l}{l^\alpha }.$$
(56)
Another useful momentum integral that also can be obtained from (55) is
$`{\displaystyle \frac{d^3q}{(2\pi )^3}n(\omega )[1+n(\omega )]}=\left({\displaystyle \frac{m}{2\pi \beta }}\right)^{3/2}\mathrm{Li}_{1/2}\left(e^{\beta \kappa }\right).`$ (57)
We also have the results obtained from the the polylogarithmic functions in the high temperature approximation, $`\kappa T`$, and that are used in the text,
$$\mathrm{Li}_{3/2}\left(e^{\beta \kappa }\right)=\zeta (3/2)2\sqrt{\pi \frac{\kappa }{T}}+๐ช(\kappa /T),$$
(58)
and
$$\mathrm{Li}_{1/2}\left(e^{\beta \kappa }\right)=\sqrt{\pi \frac{T}{\kappa }}\zeta (1/2)+๐ช(\kappa /T).$$
(59)
Using (55) in Eqs. (51) and (52) we obtain the results quoted in the text, Eqs. (27) and (28).
The two-body effective terms are also defined analogously as
$$g_\mathrm{\Phi }(T)=\frac{^4V_{\mathrm{eff}}(\varphi _0,\psi _0)}{\varphi _0^4}|_{\varphi _0=0,\psi _0=0},$$
(60)
$$g_\mathrm{\Psi }(T)=\frac{^4V_{\mathrm{eff}}(\varphi _0,\psi _0)}{\psi _0^4}|_{\varphi _0=0,\psi _0=0},$$
(61)
and
$$g(T)=\frac{^4V_{\mathrm{eff}}(\varphi _0,\psi _0)}{\varphi _0^2\psi _0^2}|_{\varphi _0=0,\psi _0=0}.$$
(62)
From Eqs. (44) and (46) we then obtain for Eqs. (60) - (62) the results
$`g_\mathrm{\Phi }(T)=g_\mathrm{\Phi }`$ $`+`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^3๐ช}{(2\pi )^3}}\{g_\mathrm{\Phi }^2{\displaystyle \frac{1}{\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)}}{\displaystyle \frac{8}{3}}g_\mathrm{\Phi }^2{\displaystyle \frac{\omega _\mathrm{\Phi }^2(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)\right]^2}}`$ (63)
$`+`$ $`3g^2{\displaystyle \frac{1}{\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)}}6g^2{\displaystyle \frac{\omega _\mathrm{\Psi }^2(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)\right]^2}}\},`$
$`g_\mathrm{\Psi }(T)=g_\mathrm{\Psi }`$ $`+`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^3๐ช}{(2\pi )^3}}\{g_\mathrm{\Psi }^2{\displaystyle \frac{1}{\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)}}{\displaystyle \frac{8}{3}}g_\mathrm{\Psi }^2{\displaystyle \frac{\omega _\mathrm{\Psi }^2(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)\right]^2}}`$ (64)
$`+`$ $`3g^2{\displaystyle \frac{1}{\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)}}6g^2{\displaystyle \frac{\omega _\mathrm{\Phi }^2(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)\right]^2}}\},`$
and
$`g(T)=g`$ $`+`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^3๐ช}{(2\pi )^3}}\{{\displaystyle \frac{2}{3}}gg_\mathrm{\Phi }{\displaystyle \frac{1}{\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)}}{\displaystyle \frac{4}{3}}gg_\mathrm{\Phi }{\displaystyle \frac{\omega _\mathrm{\Phi }^2(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)\right]^2}}`$ (65)
$`+`$ $`{\displaystyle \frac{2}{3}}gg_\mathrm{\Psi }{\displaystyle \frac{1}{\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)}}{\displaystyle \frac{4}{3}}gg_\mathrm{\Psi }{\displaystyle \frac{\omega _\mathrm{\Psi }^2(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)\right]^2}}2g^2{\displaystyle \frac{\omega _\mathrm{\Phi }(๐ช)\omega _\mathrm{\Psi }(๐ช)}{\left[\omega _n^2+\omega _\mathrm{\Phi }^2(๐ช)\right]\left[\omega _n^2+\omega _\mathrm{\Psi }^2(๐ช)\right]}}\}.`$
The Eqs. (63) - (65) again can be expressed in terms of the free nonrelativistic propagators $`D_\mathrm{\Phi }(\omega _n,๐ช)`$ and $`D_\mathrm{\Psi }(\omega _n,๐ช)`$, which can then be identified with the corresponding one-loop diagrams that contribute here, depicted in Figs 1, 2 and 3a. All sums over the Matsubara frequencies in (63) - (65) are again evaluated with the help of the identity (53), from which we obtain the results shown in Eqs. (32), (33) and (34). |
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