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warning/0003/hep-ph0003272.html | ar5iv | text | # 1 Introduction
## 1 Introduction
A linear collider with both beams polarized will be an excellent tool not only to discover supersymmetric particles but also to determine the underlying SUSY model. Associated neutralino production $`e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ could be the first kinematically accessible process that may allow a discrimination between the Minimal Supersymmetric Standard Model (MSSM) and its nonminimal extensions. Polarization of both beams will clearly facilitate this challenging task.
An extensive analysis of polarization and spin effects in neutralino production and subsequent decay including the complete spin correlations in the MSSM has been performed in . This study also includes extended SUSY models: the Next-to-Minimal Supersymmetric Standard Model (NMSSM) with an additional Higgs singlet superfield and an $`\mathrm{E}_6`$ inspired model with a Higgs singlet and a new $`\mathrm{U}(1)^{}`$ gauge boson. Significant differences between the MSSM and SUSY models with gauge singlets arise by the neutralino singlet components which do not couple to (scalar) fermions and standard gauge bosons. Within a special scenario, where the masses of the light neutralinos were fixed in the three SUSY models and the lightest neutralino had a large singlino component in the extended models, the cross sections for neutralino production at an $`e^+e^{}`$ linear collider with polarized beams have been discussed in , and complete different polarization asymmetries have been found.
Now the study of the polarization effects is extended to a large region of the parameter space where the lightest neutralino is mainly a singlino. Comparing the longitudinally polarized cross sections we identify the parameter regions where polarization effects may help to specify nonminimal supersymmetric models. Of course for an exact determination of the SUSY parameters and discrimination between SUSY models a precise analysis also of the decay characteristics is indispensable. For a special scenario this was done in , the extension to the whole parameter space is planned.
In the following sections 2 and 3 we introduce the considered SUSY models and show the unpolarized cross section for the neutralino production. Crucial for the understanding of the polarization effects are the couplings of the left and right selectrons to the produced neutralinos that will be discussed in section 4. A comprehensive comparison of the longitudinally polarized cross sections in the MSSM and the extended models in section 5 concludes our contribution.
## 2 Extended SUSY: NMSSM and E<sub>6</sub> model
We shortly present the parameters and the neutralino sectors of the supersymmetric models. In the MSSM the neutralino sector depends on the gaugino mass parameters $`M_2`$ and $`M_1`$, the higgsino mass parameter $`\mu `$ and the ratio of the Higgs vacuum expectation values $`\mathrm{tan}\beta =v_2/v_1`$. In this paper we assume the GUT-relation $`M_1/M_2=5/3\mathrm{tan}^2\theta _W`$ and scan over the $`M`$-$`\mu `$ plane with fixed $`\mathrm{tan}\beta =3`$. The longitudinally polarized cross sections for neutralino production are compared with two extended SUSY models.
The NMSSM is the simplest extension of the MSSM by a Higgs singlet field with vacuum expectation value $`x`$ and hypercharge $`0`$ . The masses and couplings of the five neutralinos depend on $`M_2`$, $`M_1`$, $`\mathrm{tan}\beta `$, $`x`$ and the trilinear couplings $`\lambda `$ and $`\kappa `$ in the superpotential. Within the NMSSM, a light singlino-like neutralino cannot be experimentally excluded by LEP . In order to obtain a light singlino-like neutralino $`\stackrel{~}{\chi }_1^0`$ we choose in the following $`x=1000`$ GeV and $`\kappa =0.01`$. Then nearly independently of the other parameters the $`\stackrel{~}{\chi }_1^0`$ has singlino character and a mass between about 9 and 28 GeV while the masses and mixings of the heavier neutralinos correspond to the MSSM with $`\mu =\lambda x`$.
Models with additional U(1) factors in the gauge group are a further extension of the MSSM beyond the NMSSM. They can be motivated by superstring theory and imply an $`\mathrm{E}_6`$ group as effective GUT group, which is broken to an effective low energy gauge group of rank 5 with one additional $`\mathrm{U}(1)^{}`$ factor. This $`\mathrm{E}_6`$ model contains one new gauge boson $`Z^{}`$ and an extended Higgs sector with one singlet superfield with vacuum expectation value $`x`$ as in the NMSSM . To respect the experimental mass bounds of new gauge bosons the vacuum expectation value of the singlet must be larger than 1 TeV.
The extended neutralino sector in the $`\mathrm{E}_6`$ model contains six neutralinos, one $`\stackrel{~}{Z}^{}`$ gaugino and one singlino $`\stackrel{~}{N}`$ in addition to the MSSM . The $`6\times 6`$ neutralino mass matrix depends on six parameters: the $`\mathrm{SU}(2)_\mathrm{L}`$, U(1)<sub>Y</sub> and $`\mathrm{U}(1)^{}`$ gaugino mass parameters $`M_2`$, $`M_1`$ and $`M^{}`$, $`\mathrm{tan}\beta `$, $`x`$ and the trilinear coupling $`\lambda `$ in the superpotential. With the GUT relation $`M^{}=M_1`$ between the U(1) gaugino mass parameters the four lighter neutralinos always have MSSM-like character if the MSSM parameter $`\mu `$ is identified with $`\lambda x`$ , and the two heaviest neutralinos have large $`\stackrel{~}{Z}^{}`$ and $`\stackrel{~}{N}`$ components. Assuming $`|M^{}|x,M_1`$, however, the lightest neutralino can be a nearly pure singlino . Then $`\stackrel{~}{\chi }_2^0,\mathrm{},\stackrel{~}{\chi }_5^0`$ have MSSM-like character and the $`\stackrel{~}{\chi }_6^0\stackrel{~}{Z}^{}`$ has a mass $`𝒪(M^{})`$. Such scenarios where the spectrum of the lighter neutralinos is similar to the NMSSM will be discussed in the following.
Table 1 gives an overview over the described fixed parameters, while the remaining parameters $`M_2`$ and $`\mu `$ (or $`\lambda x`$) were scanned in the region $`0M_2500`$ GeV, $`500\mu (\lambda x)500`$ GeV observing the experimental bounds from the unsuccessful chargino and neutralino search at LEP2 . Also shown are the masses and singlino components of the lightest neutralino in the extended models which do not extensively vary over the scanned parameter region. The selectron masses $`m_{\stackrel{~}{e}_L}=176`$ GeV and $`m_{\stackrel{~}{e}_R}=132`$ GeV are motivated by the DESY/ECFA reference scenario and are fixed in all three models, for comparison.
## 3 Unpolarized cross sections
In this section we address the question if a neutralino with a dominant singlino component can be directly produced with a sufficient cross section at a linear collider.
The neutralino production cross section in the MSSM is well known . Fig. 1(a) shows the contour lines for the unpolarized cross section $`\sigma ^{(00)}`$ of $`e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ at a c.m.s. energy $`\sqrt{s}=500`$ GeV in the $`M_2`$-$`\mu `$ parameter space. Depending on the masses and mixing characters of the light neutralinos it takes values up to 200 fb over most of the parameter space. Note the discontinuities (bold lines) for both positive and negative $`\mu `$ where $`m_{\stackrel{~}{\chi }_2^0}=m_{\stackrel{~}{\chi }_3^0}`$ and therefore the mixing character of $`\stackrel{~}{\chi }_2^0`$ changes .
In the NMSSM with a light singlino-dominated neutralino $`\stackrel{~}{\chi }_1^0`$ the unpolarized cross section (Fig. 2(a)) is significantly smaller since the singlino component does not couple to the $`Z`$ boson and the selectrons. It is only the small contribution from the other components of $`\stackrel{~}{\chi }_1^0`$ that accounts for the cross section. Nevertheless one gets even for a neutralino with a singlino component larger than 90 % production cross sections up to 30 fb which are expected to be clearly visible at a high luminosity linear collider ($`=500/`$fb) . The use of polarized beams can still enhance these cross sections as we discuss in the following sections.
In the $`\mathrm{E}_6`$ model we explicitly study the dependence of the unpolarized cross sections on the parameter $`M^{}`$ in Figs. 3(a) – 5(a). For the large negative value $`M^{}=50`$ TeV (Fig. 3(a)) the masses and the mixings of the light neutralinos are similar as in the NMSSM. Therefore one also obtains similar cross sections up to 9 fb. For positive $`M^{}`$ (Figs. 4(a), 5(a)), however, the singlino component is larger than $`99.4\%`$ in the whole parameter space. Thus the cross sections are reduced to maximum values of about 0.2 fb for $`\lambda x<0`$ and 0.5 fb for $`\lambda x>0`$ leading to at least 100 respective 250 events at the expected high luminosity linear collider.
## 4 Neutralino couplings to selectrons and electrons
If apart from the singlino component which does not contribute to the production process both produced neutralinos $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ have dominant gaugino components, the contribution of $`Z`$ exchange can be neglected for energies above the $`Z`$ peak. Then for longitudinally polarized beams the cross section consists of two terms describing the exchange of left and right selectrons $`\sigma =\sigma _{\stackrel{~}{e}_L}+\sigma _{\stackrel{~}{e}_R}`$. The structure of $`\sigma _{\stackrel{~}{e}_{L/R}}`$ is
$$\sigma _{\stackrel{~}{e}_{L/R}}=(f_{e1}^{L/R}f_{e2}^{L/R})^2\left[(1P^{}P^+)(P^{}P^+)\right]\stackrel{~}{\sigma }(s,m_{\stackrel{~}{e}_{L/R}},m_{\stackrel{~}{\chi }_{1/2}^0})$$
(1)
where $`P^{}`$ and $`P^+`$ are the longitudinal polarizations of the electron and positron beam, respectively, and $`\stackrel{~}{\sigma }`$ is a function of the beam energy and the particle masses. For selectron masses of the same order of magnitude $`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}`$ the polarization behavior is significantly determined by the ratio
$$r_f=\frac{(f_{e1}^Rf_{e2}^R)^2}{(f_{e1}^Lf_{e2}^L)^2}$$
(2)
of the couplings of the produced neutralinos to the left handed and right handed electrons and selectrons $`f_{ei}^{L/R}`$ . The effects of different masses, especially the extreme cases $`m_{\stackrel{~}{e}_L}()`$ $`m_{\stackrel{~}{e}_R}`$, are shortly mentioned in the following section.
In this contribution we consider longitudinal beam polarizations of $`P^{}=\pm 0.85`$ for electrons and $`P^+=\pm 0.6`$ for positrons . The solid lines in Figs. 1(b), 4(b) and 5(b) show the contour for $`r_f=1`$. In the NMSSM and the $`\mathrm{E}_6`$ model with negative $`M^{}`$ (Figs. 2(b) and 3(b)) it is always $`r_f>1`$. This contour line clearly separates the gray shaded regions with maximum cross sections for polarization configurations $`(+)`$ (right handed polarized electrons and left handed polarized positrons) or $`(+)`$ (vice versa) in the MSSM (Fig. 1(b)) in the gaugino region ($`M_22|\mu |`$). In the $`\mathrm{E}_6`$ model (Figs. 4(b) and 5(b)) the situation is more complicated because the higgsino components of the $`\stackrel{~}{\chi }_1^0`$ are larger than the gaugino components in the whole parameter space. Nevertheless the contour $`r_f=1`$ approximately describes the polarization behavior. The polarization effects in detail are discussed in the following section.
## 5 Polarized cross sections
We compare the longitudinally polarized cross sections $`\sigma ^{(+)}`$ and $`\sigma ^{(+)}`$ where $`(+)`$ and $`(+)`$ denotes the polarization configuration $`(P^{}=+0.85,P^+=0.6)`$ and $`(P^{}=0.85,P^+=+0.6)`$, respectively. These polarization asymmetries may help to distinguish between the SUSY models. Figs. 1(b) – 5(b) present the regions where the different polarizations $`(+)`$ or $`(+)`$ lead to maximum cross sections. Polarization configurations with both beams polarized in the same direction or only one beam polarized are always smaller.
The following remarks apply to all models: If the gaugino-character of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ is not too small (which is the case in the MSSM for $`M_22|\mu |`$) $`\sigma ^{(+)}`$ dominates if the gaugino-coupling of the left selectron is larger than that of the right one, $`r_f<1`$, and the mass difference between the selectrons is rather small. Analogously for $`r_f>1`$, $`\sigma ^{(+)}`$ turns out to be the dominating cross section.
This situation may change for other selectron masses. If $`m_{\stackrel{~}{e}_L}()m_{\stackrel{~}{e}_R}`$ the exchange of the left (right) selectron in the neutralino production is strongly suppressed and $`\sigma ^{(+)}`$ ($`\sigma ^{(+)}`$) dominates in the whole gaugino region, cf. .
If the gaugino components of the produced neutralinos can be neglected, the ratio of the neutralino-selectron-electron couplings plays no role and the asymmetry in the MSSM and NMSSM
$$\frac{\sigma ^{(+)}\sigma ^{(+)}}{\sigma ^{(+)}+\sigma ^{(+)}}\frac{|P^{}|+|P^+|}{1+|P^{}||P^+|}\frac{L^2R^2}{L^2+R^2}>0$$
(3)
depends only on the electron-$`Z`$ couplings $`L=1/2+\mathrm{sin}^2\theta _W`$ and $`R=\mathrm{sin}^2\theta _W`$. So $`\sigma ^{(+)}`$ is generally largest in the MSSM higgsino region $`M_22|\mu |`$, whereas in the NMSSM the gaugino components of $`\stackrel{~}{\chi }_2^0`$ always cause $`\sigma ^{(+)}>\sigma ^{(+)}`$ for $`M_2<500`$ GeV. In the $`\mathrm{E}_6`$ model formula (3) becomes more complicated due to $`Z^{}`$ exchange.
Because of $`r_f>1`$ in the NMSSM and in the $`\mathrm{E}_6`$ model with $`M^{}=50`$ TeV $`\sigma ^{(+)}`$ dominates in the whole parameter space. In the $`\mathrm{E}_6`$ model with positive $`M^{}`$ both cases $`r_f>1`$ and $`r_f<1`$ occur. Thus $`\sigma ^{(+)}`$ as well as $`\sigma ^{(+)}`$ can be largest. The form of the respective parameter regions crucially depends on the value of $`M^{}`$ (compare Figs. 4(b) and 5(b)) and is only approximately given by the contours $`r_f=1`$ because here the higgsino character of $`\stackrel{~}{\chi }_1^0`$ is always rather large.
In all models the cross sections for both beams polarized are more than two times larger than the unpolarized cross sections in a large fraction of the parameter space. To describe this fact we define the ratio
$$r^{(\pm )}=\frac{\sigma ^{(\pm )}}{\sigma ^{(00)}}.$$
(4)
In the MSSM for $`M_22|\mu |`$ (gaugino-region) one nearly always obtains $`r^{(+)}>2`$ or $`r^{(+)}>2`$ in the respective regions. In the NMSSM and the $`\mathrm{E}_6`$ model with negative $`M^{}`$ it is always $`r^{(+)}>2`$ and even $`r^{(+)}>2.7`$ (NMSSM) or $`2.9`$ ($`\mathrm{E}_6`$ model) in a large fraction of parameter space. And also for positive $`M^{}`$ $`r^{(+)}>2`$ or $`r^{(+)}>2`$ holds in large parameter regions. Thus polarization of both beams is an important tool to increase the event rates.
Finally we describe the parameter regions ($`0M_2500`$ GeV, $`500\text{GeV}\mu ,\lambda x500`$ GeV) where significant differences between the SUSY models arise in neutralino production at a linear collider if both beams are polarized. In the extended models with a light singlino-like neutralino ($`m_{\chi _1^0}30`$ GeV) the polarization configuration $`\sigma ^{(+)}`$ dominates for $`\lambda x>0`$, while in the MSSM $`\sigma ^{(+)}`$ is largest apart from a narrow band between $`M_2\mu `$ and $`M_22\mu `$. If the mass of the singlino-like neutralino, however, increases also the $`\mathrm{E}_6`$ model with $`M^{}=35`$ TeV allows $`\sigma ^{(+)}>\sigma ^{(+)}`$ in some parameter regions. In the MSSM, NMSSM and $`\mathrm{E}_6`$ model with negative $`M^{}`$ the situation for $`\mu `$ or $`\lambda x<0`$ is similar to $`\mu ,\lambda x>0`$. In the $`\mathrm{E}_6`$ model with positive $`M^{}`$ the polarization asymmetries strongly depend on $`M^{}`$. For $`M^{}=50`$ TeV with a lighter singlino $`\sigma ^{(+)}`$ dominates for $`M_2\lambda x`$ and $`M_2\lambda x`$ contrary to the MSSM, whereas for $`M^{}=35`$ TeV $`\sigma ^{(+)}`$ dominates only for large $`M_2`$.
## 6 Conclusion
We summarize our results in two points:
* A high luminosity linear collider is sufficient for the direct production of neutralinos with dominant singlino components in extended supersymmetric models. Polarization of both beams is an important tool to increase the event rates by a factor of 2 or 3.
* The Minimal Supersymmetric Standard Model and extended supersymmetric models show a significant different polarization behavior for neutralino production in wide regions of the parameter space.
## Acknowledgment
We thank G. Moortgat-Pick for many valuable discussions. This work was supported by the Deutsche Forschungsgemeinschaft (DFG) under contract No. FR 1064/4-1, by the Bundesministerium für Bildung und Forschung (BMBF) under contract No. 05 HT9WWA 9 and by the Fonds zur Förderung der wissenschaftlichen Forschung of Austria, project No. P13139-PHY. |
warning/0003/physics0003043.html | ar5iv | text | # Ground and excited states of the hydrogen negative ion in strong magnetic fields
## I Introduction
The term “strong field” characterizes a situation for which the Lorentz force is of the order of magnitude or greater than the Coulomb binding force. For a hydrogen atom in the ground state the corresponding field strength cannot be reached in the laboratory, but only in astrophysical objects like white dwarfs ( B$``$ $`10^2`$$`10^5`$T) or neutron stars ( B$``$ $`10^7`$$`10^9`$T). Astrophysicists possess therefore a vivid interest in the behavior and properties of matter in strong magnetic fields: theoretically calculated data of magnetized atoms can be used for the determination of the decomposition and magnetic field configuration of astrophysical objects . On the other hand the strong magnetic field regime is accessible in the laboratory if one considers highly excited Rydberg states of e.g. atoms .
In solid state physics donor states in semiconductors with parabolic conduction bands are systems which possess a Hamiltonian equivalent to the one of hydrogen within an effective mass approximation. Due to screening effects the Coulomb force is much weaker than in the case of hydrogen. The regime where the ground state of the system is dominated by magnetic forces can therefore be reached for certain semiconductors in the laboratory. As an example we mention GaAs for which the effective mass is $`m^{}=0.067m_e`$ and the static dielectric constant $`ϵ_s=12.53ϵ_0`$. Since the Hamiltonian of the atomic ion and the negative donor are connected through a scaling transformation the values for the energies given in the present work hold for both systems equally. The reader should however keep in mind that they are given in differently scaled units.
Apart from the above atoms and molecules in strong magnetic fields are also of interest from a pure theoretical point of view. Due to the competition of the spherically symmetric Coulomb potential and the cylindrically symmetric magnetic field interaction we encounter a nonseparable, nonintegrable problem. Perturbation theory, which is possible in the weak and in the ultrastrong field regime, breaks down in the intermediate field regime. It is therefore necessary to develop new techniques to solve such problems. The neutral hydrogen atom in a strong magnetic field is now understood to a high degree (see and references therein). Recently Kravchencko has published an “exact” solution which provides an infinite double sum for the eigenvalues . With the presented method all energy values of bound states could in principle be calculated to arbitrary precision.
For two electron atoms the situation is significantly different. The problems posed by the electron-electron interaction and the non-separability on the one-particle level have to be solved simultaneously, which is much harder. The H<sup>-</sup> ion provides an additional challenge since correlation plays an important role for its binding properties. Without a field it possesses only one bound state . In the presence of a magnetic field and for the assumption of an infinitely heavy nucleus it could be shown that there exists an infinite number of bound states. For laboratory field strengths these states are, due to the binding mechanism via a one dimensional projected polarization potential, very weakly bound . Some finite nuclear mass effects can be included via scaling relations. However, the influence of the center of mass motion has not been investigated in detail so far. In the present work we assume an infinitely heavy nucleus which represents a good approximation for the slow H<sup>-</sup> atomic ion in strong magnetic fields and describes simultaneously the situation of negatively charged donors D<sup>-</sup> in the field. Relativistic corrections were neglected since they are assumed to be small compared to the electron detachment energy of the system. We will use in the following the spectroscopic notation $`{}_{}{}^{2S+1}M`$ for the electronic states of the ion where $`M`$ and $`S`$ are the total magnetic and spin quantum numbers. Since states with negative z-parity are not considered here we omit the corresponding label in our notation (see also section II A).
Many authors have tackled the quantum mechanical problem of H<sup>-</sup> in a strong magnetic field. One of the first, who pursued a variational approach to this problem, were Henry et al.. They give first qualitative insights into the weak and intermediate field regime. Mueller et al. qualitatively described the strong field ground state $`{}_{}{}^{3}(1)`$ and the $`{}_{}{}^{1}0`$ state for high fields ($`\gamma 4`$ to $`\gamma \mathrm{20\hspace{0.17em}000}`$, where $`\gamma =1`$ a. u. corresponds to $`2.355410^5`$T).
Larsen has published a number of papers on this problem . On the one hand he created very simple and physically motivated trial functions with only a small number of variational parameters. On the other hand his energies were “state of the art” in variational calculations for a long time. In he provides binding energies of the lowest $`{}_{}{}^{1}0`$ state in the field regime $`\gamma =05`$ and of the $`{}_{}{}^{3}(1)`$ state in the regime $`\gamma =03`$. He also presents figures showing the binding energies of the singlet and triplet state for $`M=2`$ and $`M=3`$. Later he presents total and electron detachment energies for the lowest $`{}_{}{}^{1}0`$, $`{}_{}{}^{3}(1)`$ and $`{}_{}{}^{3}(2)`$ state in the high field regime. More specifically the regime $`\gamma =20\mathrm{1\hspace{0.17em}000}`$ for the $`{}_{}{}^{3}(1)`$ state and $`\gamma =20200`$ for the other states were investigated. Furthermore Park and Starace provided upper and lower bounds for energies and binding energies of the ground state $`{}_{}{}^{1}0`$ for weak fields.
In the nineties several authors improved the accuracy of the binding energies and total energies by new techniques. Vincke and Baye report total ionization energies for the lowest singlet and triplet states with $`M=0,1`$ and $`2`$ for a few field strengths in the regime $`\gamma =4400`$. They are to our knowledge the first who reported that the $`{}_{}{}^{1}(1)`$ state becomes bound for sufficiently high field strengths and realized that the $`{}_{}{}^{1}(2)`$ state is slightly stronger bound than the corresponding triplet state in the high field regime. Larsen and McCann present in one-particle binding energies for the $`{}_{}{}^{1}0`$ state in the broad magnetic field regime $`\gamma =0200`$. In the same authors consider furthermore the singlet and triplet states of $`M=1,2`$. The triplet states are calculated for $`\gamma =0.5200`$, the $`{}_{}{}^{1}(1)`$ state in the field regime $`\gamma =55\mathrm{2\hspace{0.17em}000}`$ and the $`{}_{}{}^{1}(2)`$ state is calculated for a few field strength in the range $`\gamma =1100`$. Blinowski and Szwacka have subsequently used a Gaussian basis set, similar to the one used in our calculation. They present results for the $`{}_{}{}^{1}0`$ state, which are less accurate than those of ref. .
We also mention some Hartree–Fock calculations: very early Virtamo has investigated the ground state energies from $`\gamma 20`$ to $`\gamma \mathrm{20\hspace{0.17em}000}`$. Thurner et al. (results published in ) have calculated triplet states for $`M`$=$`1`$,$`2`$ and $`3`$ for many field strength in the broad range $`\gamma =210^4210^3`$. However since they use spherical wave functions for weak fields and cylindrical ones for high fields, there remains a gap of inaccurate results in the intermediate field regime.
In the present investigation we provide lower variational energies and higher one-particle binding energies for the atomic H<sup>-</sup> problem and respectively the negatively charged donor center D<sup>-</sup> problem in a strong magnetic field compared to all other published data sofar. An exception is the field free situation: the calculation by Pekeris gives $`0.52775`$ a.u. for the ground state binding energy whereas we obtain $`0.5275488`$ a.u. Clearly the field-free situation is much better understood than the case of a strong field.
The paper is organized as follows: in section II we consider the symmetries of the Hamiltonian and the basis set we use in our calculations. In section III we will report on the strategy we employed for the selection of basis functions in order to obtain accurate results. Section IV contains the discussion of our results and a comparison with the literature.
## II Hamiltonian, Symmetries and basis set
### A Hamiltonian and Symmetries
In the following we assume an infinite nuclear mass (fixed donor). The magnetic field is chosen to point along the z-direction. The nonrelativistic Hamiltonian takes in atomic units the form
$$H=H_1+H_2+\frac{1}{\left|𝒓_1𝒓_2\right|}$$
(1)
with
$$H_i=\frac{1}{2}𝒑_i^2+\frac{1}{2}\gamma l_{z_i}+\frac{\gamma ^2}{8}\left(x_i^2+y_i^2\right)\frac{1}{\left|𝒓_i\right|}+\gamma s_{zi}.$$
(2)
The Hamiltonian is splitted in its one-particle operators, where $`1/2\gamma l_{z_i}`$ is the Zeeman term, $`\gamma ^2/8\left(x_i^2+y_i^2\right)`$ is the diamagnetic term, $`1/\left|𝒓_i\right|`$ is the attractive Coulomb interaction with the nucleus (donor) and $`\gamma s_{zi}`$ the spin Zeeman term (we take the g-factor equal 2). The two-particle operator $`1/\left|𝒓_1𝒓_2\right|`$ represents the repulsive electron-electron interaction.
The Hamiltonian (1) possesses four independent symmetries and associated quantum numbers: the total spin $`S^2`$, the total z-projection of the spin $`S_z`$, the z-component of the total angular momentum $`M`$ and the total z-parity $`\mathrm{\Pi }_z`$ (parity is also conserved but not a further independent symmetry).
### B One-particle basis set
For our calculation we use an anisotropic Gaussian basis set, which has been put forward by Schmelcher and Cederbaum in ref. , for the purpose of investigating atoms and molecules in strong magnetic fields. It has already successfully been applied to helium , H$`{}_{}{}^{+}{}_{2}{}^{}`$ and H<sub>2</sub> .
Adapted to the problem discussed here this one-particle basis set for the spatial part reads in the cylindrical coordinates as follows
$$\mathrm{\Phi }_i(\rho ,\varphi ,z)=\rho ^{n_{\rho _i}}z^{n_{z_i}}e^{\alpha _i\rho ^2\beta _iz^2}\mathrm{exp}(im_i\varphi ).$$
(3)
These functions are eigenfunctions of the symmetry operations of the one-particle Hamiltonian $`H_i`$, i.e. eigenfunctions of $`l_z`$ and $`\pi _z`$. The additional parameters $`n_{\rho _i}`$ and $`n_{z_i}`$ obey the following restrictions:
$`n_{\rho _i}=|m_i|+2k_i;`$ $`k_i=0,1,2,\mathrm{}`$ $`\text{and}m_i=\mathrm{},2,1,0,1,2,\mathrm{}`$ (4)
$`n_{z_i}=\pi _{z_i}+2l_i;`$ $`l_i=0,1,2,\mathrm{}`$ $`\text{and}\pi _{z_i}=0,1.`$ (5)
The exponents $`\alpha _i`$ and $`\beta _i`$ serve as positive, nonlinear variational parameters. Due to these parameters, the one-particle functions are flexible enough to be adapted to the situation of an arbitrary field strength: in the weak magnetic field regime a basis set with an almost isotropic choice of parameters $`\alpha _i\beta _i`$ describes the slightly perturbed spherical symmetry. For very high magnetic fields it is appropriate to choose $`\alpha =\gamma /4`$ since $`\rho ^{|m_i|}\mathrm{exp}(\gamma /4\rho ^2)`$ yields the $`\rho `$-dependence of the lowest Landau level for a given magnetic quantum number. The $`\beta _i`$ will be well tempered in a wide region. In the intermediate field regime the basis is composed of functions with certain magnetic field dependent sets of $`\{\alpha _i,\beta _i\}`$ which mediate the extreme cases. The optimal choice is found by searching the set of $`\{\alpha _i`$, $`\beta _i\}`$ which yields the lowest eigenvalues of the one-particle Hamiltonian. The parameters $`\{\alpha _i`$,$`\beta _i\}`$ are successively optimized using the pattern search algorithm. In this manner we have optimized up to five excited states in every symmetry subspace. The starting values for the parameters $`\{\alpha _i`$, $`\beta _i\}`$ have to be chosen very carefully to find a deep local or even the global minimum. Since the search in this high dimensional space is very time consuming, an optimal choice of the $`k_i`$ and $`l_i`$ is crucial: for every new $`k_i`$, $`l_i`$ configuration a new optimization procedure has to be started. The resulting binding energies for the neutral hydrogen atom were identical to 7 – 9 digits with the one given in for almost all field strengths for the ground state and 5 – 7 digits were recovered for states with higher magnetic quantum number $`|m_i|`$.
We point out that Blinowski and Szwacka have used a similar basis set, but without the monomers $`\rho ^{2k_i}`$ and $`z^{2l_i}`$. The additional monomers however decisively enhance the flexibility and accuracy of the calculations.
### C Two-particle configurations
As a next step we build two-particle configurations from our optimized one-particle basis set and represent the Hamiltonian (1) in this configuration space. This is done for each total symmetry $`(S^2,\mathrm{\Pi }_z,L_z)`$ separately. The corresponding spectrum of H<sup>-</sup> is then obtained by diagonalizing the Hamiltonian matrix. We hereby use all possible excited two-particle configurations constructed from our optimized one-particle basis set, i.e. our approach is a full configuration interaction method (full CI). The two-particle functions are constructed from the one-particle functions by selecting combinations for $`m_i+m_j=M`$ and $`\pi _{zi}+\pi _{zj}=\mathrm{\Pi }_z`$. The spin part can be trivially separated. Due to the antisymmetrization of the spatial wave function the configuration space of the triplet states is slightly smaller than that of the singlet states since for triplet configurations there are no combinations with $`i=j`$.
As our basis set is not orthogonal we have to solve a finite-dimensional generalized real symmetric eigenvalue problem
$$(\underset{¯}{\underset{¯}{H}}E\underset{¯}{\underset{¯}{S}})\underset{¯}{c}=0$$
(6)
where $`\underset{¯}{\underset{¯}{H}}`$ is the matrix representation of the Hamiltonian and $`\underset{¯}{\underset{¯}{S}}`$ the overlap matrix. The resulting energies $`E`$ are strict upper bounds to the exact eigenvalues in the given subspace of symmetries.
Some technical remarks concerning the calculation of the matrix elements are in order. All matrix elements can be evaluated analytically. With the exception of the electron-electron integrals all expressions can be calculated very rapidly. The electron-electron integrals, however, deserve a special treatment: through a combination of transformation techniques as well as analytical continuation formulae for the series of involved transcendental functions their representation has been simplified enormously (for details see ref. and in particular ). It is due to this extremely efficient implementation of the electron-electron integral that large basis sets of the order of $`25004000`$ could be used in the present work to perform CI calculations for many field strengths.
## III Selection of the basis functions
Since the single bound state in the absence of the external field is bound only due to correlation, and all the other states in the presence of the magnetic field are only weakly bound, it is very important to include correlation by a proper choice of the one-particle basis functions building up the two-particle configurations. For the $`M=0`$ singlet state this was achieved by selecting one-particle basis functions not only with $`m_1=m_2=0`$ but also with $`m_1=m_20`$. This allows one to describe the angular correlation which is particular important for the $`{}_{}{}^{1}0`$ state. In general the enhanced binding properties of negative ions in the presence of a magnetic field are due to a balanced competition of the different interactions. On the one hand the confinement due to the magnetic field raises the kinetic energy and the electrostatic repulsion due to the electron-electron interaction. These effects tend to lower the binding energy. On the other hand the confinement raises the nuclear attraction energy, the exchange energy and to some extent also the correlation energy which tend to enhance the binding energy. Of course one has to distinguish between, for example, the $`{}_{}{}^{1}0`$ state whose binding properties are dominated by correlation effects and the excited bound states with nonzero magnetic quantum numbers which possess a significant contribution to their binding energy through exchange effects and due to the occupation of the series of tightly bound hydrogenic orbitals $`1s`$, $`2p_1`$, $`3d_2`$,…etc.
For the description of the lowest states with $`|M|>0`$ an effective one-particle picture can be employed : the hydrogen negative ion consists of a tightly bound core electron with magnetic quantum number zero and a significantly less bound electron which carries the magnetic quantum number of the ion. The core electron is then described by one-particle basis functions with $`m_1=0`$. The outer electron is described by one-particle functions with $`m_2=M`$ in order to take into account the fact that it is weakly bound and thus spatially extended. In order to go beyond this effective one-particle picture we used, similar to the case $`M=0`$ one-particle functions with other magnetic quantum numbers to obtain in particular the correlation behavior.
The above picture is not valid for the tightly bound states in the high field regime: the number of functions with different magnetic quantum numbers can be reduced as we increase the field strength. This reduction in the number of basis functions is also suggested by the occurrence of linear dependencies for strong fields. The extent of this reduction can be seen from the fact that the number of two-particle basis functions drops from $`\mathrm{4\hspace{0.17em}000}`$ for $`\gamma =0`$ to less than $`\mathrm{3\hspace{0.17em}000}`$ for $`\gamma =4000`$ for the $`{}_{}{}^{1}0`$ state.
## IV Results and discussion
As already mentioned the H<sup>-</sup> ion possesses only one bound state in the absence of the magnetic field . Turning on the field it has been shown that there exists (for infinite nuclear mass) for any nonzero field an infinite number of bound states. The corresponding proof relies on the physical picture that the external electron is for weak fields far from the neutral atomic core and experiences therefore to lowest order a polarization potential due to the induced dipole moment of the core. Perpendicular to the field the motion of the external electron is dominated by the field and it occupies approximately Landau orbitals whereas parallel to the field it is weakly bound due to the projection of the mentioned polarization potential on the Landau orbitals which yields an one-dimensional binding along the field. For typical strong laboratory fields the corresponding binding energies are of the order of $`10^6`$ eV for the hydrogen atom negative ion and significantly larger for more electron atoms with a larger polarizability. To investigate theses states in the weak field regime goes clearly beyond the feasibility of the present method. Instead we will investigate a number of states, starting from the value of the field strength for which they become significantly bound, which means that the outer electron is already relatively close to the core and possesses a binding energy of at least a few meV. Clearly in that case the picture of the polarization potential is no more valid since exchange and correlation effects rule the binding properties of the ion. Within our approach we could find one bound state for each negative magnetic quantum number of the ion considered ($`3M0`$) for both singlet and triplet states, except the $`{}_{}{}^{3}0`$ state, which is unbound. Their behavior has been studied for the complete range of field strengths $`0.01\gamma 4000`$. The one bound state of the H<sup>-</sup> ion in the absence of the field represents, in the above sense, an exception since it is already significantly bound without the field. All these states possess positive z-parity and *no bound states could be found for negative z-parity*.
### A Threshold energies
The electron detachment energy is defined to be the energy we need to remove one electron from the atom without changing the quantum numbers of the total system. The corresponding lowest threshold energy $`E_T`$ for the H<sup>-</sup> ion can be expressed as:
$$E_T=\frac{\gamma }{2}\left(|M|+M+2+g_eM_s\right)I(\text{H})$$
(7)
where $`I(`$H$`)`$ is the binding energy of the ground state of the neutral hydrogen atom in a magnetic field. The term $`\gamma /2(|M|+M+2)`$ is the energy of an electron in the lowest Landau level with magnetic quantum number $`m=M`$ where the spin part is omitted. This means that the free electron carries the whole angular momentum of the state. For magnetic quantum numbers $`M0`$ the threshold energy $`E_T`$ is independent of the angular momentum $`M`$, i.e. there is only a singlet and a triplet threshold. The threshold energy is then $`E_T=\gamma I(`$H$`)`$ for singlet states and $`E_T=I(`$H$`)`$ for triplet states. We denote the electron detachment energy by $`I(H^{})`$ which is given by $`I(`$H$`{}_{}{}^{})=E_TE_{tot}`$ where $`E_{tot}`$ is the total energy of the considered state of H<sup>-</sup>.
### B Total, electron detachment and transition energies
Before we discuss the individual states and their properties let us describe some general features of the states considered here. The total energy of the singlet states is monotonically increasing with increasing field strength. This fact is caused by the increase of the field-dependent kinetic energy. In contrast to this the total energy of the triplet states is monotonically decreasing with increasing field strengths. This is a consequence of the additional spin Zeeman term (we consider here only the $`S_z=1`$ component of the spin triplet states). The electron detachment energies are monotonically increasing with increasing field strength for all states considered here, i.e. both singlet and triplet states. This has to be seen in view of the above-mentioned fact that the zero-point kinetic (Landau) energy of the electrons is raised in the presence of the magnetic field and therefore the threshold energy for loosing one electron is raised in the same way.
For the $`{}_{}{}^{1}0`$ state the total energy raises from $`0.52754875`$ at $`\gamma =0`$ to $`3986.49870`$ at $`\gamma =4000`$. This state is the most tightly bound state for all field strengths. The detachment energy increases from $`0.027549`$ a.u. at $`\gamma =0`$ to $`2.29805`$ a.u. at $`\gamma =4000`$. There are two reason which give rise to the fact, that this state is the most tightly bound one. On the one hand the electrons are in this state much closer to the nucleus than in other states. This increases the binding due to the attractive nuclear potential energy. On the other hand correlation has an important impact on the binding energy. Both effects are reinforced with increasing field strength as the electrons become more an more confined in the x-y plane perpendicular to the magnetic field. These effects overcome the influence of the static electron-electron repulsion. The total energies and the detachment energies of the $`{}_{}{}^{1}0`$ state are presented in table I. It can be seen that the detachment energies for this most tightly bound state could be improved by 1-2% for all field strengths compared to the existing literature. This is not correct for a vanishing field, where much more efficient basis sets like the Hylleraas basis set are available. For numerical reasons the relative accuracy for the detachment energies is largest in the intermediate field regime.
The $`{}_{}{}^{3}0`$ state is not bound for all considered field strengths. This can be understood in an effective particle picture as follows: for triplet states the spatial two-particle wave function is antisymmetric with respect to particle exchange and therefore the two particles have to occupy different spatial orbitals, i.e. we are exclusively dealing with excited configurations. For $`M0`$ it is (see later) possible to obtained tightly bound triplet states in a strong magnetic field by occupying different orbitals of the hydrogenic series ($`1s,2p_1,3d_2,\mathrm{}`$) which yields the one-particle excited configurations of the type $`1s2p_1,1s3d_2,\mathrm{}`$. For the case of the $`{}_{}{}^{3}0`$ state, however, we have $`M=0`$ and only configurations constructed from pairs of two orbitals with $`(m,m)`$ are allowed which are either of doubly excited character ($`m0`$) or a singly excited configuration with $`m=0`$. Therefore no magnetically tightly bound configurations are allowed for the $`{}_{}{}^{3}0`$ state which illuminates its unbound character for any field strength. All singlet and triplet electron detachment energies of all the considered bound states are presented also graphically: Figure 1 shows the singlet detachment energies and figure 2 the corresponding energies for the triplet states.
It is important to mention that the global ground state of the ion undergoes a crossover with respect to its symmetry with increasing field strength. For weak fields the $`{}_{}{}^{1}0`$ state is the ground state of the system, whereas in strong fields the $`{}_{}{}^{3}(1)`$ state becomes the ground state which was first shown in ref.. This is caused by the spin Zeeman term, which lowers the total energy of the triplet states. The crossover takes place at $`\gamma _c0.05`$ which corresponds to approximately $`10^4`$ T for the H<sup>-</sup> ion. The $`{}_{}{}^{3}(1)`$ state is very weakly bound when it becomes the ground state (at $`\gamma _c`$ the detachment energy is $`310^4`$ a.u.). This prevents us from localizing more exactly the field strength at which the crossover takes place. The $`{}_{}{}^{3}(1)`$ state, being the ground state of the anion for $`\gamma >\gamma _c`$ never becomes the most tightly bound state. At $`\gamma =\mathrm{4\hspace{0.17em}000}`$ its electron detachment energy is $`1.25`$ a.u. and therefore much less than the detachment energy of the $`{}_{}{}^{1}0`$ state. This is due to the fact that the tightly bound states are formed by occupying the hydrogenic series $`1s,2p_1,3d_2,\mathrm{}`$ (as mentioned above) and the $`{}_{}{}^{1}0`$ states allows for the $`1s^2`$ configuration yielding the strongest binding although it represents an excited state for $`\gamma >\gamma _c`$ due to its spin character.
The singlet state $`{}_{}{}^{1}(1)`$ is not bound for weak fields. It becomes bound in the regime $`\gamma 15`$ which is an unexpected behavior. The $`{}_{}{}^{1}(1)`$ state lies higher in the spectrum than the bound $`{}_{}{}^{1}(2)`$ and $`{}_{}{}^{1}(3)`$ states for the intermediate field region. In the high field region it however crosses both states. The crossing with the $`{}_{}{}^{1}(3)`$ takes place at $`\gamma 300`$, the crossing with the $`{}_{}{}^{1}(2)`$ state is at $`\gamma \mathrm{4\hspace{0.17em}000}`$. Unfortunately the accuracy of our method is not sufficient to provide a closer look at this crossing. The fact that the $`{}_{}{}^{1}(1)`$ state is not bound for weak fields but bound for strong fields is a consequence of the complicated interplay of the different interactions. The Coulomb repulsion of the two electrons is much weaker for the spatially antisymmetric triplet states compared to the singlet states. The electron-electron repulsion is higher for the states with $`M=1`$ compared to the states with $`M<1`$. This pushes the $`|M|=1`$ singlet states for weak fields beyond the threshold energy, i.e. makes them unbound. The total ionization and the detachment energies of the singlet and triplet states with $`M=1`$ are presented in table II. The suppression of the binding for the singlet state can clearly be seen from this table: the detachment energy of the singlet is $`100`$ times lower than for the triplet at $`\gamma =10`$, but at $`\gamma =4000`$ the ratio is of the order $`2`$. The comparison with the literature (see table II) shows that our detachment energies are variationally lower by several percent than the best available data. For the situation of weakly bound states the improvement is significantly larger.
Let us now consider the energies for the states with $`M=2`$ which are presented in table III. Focusing on the detachment energies we realize that for weak fields the triplet state possesses a larger detachment energy than the singlet state, but for intermediate and high fields the singlet state is stronger bound than the triplet one, i.e. we encounter a crossover which is presented in figure 3. Compared to the data of ref., our method yields $`510`$% higher variational detachment energies for the triplet state and several times higher detachment energies for the singlet one. If we consider the singlet-triplet splitting which is the difference of the total energies between the singlet and the triplet state, where the spin-Zeeman shift is omitted, it can be observed, that for all states this splitting behaves monotonically increasing with increasing field strength in the weak field regime. The splitting for the states with $`M=2`$ and $`M=3`$ are shown in figure 4. The splitting for the $`M=2`$ states increases in weak fields, but for high fields this splitting decreases and becomes negative above some critial field strength. It seems that the Coulomb repulsion, due to antisymmetrization of the wave function is dominated by correlation effects. That the above observation is in fact a consequence of correlation is supported by Vincke and Baye : the reversed order concerning the detachment energies (see figure 3) occurs if they include so-called transverse mixing, which simulates correlations in their approach.
For states with $`M=3`$ only a few published data are available. These states are only weakly bound, although they are stronger bound for $`\gamma 300`$ than the $`{}_{}{}^{1}(1)`$ state. The singlet state has for $`\gamma =0.2`$ a detachment energy of $`7.110^5`$ a.u. and at $`\gamma =1000`$ its detachment energy is $`0.19`$. The electron detachment energies of the triplet state are of the same order of magnitude and the absolute value of the singlet triplet splitting is the lowest of the states considered here. As a consequence a careful convergence study of the results (detachment energy) is indispensable. Our data are given in table IV.
The wavelengths of the transitions of the singlet states are presented in figure 5. The wavelengths are monotonically decreasing with incresing field strength except for the transition from the $`{}_{}{}^{1}(1)`$ state to the $`{}_{}{}^{1}(2)`$ state. As mentioned above these states cross at $`\gamma \mathrm{4\hspace{0.17em}000}`$. Therefore the corresponding wavelength for this transition diverges at the crossing field strength. The transition wavelengths for the triplet states shown in figure 6 are also monotonically decreasing with increasing field strength.
Finally we comment on corrections due to the finite nuclear mass. There are two kinds of corrections, which are relevant here. One, which is special for ions in strong magnetic fields and which describes the coupling between the center of mass motion and the electronic motion. This coupling is due to a motional electric field of intrinsic dynamical origin seen by the moving ion in a magnetic field . Second there are corrections due to the replacement of the naked masses by reduced ones which can be easily included in our data by performing the corresponding shifts . A full dynamical treatment of the atomic ion including the collective motion goes clearly beyond the scope of the present investigation. It is important to note that for the case of the fixed negative donors there naturally occur no such corrections.
## V Brief Summary
We have investigated the H<sup>-</sup> ion, negative donors D<sup>-</sup> respectively, in a strong magnetic field via a fully correlated approach. The key ingredient is an anisotropic Gaussian basis set, whose one-particle wave functions are nonlinearly optimized in order to obtain the spectrum of the one-particle Hamiltonian. In contrast to other basis sets, which are appropriate either for the low field or for the high field regime, our basis set is flexible enough to be adapted to the situation of arbitrary field strength and especially suited for the intermediate field regime. All calculations were performed in the infinite mass frame neglecting relativistic corrections.
We have investigated the low field ground state $`{}_{}{}^{1}0`$, as well as singlet and triplet states for $`M=1,2,3`$ for the broad field regime $`\gamma =810^4410^3`$. For all states and almost all field strengths we could reach at least $`12\%`$ higher binding energies, compared to all other published data. For some states our binding energies were larger by a factor up to two. The global ground state undergoes a crossover with respect to its symmetry which is well-known in the literature : for weak fields $`\gamma 510^2`$ the global ground state is the $`{}_{}{}^{1}0`$ state, whereas for $`\gamma 510^2`$ it is the $`{}_{}{}^{3}(1)`$ state, which is much weaker bound than the $`{}_{}{}^{1}0`$ state for all field strengths. The $`{}_{}{}^{1}(1)`$ state becomes bound for $`\gamma 5`$ and it crosses the $`{}_{}{}^{1}(3)`$ state at $`\gamma 300`$ and the $`{}_{}{}^{1}(2)`$ state at $`\gamma 4000`$. We have also investigated the electronic states with $`M=2`$ in detail. For $`\gamma 1`$ the triplet state is stronger bound than the singlet, whereas for $`\gamma 1`$ the singlet is stronger bound than the triplet. Explanations for the binding mechanisms of the considered states have been provided. The transition wavelengths for all allowed transitions as a function of the field strength are thereby obtained. No stationary transitions which could be of relevance to the astrophysical observation in magnetized white dwarfs have been observed.
### Acknowledgments
The Deutsche Forschungsgemeinschaft (OAA) is gratefully acknowledged for financial support. We thank W. Becken for many fruitful discussions and for his help concerning computational aspects of the present work.
Tables
Figure Captions |
warning/0003/gr-qc0003033.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Inhomogeneous spacetime models generally inherit from the full relativistic dynamics the strong nonlinearity in straightest ways, so that a moderately simple inhomogeneous model can be a good testing ground to obtain an insight toward understandings of generic properties of the relativistic dynamics. One of the best-known such models is the Gowdy model, where, first of all, two commuting spatial Killing vectors are assumed. By this assumption the spatial manifold is simplified to a one dimensional reduced manifold. Second, it is assumed that the spatial manifold is compact without boundary, i.e., closed. This assumption is favorable in view of the fact that no ambiguities occur in the boundary conditions once a topology is specified. A closed space is physically natural, as well, due to, e.g., the finiteness of the gravitational action.
A striking property resulting from these two assumptions is the fact that the possible spatial topologies are restricted. When Gowdy models first appeared in 1971 as global generalizations of cylindrical plane wave solutions, only two spatial topologies, $`S^2\times S^1`$ and $`S^3`$, were considered. After a few years from the first paper, Gowdy became aware of the fact that any closed three dimensional Riemannian manifold which admits two commuting Killing vectors is homeomorphic to one of $`T^3`$, $`S^2\times S^1`$, $`S^3`$ or a lens space $`L(p,q)`$. So, he added $`T^3`$ model to his consideration. (Since lens spaces are finitely covered by the sphere $`S^3`$, he argued little about lens space models.) This class of solutions or cosmological models consists of the original version of the Gowdy models.
Recently, it has been pointed out that a natural generalization is possible, where the two commuting Killing vectors are local, i.e., the Killing vectors are defined in a neighborhood of every point but are not necessarily globally defined. (More precise definition of local Killing vectors will be presented in the next section.) This generalization gives us a set of new favorable models to test the dynamical properties of relativity, especially, as we will see, in connection with the spatial topology.
Part of the motivation of this work comes from a plan to systematically investigate how the spatial topology influence the dynamics of a spacetime. In fact, the relativistic dynamics of a spacetime often seems to be influenced by its spatial topology. For example, we know the recollapsing conjecture , where it is claimed that the well-known recollapsing property of a positive curvature (topologically $`S^3`$) homogeneous and isotropic cosmological model continues to hold also when the spacetime is inhomogeneous (if only appropriate energy and pressure conditions are fulfilled). Thus we may interpret the recollapsing property is a direct consequence derived from the spatial topology.
As another example, note (e.g., ) that the dynamics of a vacuum homogeneous universe can be well characterized by the “potential” (given by the spatial scalar curvature), which can be determined if we know the Bianchi type. The local dynamics of a compactified (locally) homogeneous model can also be understood in the same manner (see ). Moreover in most cases, the corresponding Bianchi type is unique if the spatial manifold is closed. Thus we can say that the spatial topology determines the dynamical behavior of the universe at least at the locally homogeneous limit.
How about when the spacetime becomes inhomogeneous? Can one, as in the locally homogeneous cases, characterize the dynamics of an inhomogeneous universe by the spatial topology? When the spacetimes become inhomogeneous distinctions in local properties will obviously disappear, as the Einstein field equations can be written in the same form, irrespective of the spatial topology. However, the boundary conditions may be different in general. So, we can naturally expect that the distinctions in the dynamical properties due to the differences in the spatial topologies manifest themselves in their global, discrete, statistical, or average properties. (Note that the recollapsing property mentioned above itself is also regarded as a global property.) It seems fascinating to find such properties. Also, such an investigation might provide us with some clues to the topology of our universe.
In this paper we show, after some preparations availing ourselves of the generalized Gowdy models, an example of such properties using those models. More precise contents and contributions should be the following.
We first reveal all the possible spatial topologies of the generalized Gowdy models and classify them. In particular, we categorize all the models into two kinds, then subdivide each of them into finite types each of which has a correspondence to a particular type of locally homogeneous manifolds. Although we show the best classification we have, which will require further considerations (in fact, the number of the possible topologies is infinite), the set of the (finite) types mentioned above will be found to present the most relevant set of classes in an appropriate sense. (To be precise, however, some of the classes can have a discrete parameter.) While by our systematic analysis we find no new topologies other than those in the literature , one of our contributions is that we first asserts that those models do exhaust all the possible ones.
Then we will concentrate on the models belonging to one kind, called the first kind, where each spatial manifold is a $`T^2`$-bundle over $`S^1`$. For these models we argue how we can represent the metric, and discuss the question of natural reduction (to $`S^1`$ as the spatial part). As a result of the reduction, we find that some models give rise to the same reduced Einstein equations with the same boundary conditions for their metric functions. We say that such two models are dynamically equivalent to each other. The dynamical equivalence greatly diminishes the number of representative models (or topologies) we should consider. The basic idea for the reductions is the same as the one presented in Ref., while the dynamical equivalence is first introduced in this paper. All of these settings are an indispensable step for the subsequent study of the models.
Finally, motivated by the reason explained above, we consider varieties of spatial translation and reflection symmetries, and then ask if these symmetries imposed on initial data sets are preserved in time or not for each reduced model. In fact, we find remarkable differences in the reflections.
These are done in Secs.2 to 4. Section 5 is devoted to a summary and comments. In particular, comments on AVTD behavior and the local $`\mathrm{U}(1)`$ models are made. Appendix A gives a summary of the standard description of the Gowdy models with generalizations. Appendix B gives a summary of a calculation for the compactification of Sol geometry.
We adopt the abstract index notation , that is, we use small Latin indices $`a,b,\mathrm{}`$ not to denote components but to represent the type of a tensor explicitly. To denote components we use Greek indices $`\mu ,\nu \mathrm{}`$ or capital Latin indices $`A,B,\mathrm{}`$. Conventionally we write tilde to denote a metric on a universal cover like $`\stackrel{~}{\mathrm{g}}_{ab}`$, while a metric on a quotient space is simply written as $`\mathrm{g}_{ab}`$. In the component representation, however, these metrics are possibly represented in the same way like $`\mathrm{g}_{\mu \nu }`$. We assume that all the spatial three-manifolds are orientable in this paper.
## 2 Topologies and Geometries of Locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric spaces and spacetimes
Let $`(M,h_{ab})`$ be a Riemannian manifold . Suppose for any point $`p`$ in $`M`$ there exit neighborhood $`U`$ which admits Killing vector fields, but these Killing vectors are not necessarily defined on whole $`M`$. For example, consider the manifold $`𝐑^3`$ with metric
$$\mathrm{d}s^2=e^{2\alpha (x)}\mathrm{d}x^2+e^{2\beta (x)}(\mathrm{d}y^2+\mathrm{d}z^2),$$
(1)
where $`\alpha (x)=\alpha (x+2\pi )`$ and $`\beta (x)=\beta (x+2\pi )`$ are real periodic functions. This Riemannian manifold, denoted as $`\stackrel{~}{M}_1`$ hereafter, possesses three independent Killing vectors
$$\xi _2=\frac{}{y},\xi _3=\frac{}{z},\xi _4=z\frac{}{y}+y\frac{}{z}.$$
(2)
Introducing identifications in $`\stackrel{~}{M}_1`$ by infinite actions generated by
$`g_1e^{\pi \xi _4}e^{2\pi \xi _1}`$ $`:`$ $`(x,y,z)(x+2\pi ,y,z),`$
$`g_2e^{2\pi \xi _2}`$ $`:`$ $`(x,y,z)(x,y+2\pi ,z),`$
$`g_3e^{2\pi \xi _3}`$ $`:`$ $`(x,y,z)(x,y,z+2\pi ),`$ (3)
where $`\xi _1\frac{}{x}`$ is a vector field defined for convenience, we obtain a closed manifold $`M_1`$ homeomorphic to $`T^3/𝐙_2`$. (Here, an exponential of a vector represents the diffeomorphism generated by the vector.) See Fig.1. Note that $`\xi _2`$ (and $`\xi _3`$) on the bottom points the opposite direction to that on the top, which fact tells us that $`\xi _2`$ and $`\xi _3`$ cannot be defined on the whole $`M_1`$. $`\xi _4`$ is also incompatible with the identifications by the translations by $`g_2`$ and $`g_3`$. Thus, $`M_1`$ admits no Killing vectors, though we can define Killing vectors on a neighborhood of every point.
If we consider a spacetime manifold $`(M_1\times 𝐑,\mathrm{g}_{ab})`$ whose spatial metric coincides with Eq.(1), this gives a simplest (but rather uninteresting) example of a generalized Gowdy model.
Definition Suppose a Riemannian manifold $`(M,h_{ab})`$ possesses the following two properties:
There exists an open cover $`\left\{O_i\right\}`$ of $`M`$ such that every $`(O_i,h_{ab}^{(i)})`$ admits $`m`$ independent Killing vectors $`\xi _1^{(i)}\mathrm{}\xi _m^{(i)}`$, where $`h_{ab}^{(i)}`$ is the natural metric on $`O_i`$ inherited from $`h_{ab}`$.
On each $`O_iO_j(\mathrm{})`$, $`\xi ^{(i)}`$ and $`\xi ^{(j)}`$ are related through a linear transformation
$$\xi _\alpha ^{(i)}=\underset{\beta =1}{\overset{m}{}}f_\alpha ^{(ij)\beta }\xi _\beta ^{(j)},\alpha =1n,$$
(4)
where $`f_\alpha ^{(ij)\beta }`$ is a nondegenerate constant matrix (for fixed $`i`$ and $`j`$).
By local Killing vectors on $`(M,h_{ab})`$, we mean the collection of the pairs $`(\{O_i\},\{\xi _1^{(i)}\mathrm{}\xi _m^{(i)}\})`$ of the cover and the set of Killing vectors on each patch. For simplicity, we also simply denote them as $`\xi _1\mathrm{}\xi _m`$ if we do not have to say on which patch these Killing vectors are defined.
From the property (ii), the Killing orbits are well defined globally, even though the Killing vectors are local.
Definition A three-dimensional closed Riemannian manifold $`(M,h_{ab})`$ is called a locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric space (or an $`𝒰^2`$-symmetric space in short) if $`(M,h_{ab})`$ admits two commuting local Killing vectors $`_i(\{O_i\},\{\xi _1^{(i)},\xi _2^{(i)}\})`$, i.e., $`[\xi _1^{(i)},\xi _2^{(i)}]=0`$ on every $`O_i`$. We also impose the condition that this manifold must have no other local Killing vectors than those mentioned. We call this condition the genericity condition.
As we will see later, the Killing orbits for this kind of space are flat 2-tori if the local Killing vectors are nondegenerate. We can choose each patch $`O_i`$ in the above definition as a regular neighborhood of such a $`T^2`$ (i.e., a neighborhood such that the boundaries coincides with another orbits), since the two commuting Killing vectors can be defined on such a neighborhood. In this case, the group $`\mathrm{U}(1)\times \mathrm{U}(1)`$ acts on this patch as isometries, hence the word “locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric” for the whole manifold $`(M,h_{ab})`$. However, this terminology is rather lengthy to use frequently, so we also use “$`𝒰^2`$-symmetric” or more simply “$`𝒰^2`$-” in this paper.
Definition A spacetime manifold $`(M\times 𝐑,\mathrm{g}_{ab})`$ is called a locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric spacetime (or a generalized Gowdy spacetime) if the spatial manifold $`M`$ is closed and the spacetime admits two commuting spatial local Killing vectors. The genericity condition is understood, as in the case of the $`𝒰^2`$-symmetric spaces.
Our first task is to determine all the possible topologies for the locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric spaces. It is convenient to consider separately the case where the local Killing vectors degenerate on some points and the case where they do not on any points. We call the latter (nondegenerate) type of manifolds the first kind, and the former (degenerate) ones the second kind.
### 2.1 The first kind
###### Lemma 1
If $`(M,h_{ab})`$ is an orientable locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric space of the first kind, then all Killing orbits must be closed and homeomorphic to $`T^2`$.
Proof: Let $`(O,\eta _{ab})`$ be a Killing orbit, where $`\eta _{ab}`$ is the induced metric on the orbit from $`h_{ab}`$. Since the Killing orbit possesses two commuting Killing vectors, $`\eta _{ab}`$ is a flat metric. Since $`M`$ is orientable, $`O`$ is homeomorphic to one of $`𝐑^2`$, $`𝐑\times S^1`$, or $`T^2`$. First, let us assume $`O𝐑^2`$. Consider a geodesic $`l`$ in $`(O,\eta _{ab})`$, and consider a point sequence $`\{q_i\}`$ such that all $`q_i`$ are on $`l`$ and $`q_i`$ and $`q_{i+1}`$ are unit distant with respect to $`\eta _{ab}`$. In $`M`$ the sequence $`\{q_i\}`$ must have a converging subsequence $`\{p_i\}`$, since $`M`$ is compact. Let $`V_ϵ`$ be the neighborhood of the limit point $`p_{\mathrm{}}`$ with radius $`ϵ`$, i.e., $`|pp_{\mathrm{}}|<ϵ`$ for all $`pV_ϵ`$, where $`||`$ is the geodesic distance with respect to $`h_{ab}`$. $`V_ϵ`$ contains all $`p_i`$ $`(i>N_ϵ)`$ for an integer $`N_ϵ`$, so that for any $`p_i`$ and $`p_j`$, $`(i,j>N_ϵ)`$, $`|p_ip_j|<2ϵ`$ as a result of the triangle inequality. On the other hand, the geodesic $`l`$ is a “straight line” in the Euclid plane $`(𝐑^2,\eta _{ab})`$, so that the distant between any two points $`p_i`$ and $`p_j`$ $`(ij)`$ in $`O`$ is larger than or equal to unity. Hence we can choose a sequence of “isometric disconnected neighborhoods” $`\{U_i\}`$ in $`M`$ such that $`p_iU_i`$, $`S_iS_j=\mathrm{}`$, where $`S_iU_iO`$, and there exist an isometry $`\varphi _{ij}:(p_i,S_i,U_i)(p_j,S_j,U_j)`$ for all $`i,j`$. Since $`S_iS_j=\mathrm{}`$ and $`|p_ip_j|<2ϵ`$, we must have the consequence $`p_jS_i`$ but $`p_jU_i`$ for a sufficiently small $`ϵ`$ (i.e., for sufficiently large $`i`$ and $`j`$). Moreover, this implies that for an arbitrarily small $`ϵ^{}<ϵ`$ there exist two points in $`M`$, $`p_i`$ and $`p_j`$, such that $`|p_ip_j|<2ϵ^{}`$, $`p_jS_i`$ but $`p_jU_i`$, and there exist the isometry $`\varphi _{ij}:p_ip_j`$. This in turn implies that there exist a third motion in $`(M,h_{ab})`$ off the orbit, which fact is against the genericity condition. (That is, $`𝐑^2`$-orbits can be formed only if $`(M,h_{ab})`$ is locally homogeneous.) In the case of $`O𝐑\times S^1`$, we can choose a geodesic $`l`$ which is open as in the case of $`O𝐑^2`$. Repeating a similar argument we conclude that this case is also against the genericity. Thus, the only possibility is the case $`OT^2`$.
###### Theorem 2
If $`(M,h_{ab})`$ is an orientable locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric space of the first kind, then $`M`$ is homeomorphic to a $`T^2`$-bundle over $`S^1`$, for which the local Killing vectors generate the $`T^2`$-fibers.
Proof: Consider a normal vector field that is everywhere nonvanishing and normal to the Killing orbit, and consider the integration curve $`c`$ of this field passing through an arbitrarily given point $`pM`$. Let $`O`$ be the Killing orbit to which $`p`$ belongs, so $`c`$ intersects with $`O`$ at $`p`$. If $`c`$ did not intersect again with $`O`$, then it would imply that the total volume of $`M`$ is infinite, which fact contradicts with the compactness of $`M`$. Hence $`c`$ intersects with $`O`$ again. This implies that if we regard each Killing orbit as a point, then $`M`$ reduces to $`S^1`$. Since all orbits are homeomorphic to $`T^2`$ as in Lemma 1, $`M`$ is a $`T^2`$-bundle over $`S^1`$.
Topologically, any $`T^2`$-bundle over $`S^1`$ can be obtained by first considering the product $`T^2\times I`$, where $`I=[0,1]`$ is the unit interval, then identifying $`T^2\times \{0\}`$ and $`T^2\times \{1\}`$ by the action of a gluing map $`\varphi :T^2T^2`$. Since any two gluing maps which are homotopic to each other results in topologically the same manifold, we usually think of the gluing map $`\varphi `$ as an element of the mapping class group of $`T^2`$, $`\mathrm{mcg}(T^2)`$. Here, the mapping class group of a manifold $`M`$ is the group of diffeomorphisms of $`M`$ modulo diffeomorphisms which are homotopic to the identity, $`\mathrm{mcg}(M)=\mathrm{Diff}(M)/\mathrm{Diff}_0(M)`$. The group $`\mathrm{mcg}(T^2)`$ is well known as the modular group, $`\mathrm{GL}(2,𝐙)`$. However, since we are interested only in orientable manifolds, we consider the orientation-preserving mapping class group of the torus, $`\mathrm{mcg}_+(T^2)\mathrm{SL}(2,𝐙)`$. Letting
$$A=\left(\begin{array}{cc}p& q\\ r& s\end{array}\right)\mathrm{SL}(2,𝐙),$$
(5)
the fundamental group of the corresponding space is represented by
$$\pi _1=g_1,g_2,g_3;[g_2,g_3]=1,g_1g_2g_1^1=g_2^pg_3^r,g_1g_3g_1^1=g_2^qg_3^s,$$
(6)
where $`g_2`$ and $`g_3`$ are generators of the fiber $`T^2`$ and $`g_1`$ is that of the base $`S^1`$. We therefore have a natural map $`\omega _1:\mathrm{SL}(2,𝐙)W_1`$, where $`W_1`$ is the set of orientable $`T^2`$-bundles over $`S^1`$. While $`\omega _1`$ is not injective, we have the following useful theorem, of which proof can be found in Ref. (Theorem 5.5 therein);
###### Theorem 3 ()
Let $`M`$ be the total space of a $`T^2`$-bundle over $`S^1`$ with gluing map $`\varphi `$, and let $`A\mathrm{GL}(2,𝐙)`$ represent the automorphism of $`\pi _1(T^2)`$ induced by $`\varphi `$. Then, $`M`$ admits a Sol-structure if $`A`$ is hyperbolic, an $`E^3`$-structure if $`A`$ is periodic, or a Nil-structure otherwise. In particular, if $`|\mathrm{Tr}A|>2`$, then $`A`$ is hyperbolic, so $`M`$ admits a Sol-structure.
This theorem tells us that every $`T^2`$-bundle over $`S^1`$ admits a “geometric structure” , in other words, admits a locally homogeneous metric. This structure is, as in the above theorem, one of the three types, $`E^3`$, Nil, or Sol. Conversely, we can find all the $`T^2`$-bundles over $`S^1`$ from the closed quotients of $`E^3`$, Nil, and Sol, so we can refer to known classifications of the closed quotients of the three geometries to classify $`T^2`$-bundles over $`S^1`$. This procedure simultaneously determines the corresponding geometric structure for each representative. (To avoid confusions we remark that locally homogeneous metrics are used just for convenience here to discuss topologies, but we will see that appropriate inhomogeneous metrics are obtained by “relaxing” such a metric.)
The case of $`E^3`$: All the orientable closed manifolds modeled on $`E^3`$ are classified into six manifolds , $`T^3`$, $`T^3/𝐙_2`$, $`T^3/𝐙_3`$, $`T^3/𝐙_4`$, $`T^3/𝐙_6`$, and $`T^3/𝐙_2\times 𝐙_2`$. All of these except the last one can be regarded as $`T^2`$-bundle over $`S^1`$. Associated with these five are the representative elements $`E_\lambda \mathrm{SL}(2,𝐙)`$ $`(\lambda =1,2,3,4,6,\mathrm{respectively})`$:
$$E_1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),E_2=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),E_3=\left(\begin{array}{cc}0& 1\\ 1& 1\end{array}\right),E_4=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),E_6=\left(\begin{array}{cc}0& 1\\ 1& 1\end{array}\right).$$
(7)
The case of Nil: All the orientable closed manifolds modeled on Nil can also be completely classified (See page 4878 of Ref.). Among them, ones which can be regarded as $`T^2`$-bundle over $`S^1`$ are given by the following two families $`N_{\pm 1}(n)\mathrm{SL}(2,𝐙)`$:
$$N_1(n)=\left(\begin{array}{cc}1& n\\ 0& 1\end{array}\right),N_1(n)=\left(\begin{array}{cc}1& n\\ 0& 1\end{array}\right),$$
(8)
where $`n`$ is a nonzero integer. The parameter $`n`$ in $`N_1(n)`$ can be chosen to be positive, since the relabeling $`(g_1,g_3)(g_3,g_1)`$ reverses the sign of $`n`$. For $`N_1(n)`$, the opposite sign of $`n`$ corresponds to a distinct topology.
The case of Sol: As in Theorem 3, if $`|\mathrm{Tr}A|>2`$, we have a $`T^2`$-bundle over $`S^1`$ modeled on Sol. All those for $`|\mathrm{Tr}A|2`$ are non-orientable , so we will not consider them. Any two ones with distinct values of $`\mathrm{Tr}A`$ such that $`|\mathrm{Tr}A|>2`$ are not homeomorphic, so one parameter family of $`\mathrm{SL}(2,𝐙)`$ ,
$$S(n)=\left(\begin{array}{cc}0& 1\\ 1& n\end{array}\right),$$
(9)
where $`n`$ is an integer such that $`|n|>2`$, gives a one parameter family of distinct $`T^2`$-bundles over $`S^1`$ modeled on Sol. However, this does not implies that the ones modeled on Sol can be classified only by $`\mathrm{Tr}A`$ like Eq.(9), since there exist topologically distinct manifolds with the same $`\mathrm{Tr}A`$ . No complete classification of the closed manifolds modeled on Sol seems to be known. Nevertheless, as we will see in the next section, the sequence (9) gives nice representatives as reduced relativistic models.
### 2.2 The second kind
###### Theorem 4
If $`(M,h_{ab})`$ is a locally $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric space of the second kind, then $`M`$ is homeomorphic to one of $`S^3`$, $`S^2\times S^1`$, or a lens space $`L(p,q)`$.
Proof: Since $`M`$ is closed, if removing an open neighborhood of the degenerate points from $`M`$, the resulting manifold $`\widehat{M}`$ with boundaries is compact. Moreover, the boundaries can be chosen so as to coincide with (regular) Killing orbits. Let $`\widehat{M}`$ be the boundaries of $`\widehat{M}`$ so chosen. Note that in the proof of Lemma 1, only the compactness of $`M`$ is assumed to show that the Killing orbits are $`T^2`$. Hence all the Killing orbits for $`(\widehat{M},\widehat{h}_{ab})`$ are also $`T^2`$. (Here, $`\widehat{h}_{ab}`$ is the restriction of $`h_{ab}`$ to $`\widehat{M}`$.) In particular, every connected component of $`\widehat{M}`$ is $`T^2`$. Consider a (nonvanishing) normal vector field with respect to the Killing orbits in $`(\widehat{M},\widehat{h}_{ab})`$. This flow of the normal vector field defines a unique one-to-one correspondence between boundary components, since the image of a boundary component by this flow must end up with another boundary component, due to the compactness of $`\widehat{M}`$. Since $`M`$ and therefore $`\widehat{M}`$ are assumed to be connected, $`\widehat{M}`$ must have only two boundary components and is naturally homeomorphic to the product $`T^2\times I`$, where $`I[0,1]`$ is the unit interval. Next, consider the neighborhood of a degenerate Killing orbit removed. An action of $`\mathrm{U}(1)\times \mathrm{U}(1)`$ can degenerate only to $`\mathrm{U}(1)`$, so a connected set of degenerate points forms a circle, and a neighborhood of such a circle with boundary forms a solid torus. The boundary, which is homeomorphic to $`T^2`$, can again be chosen so as to coincide with a (regular) Killing orbit. Consider two such neighborhoods, $`V_1`$ and $`V_2`$. Then, the original manifold is recovered by attaching them to $`\widehat{M}`$ along the boundaries: $`MV_1\widehat{M}V_2`$. Note, however, that $`V_1\widehat{M}`$ is another solid torus with the $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetry. Let us rewrite this manifold as $`V_1`$. The original manifold $`M`$ is now obtained simply by identifying the boundaries of the two solid tori $`V_1`$ and $`V_2`$, $`MV_1V_2`$. Finally, it is a well-known fact (e.g. ) that the sum of two solid tori is homeomorphic to one of $`S^3`$, $`S^2\times S^1`$, or a lens space $`L(p,q)`$.
Cosmological models based on $`𝒰^2`$-symmetric spatial manifolds of the second kind can be thought of as a global generalization of cylindrically symmetric (plane wave) models. In fact, the symmetry axis of a cylindrically symmetric space is a degenerate orbit for the $`\mathrm{U}(1)\times \mathrm{U}(1)`$ action of the cylindrical symmetry. If, as usual, taking this axis as $`z`$-coordinate axis, and introducing identifications $`zz+z_0`$ for a constant $`z_0`$, we have the symmetry axis that is a circle. A regular neighborhood of this circle is a solid torus with the $`\mathrm{U}(1)\times \mathrm{U}(1)`$ action. An $`𝒰^2`$-manifold $`M`$ can be obtained by identifying (the boundaries of) such two $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric solid tori.
Now, it should be remarked that the local Killing vectors for an $`𝒰^2`$-manifold of the second kind are always globally defined. This is because this kind of space can also be regarded as the product $`T^2\times I`$ (with a metric that degenerates on the boundaries). The original Gowdy models in fact include all the models based on $`S^3`$, $`S^2\times S^1`$, and the lens spaces $`L(p,q)`$. We therefore do not discuss much about these models in subsequent sections.
We end this subsection with a summary of the classification (e.g. ) of the lens spaces, for completeness. Consider a gluing map $`\varphi :T^2T^2`$ to identify the boundaries of two $`\mathrm{U}(1)\times \mathrm{U}(1)`$ symmetric solid tori, $`V_1`$ and $`V_2`$. Since if two gluing maps are homotopic the resulting manifolds are homeomorphic, it is again natural to think of an identification map as an element of the mapping class group, $`\varphi \mathrm{mcg}(T^2)\mathrm{GL}(2,𝐙)`$. Let $`\left(\begin{array}{c}1\\ 0\end{array}\right)`$ and $`\left(\begin{array}{c}0\\ 1\end{array}\right)`$ be, respectively, meridian and longitudinal loops of $`V_1`$ or $`V_2`$. We define the action of $`A=\left(\begin{array}{cc}p& r\\ q& s\end{array}\right)\mathrm{GL}(2,𝐙)`$ by the left action to these vectors. Then, the closed manifold $`M`$ made by gluing the two solid tori with respect to $`A`$ is the lens space $`L(p,q)`$. (The topology of $`M`$ is regardless to $`r`$ and $`s`$, i.e., it is determined only by the mapping of the meridian loop of $`V_1`$.) The integers $`p`$ and $`q`$ are coprime and we can set $`p0`$ (since $`L(p,q)L(p,q)`$). $`S^3`$ and $`S^2\times S^1`$ correspond, respectively, to $`L(1,q)`$ and $`L(0,1)`$. Lens spaces are completely classified by the fact that $`L(p,q)`$ and $`L(p^{},q^{})`$ are homeomorphic if and only if $`p=p^{}`$, and $`q\pm q^{}\mathrm{mod}p`$ or $`qq^{}\pm 1\mathrm{mod}p`$.
## 3 Relaxation, Reduction and Dynamical Equivalence
In this section we present an detailed account of the reduction procedure for the $`𝒰^2`$-symmetric models of the first kind, with emphasis on the significance of the corresponding geometric structures. We will also obtain useful results, the “dynamical equivalences”.
First, we remark that since the spacetime manifold is the product $`M\times 𝐑`$, where the spatial manifold $`M`$ is a ($`T^2`$-)bundle over $`S^1`$ with fibers generated by the local Killing vectors, the natural reduced manifold obtained by contracting the Killing orbits should be $`S^1\times 𝐑`$. However, as we see bellow, when the bundle is not (covered by) a direct product, this reduction is not trivial.
We represent an $`𝒰^2`$-symmetric spacetime (of the first kind) $`(M\times 𝐑,\mathrm{g}_{ab})`$ with the covering map
$$\pi :(\stackrel{~}{M}\times 𝐑,\stackrel{~}{\mathrm{g}}_{ab})(M\times 𝐑,\mathrm{g}_{ab})=(\stackrel{~}{M}\times 𝐑,\stackrel{~}{\mathrm{g}}_{ab})/\mathrm{\Gamma },$$
(10)
where $`(\stackrel{~}{M}\times 𝐑,\stackrel{~}{\mathrm{g}}_{ab})`$ is the universal covering manifold of $`(M\times 𝐑,\mathrm{g}_{ab})`$, and $`\mathrm{\Gamma }`$ is a covering group, which acts on $`\stackrel{~}{M}`$ discretely. Since the spatial manifold $`M`$ is a $`T^2`$-bundle over $`S^1`$, $`\stackrel{~}{M}`$ is homeomorphic to $`𝐑^3`$. We can therefore use globally defined coordinates, say $`(x^0,x^1,x^2,x^3)(t,x,y,z)`$, to represent $`\stackrel{~}{\mathrm{g}}_{ab}`$. A natural form of the metric admitting two commuting Killing vectors represented with these coordinates is $`\mathrm{g}_{\mu \nu }(t,x)dx^\mu dx^\nu `$, where Greek indices $`\mu ,\nu ,\mathrm{}`$ run over $`0`$ to $`3`$. Moreover, fixing the freedom of diffeomorphisms in $`\mathrm{g}_{\mu \nu }(t,x)dx^\mu dx^\nu `$ as much as possible, we can represent the metric $`\stackrel{~}{\mathrm{g}}_{ab}`$ as
$$\mathrm{d}s^2=e^{\gamma /2}(\mathrm{d}t^2+\mathrm{d}x^2)+2N_A\mathrm{d}x^A\mathrm{d}t+Re_{AB}\mathrm{d}x^A\mathrm{d}x^B,$$
(11)
where $`\gamma `$, $`R`$, $`N_A`$ and $`e_{AB}`$ (Capital indices $`A,B,\mathrm{}`$ run 2 to 3) are functions of $`t`$ and $`x`$. We set $`dete_{AB}=1`$ so that $`R`$ describes the area of the Killing orbit. This metric is called the generic metric of the canonical form. We will also consider the restricted metric with $`N_A=0`$ (the “two-surface orthogonality” )
$$\mathrm{d}s^2=e^{\gamma /2}(\mathrm{d}t^2+\mathrm{d}x^2)+Re_{AB}\mathrm{d}x^A\mathrm{d}x^B,$$
(12)
which is called the two-surface orthogonal metric of the canonical form. For a standard prescription of Einstein’s equations for the last metric, see Appendix A. (Note that the metric functions are represented with “bars” in this Appendix. We will take this notation from the next subsection to avoid conflicts with “noncanonical” metrics which will appear.)
Both metrics admit two dimensional isometry group consisting of translations along $`y`$ and $`z`$ axes. Note that the $`y`$-$`z`$ planes descend to the $`T^2`$-fibers after appropriate identifications in each $`y`$-$`z`$ plane. These identifications are generated by two independent vectors, which we can choose without loss of generality as the unit coordinate generators, $`/y`$ and $`/z`$. Then we can naturally think of the actions of these generators as the representation of the generators, $`g_2`$ and $`g_3`$, of $`\pi _1`$ into the isometry group; $`g_2e^\frac{}{y}`$, $`g_3e^\frac{}{z}`$. The remaining generator $`g_1`$ must induce a translation along $`x`$-axis plus the modular transformation induced by an element $`A=n^A{}_{B}{}^{}\mathrm{SL}(2,𝐙)`$. Thus, $`g_1:(x,y,z)(x+2\pi ,n^A{}_{B}{}^{}x_{}^{B})`$ on the $`t=`$constant space. Together with $`g_3:(x,y,z)(x,y+2\pi ,z)`$ and $`g_3:(x,y,z)(x,y,z+2\pi )`$, all the relations in the fundamental group (6) are fulfilled. The boundary conditions for the metric (11) or (12) is determined by the requirement that the action of $`g_1`$ be an isometry of the metric, which can be easily found;
$`\lambda (t,x)=\lambda (t,x+2\pi ),R(t,x)=R(t,x+2\pi ),`$
$`N_A(t,x)=n^B{}_{A}{}^{}N_{B}^{}(t,x+2\pi ),e_{AB}(t,x)=n^C{}_{A}{}^{}n_{}^{D}{}_{B}{}^{}e_{CD}^{}(t,x+2\pi ).`$ (13)
Note that the components $`e_{AB}`$ that describe each $`T^2`$-fiber are not periodic functions in general, so that the reduced manifold (spanned by $`t`$ and $`x`$) cannot be naturally regarded as $`S^1\times 𝐑`$, if we represent the metric as Eq.(11) or (12). The reason why $`e_{AB}`$ do not automatically become periodic is that what two independent components in $`e_{AB}`$ themselves describe are Teichmüller parameters for the $`T^2`$-fiber, rather than moduli parameters. (For the difference between Teichmüller and moduli parameters, see, e.g., Ref..)
An idea to remedy the situation is to choose the metric functions so that they become constant when the spatial metric is at a locally homogeneous limit. That is, conversely, we “relax” a locally homogeneous metric in an appropriate way to obtain a suitable metric. Note that a constant function is naturally a function on $`S^1`$. As far as considering smooth deformations of the metric, metric functions chosen in such a way must continue to be (smooth) functions on $`S^1`$ even when the metric becomes inhomogeneous (with the $`𝒰^2`$ symmetry).
The scheme One nice way to realize this idea is to expand the spatial metric in terms of the (left) invariant one-forms $`\sigma ^i`$ $`(i=13)`$ of a Bianchi group $`G`$. As we have seen in Sec.2.1, any $`T^2`$-bundle over $`S^1`$ admits one of $`E^3`$, Nil, or Sol-structure. These geometric structures correspond, respectively, to Bianchi I(VII<sub>0</sub>), II, and VI<sub>0</sub>. ($`E^3`$ has multiple correspondences.) The appropriate Bianchi group $`G`$ is determined from this correspondence. Let $`\xi _i`$ $`(i=13)`$ be independent generators of $`G`$, that is, $`\xi _i`$ are Killing vectors of a $`G`$-invariant metric. Since $`\sigma ^i`$ are by definition invariant under the action of $`G`$, the homogeneous metric is written as $`h_{ij}\sigma ^i\sigma ^j`$ with (spatially) constant components $`h_{ij}`$. The Killing vectors of this homogeneous metric include a commuting pair. Suppose $`[\xi _2,\xi _3]=0`$. Let $`\chi `$ be a vector field such that $`\chi `$ is independent of $`\xi _2`$ and $`\xi _3`$ at every point and is invariant under the action generated by $`\xi _2`$ and $`\xi _3`$, i.e., $`_{\xi _2}\chi =_{\xi _3}\chi =0`$. Here, $``$ represents the Lie derivative. Since $`[\xi _2,\xi _3]=0`$, we can choose coordinates ($`y`$ and $`z`$) generated by $`\xi _2`$ and $`\xi _3`$. Moreover, since $`_{\xi _2}\chi =_{\xi _3}\chi =0`$ implies the commutativity $`[\xi _2,\chi ]=[\xi _3,\chi ]=0`$, we can choose a third coordinate ($`x`$) as that of generated by $`\chi `$. Consider an arbitrary function $`f(x)`$ that depends only on $`x`$ chosen in this way. Then, the level set of $`f(x)`$ is invariant under the diffeomorphisms generated by $`\xi _2`$ and $`\xi _3`$, and it coincides with the set of the orbits generated by $`\xi _2`$ and $`\xi _3`$. Hence we can write the inhomogeneous metric as $`h_{ij}(x)\sigma ^i\sigma ^j`$, which is invariant under the diffeomorphisms generated by $`\xi _2`$ and $`\xi _3`$, and inhomogeneous in the desired manner. (Recall that all $`\sigma ^i`$ are invariant under $`\xi _2`$ and $`\xi _3`$.) The metric functions $`h_{ij}(x)`$ should be periodic in $`x`$ because of the reason explained above. Hence, a spacetime metric with $`h_{ij}(t,x)\sigma ^i\sigma ^j`$ as the spatial part naturally gives rise to a reduction onto $`S^1\times 𝐑`$. We call this scheme the relaxation method. If $`\xi _1=\chi `$ as in the Bianchi I case, then this scheme descends to the usual coordinate representation like Eq.(11) or (12), but otherwise it is nontrivial.
Recall that a Bianchi group $`G`$ is a three-dimensional simply transitive group acting on a three-dimensional simply connected manifold $`\stackrel{~}{M}`$, and therefore we can identify $`G`$ with $`\stackrel{~}{M}`$ (e.g. ). We adhere to this viewpoint in this paper and use the components of $`G`$ to represent the coordinates of $`\stackrel{~}{M}`$, as well.
One comment should be made here. As pointed out in Ref., all Bianchi homogeneous spaces except for VIII and IX possess a commuting pair of Killing vectors. Even for VIII and IX, if considering higher symmetry there can exist commuting Killing vectors. Moreover, these homogeneous spaces can be compactified to closed spaces, except for IV and VI<sub>a</sub> types. Therefore all the locally homogeneous manifolds of the Bianchi types except IV and VI<sub>a</sub> seems to be locally homogeneous limits of $`𝒰^2`$-manifolds. However, as we mentioned above, only Bianchi I, VII<sub>0</sub>, II, and VI<sub>0</sub> types can be such a limit for the first kind manifolds. The second kind manifolds correspond to Bianchi IX (in case of $`MS^3`$ and $`L(p,q)`$) and the Nariai-Kantowski-Sachs type (in case of $`MS^2\times S^1`$). Thus, Bianchi III, V, VII<sub>a</sub>, VIII types are missing in the list of possible manifolds as the locally homogeneous limits of our $`𝒰^2`$-manifolds. The reason is because each orbit of the two commuting local Killing vectors for a locally homogeneous manifold of these types always do not close. For example, in the case of Bianchi V, it forms $`𝐑^2`$. As we have seen in the proof of Lemma 1, such a case is possible only when the manifold is locally homogeneous. That is, we can say that, as a corollary to Theorems 2 and 4:
###### Corollary 5
Every closed locally homogeneous manifold of Bianchi III, V, VII<sub>a</sub> and VIII types cannot be relaxed to be inhomogeneous with $`𝒰^2`$ symmetry.
In this sense, these manifolds are “stiff.” It is worth noting that a common feature of them is that they all contain a hyperbolic structure, $`H^2`$ or $`H^3`$.
In the following subsections, we perform the reduction procedure for each type of Nil, Sol, and $`E^3`$. The spacetime manifold $`M\times 𝐑`$ reduces to $`S^1\times 𝐑`$ for any spatial manifold $`M`$. One important feature that is found as a result of the reductions is that, even for models with topologically distinct spatial manifolds, distinctions at the reduced level can completely degenerate. In other words, we obtain for some models the same set of PDEs (as the reduced Einstein equations) with the same boundary conditions. For such models the universal covering spacetime metrics can be represented in exactly the same form with only the difference contained in the covering group $`\mathrm{\Gamma }`$. We call that such models are dynamically equivalent to each other.
### 3.1 Case of Nil
The Bianchi II group $`G_{\mathrm{II}}`$ is the three-dimensional simply transitive group with the multiplication rule
$$\left(\begin{array}{c}a\\ b\\ c\end{array}\right)\left(\begin{array}{c}x\\ y\\ z\end{array}\right)=\left(\begin{array}{c}a+x\\ b+y\\ c+z+ay\end{array}\right).$$
(14)
This group is generated by
$$\xi _1=\frac{}{x}+y\frac{}{z},\xi _2=\frac{}{y},\xi _3=\frac{}{z}.$$
(15)
The one-forms which are invariant under the action of $`G_{\mathrm{II}}`$ are
$$\sigma ^1=\mathrm{d}x,\sigma ^2=\mathrm{d}y,\sigma ^3=\mathrm{d}zx\mathrm{d}y.$$
(16)
Since $`\xi _2`$ and $`\xi _3`$ commute and they generate coordinates $`y`$ and $`z`$, a desired representation of the spatial metric is given by $`h_{ij}(x)\sigma ^i\sigma ^j`$.
Before proceeding further, however, we have to clarify a subtle point about the fact that we can choose another commuting pairs of group generators. For example, $`\xi _1`$ and $`\xi _3`$ also commute to each other, and the orbits they generate are distinct from those generated by $`\xi _2`$ and $`\xi _3`$. In general, since the commutation rule for linear combinations of the group generators is given by
$$[\alpha ^i\xi _i,\beta ^j\xi _j]=\alpha ^i\beta ^j[\xi _i,\xi _j]=(\alpha ^1\beta ^2\alpha ^2\beta ^1)\xi _3,$$
(17)
any two generators $`\alpha =\alpha ^i\xi _i`$ and $`\beta =\beta ^j\xi _j`$ commute if and only if $`\alpha ^1\beta ^2\alpha ^2\beta ^1=0`$. The last condition implies that $`\xi _3`$ must be tangent to the surface spanned by $`\alpha `$ and $`\beta `$ if $`\alpha `$ and $`\beta `$ are independent. Such a surface, conversely, is spanned by $`\eta _2\mathrm{sin}\theta \xi _1+\mathrm{cos}\theta \xi _2`$ and $`\eta _3\xi _3`$, where $`\theta `$ is a real parameter. Hence, we have one-parameter ($`\theta `$) family of distinct sets of orbits generated by a commuting pair. This then implies that we have one-parameter degree of freedom of relaxing a locally homogeneous manifold of Bianchi II type.
However, this fact is actually insignificant if viewing the metric as a universal cover metric. That is, the freedom of $`\theta `$ can be absorbed in the freedom of choosing the covering group $`\mathrm{\Gamma }`$. Therefore we do not have to consider the freedom of $`\theta `$, which can be fixed $`\theta =0`$ without loss of generality. (The model based on $`\theta =\pi /2`$ was presented in Ref. as ‘Type 2’, which is redundant for this reason.)
Now, we have established the fact that the relaxed (universal covering) metric can be represented by $`h_{ij}(x)\sigma ^i\sigma ^j`$ with the basis (16). The form of spacetime metric corresponding to the two-surface orthogonal class of metrics (113) is given by
$$\mathrm{d}s^2=e^{\gamma /2}(\mathrm{d}t^2+(\sigma ^1)^2)+R[e^P(\sigma ^3+Q\sigma ^2)^2+e^P(\sigma ^2)^2],$$
(18)
where the metric functions $`\gamma `$, $`R`$, $`P`$ and $`Q`$ are functions of $`t`$ and $`x`$, and periodic with respect to $`x`$. (We have changed the choice of $`P`$ and $`Q`$ from that in Ref., since the present choice gives us the most natural and simplest correspondence to $`\overline{P}`$ and $`\overline{Q}`$. See bellow.) These metric functions are related to those for the canonical representation (113) through
$$\overline{\gamma }=\gamma ,\overline{R}=R,\overline{P}=P,\overline{Q}=Qx.$$
(19)
The reduced Einstein equations for the unbarred variables $`\gamma `$, $`P`$, $`Q`$, and $`R`$ are therefore easily obtained by substituting these equations into those for the canonical (i.e. barred) variables presented in Appendix A.
Our next task is to check that we can indeed compactify the universal cover. We choose the period for the metric functions along the $`x`$ axis as $`2\pi `$:
$$f(t,x)=f(t,x+2\pi ),\text{for }f=\gamma ,R,P\text{ and }Q.$$
(20)
We denote the isometry group for the spacetime metric (18) as $`H_{\mathrm{II}}^2`$, and its subgroup which is a subgroup of $`G_{\mathrm{II}}`$ as $`H_{\mathrm{II}}G_{\mathrm{II}}`$. We have
$$H_{\mathrm{II}}=\left\{\left(\begin{array}{c}2\pi n\\ b\\ c\end{array}\right)G_{\mathrm{II}}\right|n𝐙,b,c𝐑\}.$$
(21)
The full isometry group $`H_{\mathrm{II}}^2`$ is generated by $`H_{\mathrm{II}}`$ and the $`𝐙_2`$-isometry
$$k:(x,y,z)(x,y,z).$$
(22)
Recall that the fundamental groups for $`N_1(n)`$ and $`N_1(n)`$ are given by Eqs.(8) with (6). Putting
$$g_i=\left(\begin{array}{c}2\pi n_i\\ g_i^2\\ g_i^3\end{array}\right),i=13,$$
(23)
we have to solve the relations in these fundamental groups for the parameters $`g_i^j`$ and $`n_i`$. The solutions are easily found for $`N_1(n)`$. For $`N_1(n)`$ we have to think of $`g_1`$ as a composite of $`k`$ and an element of $`H_{\mathrm{II}}`$. The solutions are given by
For $`N_1(n)`$:
$`\mathrm{\Gamma }_n`$ $`=`$ $`\{g_1,g_2,g_3\}`$ (33)
$`=`$ $`\{\left(\begin{array}{c}2\pi p\\ g_1^2\\ g_1^3\end{array}\right),\left(\begin{array}{c}0\\ 0\\ \frac{2\pi }{n}(pg_3{}_{}{}^{2}qg_1{}_{}{}^{2})\end{array}\right),\left(\begin{array}{c}2\pi q\\ g_3^2\\ g_3^3\end{array}\right)\},`$
where $`p,q𝐙`$, $`g_i{}_{}{}^{j}𝐑`$, and $`pg_3{}_{}{}^{2}qg_1{}_{}{}^{2}0`$;
For $`N_1(n)`$:
$`\mathrm{\Gamma }_n`$ $`=`$ $`\{g_1,g_2,g_3\}`$ (43)
$`=`$ $`\{k\left(\begin{array}{c}2\pi p\\ g_1^2\\ g_1^3\end{array}\right),\left(\begin{array}{c}0\\ 0\\ \frac{2\pi pg_3^2}{n}\end{array}\right),\left(\begin{array}{c}0\\ g_3^2\\ g_3^3\end{array}\right)\},`$
where $`p𝐙`$, $`g_i{}_{}{}^{j}𝐑`$, and $`pg_3{}_{}{}^{2}0`$. The point here is that there exist solutions, ensuring that we can compactify the universal cover for any closed spatial manifold $`N_{\pm 1}(n)`$. For completeness we remark that the parameters appearing in $`\mathrm{\Gamma }_n`$ can be fixed arbitrarily, since redefinitions of metric functions can make the parameters equal to arbitrary values. For example, we can take $`p=g_3{}_{}{}^{2}=1`$, and $`q=g_1{}_{}{}^{2}=g_1{}_{}{}^{3}=g_3{}_{}{}^{3}=0`$ for $`N_1(n)`$, and $`p=g_3{}_{}{}^{2}=1`$, and $`g_1{}_{}{}^{2}=g_1{}_{}{}^{3}=g_3{}_{}{}^{3}=0`$ for $`N_1(n)`$.
Note that we did not have to impose any other condition than Eq.(20) for all spatial topologies $`N_{\pm 1}(n)`$. This results in obtaining the same reduced Einstein equations with the same boundary conditions (20). Thus:
###### Proposition 6
All (two-surface orthogonal) $`𝒰^2`$-symmetric models of Nil type are dynamically equivalent.
Here, the two-surface orthogonality is necessary for $`N_1(n)`$ models to ensure that the map $`k`$ is an isometry. That is, generic $`N_1(n)`$ models with shift functions are not allowed. On the other hand, all generic $`N_1(n)`$ models are dynamically equivalent.
### 3.2 Case of Sol
The invariant one-forms are given by
$$\sigma ^1=\mathrm{d}x,\sigma ^2=\frac{1}{\sqrt{2}}(e^{\mathrm{q}x}\mathrm{d}y+e^{\mathrm{q}x}\mathrm{d}z),\sigma ^3=\frac{1}{\sqrt{2}}(e^{\mathrm{q}x}\mathrm{d}y+e^{\mathrm{q}x}\mathrm{d}z),$$
(44)
where $`\mathrm{q}>0`$ is a positive parameter introduced for convenience.
Using these 1-forms we can write the two-surface orthogonal metric as
$$\mathrm{d}s^2=e^{\gamma /2}(\mathrm{d}t^2+(\sigma ^1)^2)+R[e^P(\sigma ^2{}_{}{}^{}+Q\sigma ^3{}_{}{}^{})^2+e^P(\sigma ^3{}_{}{}^{})^2],$$
(45)
where $`\gamma `$, $`P`$, $`Q`$, and $`R`$, are functions of $`t`$ and $`x`$, and they are assumed to be periodic in $`x`$. We have defined
$$\sigma ^2{}_{}{}^{}\frac{1}{\sqrt{2}}(\sigma ^2\sigma ^3)=e^{\mathrm{q}x}\mathrm{d}y,\sigma ^3{}_{}{}^{}\frac{1}{\sqrt{2}}(\sigma ^2+\sigma ^3)=e^{\mathrm{q}x}\mathrm{d}z.$$
(46)
(The choice of metric functions $`P`$ and $`Q`$ here is different from that in Ref.. We have done so for apparent simplicity, but we should keep in mind that at the homogeneous limit the metric is not diagonalizable with respect to $`(\sigma ^1,\sigma ^2{}_{}{}^{},\sigma ^3{}_{}{}^{})`$, in contrast to with respect to $`(\sigma ^1,\sigma ^2,\sigma ^3)`$, for vacuum spacetimes.) The metric functions defined with Eq.(45) are related to those in the canonical representation (113) through
$$\overline{\gamma }=\gamma ,\overline{R}=R,\overline{P}=P+2\mathrm{q}x,\overline{Q}=e^{2\mathrm{q}x}Q.$$
(47)
The reduced Einstein equations for $`\gamma `$, $`P`$, $`Q`$, and $`R`$ can be obtained by substituting these equations into those presented in Appendix A.
For compactifications, see Appendix B. In particular, the isometry group of the metric (45) is given by the $`H_{\mathrm{VI}_0}^2`$ presented in the Appendix. Therefore we can compactify the universal cover possessing this metric for all $`A`$. If we choose the period along $`x`$-axis for the metric functions as $`2\pi `$, we must put $`c_3=2\pi `$. Accordingly, we must choose the parameter $`\mathrm{q}`$ so that $`e^{2\pi \mathrm{q}}`$ equals to the greater absolute eigenvalue of the matrix $`A`$, since we have assumed $`\mathrm{q}>0`$. Note that the characteristic polynomial (135) depends only on $`\mathrm{Tr}A`$. Let $`n\mathrm{Tr}A`$. Then we obtain
$$\mathrm{q}=\frac{1}{2\pi }\mathrm{log}\frac{\left|n\right|+\sqrt{n^24}}{2}.$$
(48)
The reduced Einstein equations for the (unbarred) metric functions depends (only) on the parameter $`\mathrm{q}`$, and the boundary conditions are the common simple periodic boundary condition. Moreover, $`\mathrm{q}`$ depends only on $`\left|n\right|`$. We thus conclude:
###### Proposition 7
An $`𝒰^2`$-symmetric model of Sol type is specified with an element $`A\mathrm{mcg}_+(T^2)\mathrm{SL}(2,𝐙)`$ such that $`\left|\mathrm{Tr}A\right|>2`$. Let $`A_1`$ and $`A_2`$ be such matrices. Then, the corresponding models are dynamically equivalent if $`\left|\mathrm{Tr}A_1\right|=\left|\mathrm{Tr}A_2\right|`$.
From this fact, we find it suffices to consider the models with the one-parameter family $`S(n)`$ ($`n>2`$), defined in Eq.(9), as representatives.
### 3.3 Case of $`E^3`$
The corresponding Bianchi types are Bianchi I and VII<sub>0</sub>. More precisely, geometry $`E^3`$ is the Euclid space $`𝐑^3`$ with the standard metric $`\mathrm{d}x^2+\mathrm{d}y^2+\mathrm{d}z^2`$. The isometry group of $`E^3`$ is therefore formed by the translations $`𝐑^3`$ and rotations $`\mathrm{O}(3)`$, and is isomorphic to the Poincaré group, $`\mathrm{Isom}E^3\mathrm{IO}(3)`$. $`\mathrm{IO}(3)`$ contains two simply transitive subgroups, the Bianchi I ($`G_\mathrm{I}𝐑^3`$) and VII<sub>0</sub> ($`G_{\mathrm{VII}_0}`$) groups, so that the correspondence above is arrived at.
We can make the inhomogeneous metrics from both $`G_\mathrm{I}`$ and $`G_{\mathrm{VII}_0}`$. However, the effects are not equivalent. First, consider the metric made from $`G_\mathrm{I}`$. Since the invariant one-forms coincide with the usual coordinate basis $`\mathrm{d}x`$, $`\mathrm{d}y`$, and $`\mathrm{d}z`$, the spacetime metric and the boundary conditions are the same as the canonical ones (12) and (3). We find in particular that the metric functions for $`E_1`$ and $`E_2`$ are simply periodic. For $`E_3`$, $`E_4`$, and $`E_6`$, however, the metric functions must obey the last condition in Eqs.(3), which does not imply the simple periodic conditions, (though they have period $`6\pi `$, $`4\pi `$, and $`6\pi `$ for $`E_3`$, $`E_4`$, and $`E_6`$, respectively). As long as we adhere to the metric obtained from $`G_\mathrm{I}`$, this feature is inevitable, since the closed manifolds corresponding to $`E_3`$, $`E_4`$, and $`E_6`$ cannot be realized in the Bianchi I group $`G_\mathrm{I}`$, or $`G_\mathrm{I}^2`$. Here, $`G_\mathrm{I}^2`$ is generated by $`G_\mathrm{I}`$ and the $`𝐙_2`$-map $`h:(x,y,z)(x,y,z)`$, which is an isometry of the metric (12). Precisely, the fundamental group $`\pi _1(E_1)`$ corresponding to $`E_1`$ can be embedded in $`G_\mathrm{I}`$, and $`\pi _1(E_2)`$ can be embedded in the $`G_\mathrm{I}^2`$. However, other fundamental groups can be embedded neither in $`G_\mathrm{I}`$ nor $`G_\mathrm{I}^2`$. Since the isometry group of the metric (12) becomes a subgroup of $`G_\mathrm{I}^2`$, provided that the metric functions are all periodic with period $`2\pi `$, this implies that the embeddings of $`\pi _1`$ for $`E_3`$, $`E_4`$, and $`E_6`$ into the subgroup are impossible. Thus, the appropriate boundary conditions cannot be the simple periodic conditions for them.
On the other hand, all the fundamental groups $`\pi _1(E_\lambda )`$ $`(\lambda =1,2,3,4,6)`$ can be embedded into $`G_{\mathrm{VII}_0}`$ , which fact implies that the boundary conditions can become simple periodic ones for the metric obtained from $`G_{\mathrm{VII}_0}`$ for all $`E_\lambda `$. We present an explicit prescription bellow.
The invariant one-forms for Bianchi VII<sub>0</sub> are given by
$$\sigma ^1=\mathrm{d}x,\sigma ^2=\mathrm{cos}\mathrm{q}x\mathrm{d}y+\mathrm{sin}\mathrm{q}x\mathrm{d}z,\sigma ^3=\mathrm{sin}\mathrm{q}x\mathrm{d}y+\mathrm{cos}\mathrm{q}x\mathrm{d}z,$$
(49)
where $`\mathrm{q}>0`$ is a positive parameter. Using these 1-forms we can write the two-surface orthogonal metric as
$$\mathrm{d}s^2=e^{\gamma /2}(\mathrm{d}t^2+(\sigma ^1)^2)+R[e^P(\sigma ^2+Q\sigma ^3)^2+e^P(\sigma ^3)^2],$$
(50)
where $`\gamma `$, $`P`$, $`Q`$, and $`R`$, are functions of $`t`$ and $`x`$, and they are assumed to be periodic in $`x`$ with period $`2\pi `$. The parameter $`\mathrm{q}`$ will be chosen so that we can compactify the universal cover. (The role of $`\mathrm{q}`$ is similar to the $`\mathrm{q}`$ appearing in the Sol model.) The correspondence to the barred metric functions in the canonical metric (113) is given by
$`e^{\overline{P}(x)}`$ $`=`$ $`\mathrm{\Phi }(P(x),Q(x),\mathrm{cos}\mathrm{q}x,\mathrm{sin}\mathrm{q}x),\overline{Q}(x)=\mathrm{\Psi }(P(x),Q(x),\mathrm{cos}\mathrm{q}x,\mathrm{sin}\mathrm{q}x),`$
$`e^{P(x)}`$ $`=`$ $`\mathrm{\Phi }(\overline{P}(x),\overline{Q}(x),\mathrm{cos}\mathrm{q}x,\mathrm{sin}\mathrm{q}x),Q(x)=\mathrm{\Psi }(\overline{P}(x),\overline{Q}(x),\mathrm{cos}\mathrm{q}x,\mathrm{sin}\mathrm{q}x).`$ (51)
where the two functions $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are defined by, for arbitrary $`P`$, $`Q`$, $`c`$, and $`s`$,
$`\mathrm{\Phi }(P,Q,c,s)`$ $`=`$ $`e^P(csQ)^2+e^Ps^2,`$
$`\mathrm{\Psi }(P,Q,c,s)`$ $`=`$ $`{\displaystyle \frac{e^P(csQ)(s+cQ)e^Pcs}{\mathrm{\Phi }(P,Q,c,s)}}.`$ (52)
Note that the vacuum Einstein equations in terms of the unbarred variables contain the parameter $`\mathrm{q}`$.
As usual, we represent the elements of the Bianchi group in the column vector form. The multiplication rule for $`G_{\mathrm{VII}_0}`$ is given by
$$\left(\begin{array}{c}a\\ b\\ c\end{array}\right)_\mathrm{q}\left(\begin{array}{c}x\\ y\\ z\end{array}\right)_\mathrm{q}=\left(\begin{array}{c}a+x\\ \left(\begin{array}{c}b\\ c\end{array}\right)+R_{\mathrm{q}a}\left(\begin{array}{c}y\\ z\end{array}\right)\end{array}\right)_\mathrm{q},\left(\begin{array}{c}a\\ b\\ c\end{array}\right)_\mathrm{q}^1=\left(\begin{array}{c}a\\ R_{\mathrm{q}a}\left(\begin{array}{c}b\\ c\end{array}\right)\end{array}\right)_\mathrm{q},$$
(53)
where $`R_{\mathrm{q}a}`$ is the rotation matrix by angle $`\mathrm{q}a`$;
$$R_{\mathrm{q}a}=\left(\begin{array}{cc}\mathrm{cos}\mathrm{q}a& \mathrm{sin}\mathrm{q}a\\ \mathrm{sin}\mathrm{q}a& \mathrm{cos}\mathrm{q}a\end{array}\right).$$
(54)
The subscripts $`\mathrm{q}`$ appearing in the column vectors are to remind that the multiplication rule is defined with respect to $`\mathrm{q}`$. The one-forms (49) are invariant under this action.
We denote the isometry group of the metric (50) which is a subgroup of $`G_{\mathrm{VII}_0}`$ as $`H_{\mathrm{VII}_0}`$. It is given by
$$H_{\mathrm{VII}_0}=\left\{\left(\begin{array}{c}2\pi n\\ b\\ c\end{array}\right)_\mathrm{q}G_{\mathrm{VII}_0}\right|n𝐙,b,c𝐑\}.$$
(55)
The (full) isometry group $`H_{\mathrm{VII}_0}^2`$ of the metric (50) is generated by $`H_{\mathrm{VII}_0}`$ and the $`𝐙_2`$-isometry $`h:(x,y,z)(x,y,z)`$.
We can perform the embedding of every $`\pi _1(E_\lambda )`$ both with and without the $`𝐙_2`$-isometry $`h`$, as we can read from Sec.V A of Ref.. (But, be careful with the differences in the choice of representations of the relations in $`\pi _1`$.) When we do not use $`h`$, i.e., embed the fundamental group $`\pi _1(E_\lambda )`$ into $`H_{\mathrm{VII}_0}`$, the appropriate value of $`\mathrm{q}`$ is simply given by
$$\mathrm{q}=\frac{1}{\lambda }.$$
(56)
Hence all the models seem to be dynamically inequivalent. However, we can choose $`\mathrm{q}=1`$ for $`\lambda =2`$, and $`\mathrm{q}=1/3`$ for $`\lambda =6`$, if we use $`h`$. For example, in the case of $`E_6`$ one can confirm that the solution of the embedding into $`H_{\mathrm{VII}_0}`$ is given by
$`\mathrm{\Gamma }_6`$ $`=`$ $`\{g_1,g_2,g_3\}`$ (67)
$`=`$ $`\{\left(\begin{array}{c}2\pi \\ g_1^2\\ g_1^3\end{array}\right)_{\frac{1}{6}},\left(\begin{array}{c}0\\ g_2^2\\ g_2^3\end{array}\right)_{\frac{1}{6}},\left(\begin{array}{c}0\\ R_{\pm \pi /3}\left(\begin{array}{c}g_2^2\\ g_2^3\end{array}\right)\end{array}\right)_{\frac{1}{6}}\}`$
with $`\mathrm{q}=1/6`$. Here, $`g_i^j`$ are real parameters. Also, we can embed the fundamental group into $`H_{\mathrm{VII}_0}^2`$ as
$$\mathrm{\Gamma }_6^{}=\{h\left(\begin{array}{c}2\pi \\ g_1^2\\ g_1^3\end{array}\right)_{\frac{1}{3}},\left(\begin{array}{c}0\\ g_2^2\\ g_2^3\end{array}\right)_{\frac{1}{3}},\left(\begin{array}{c}0\\ R_{2\pi /3}\left(\begin{array}{c}g_2^2\\ g_2^3\end{array}\right)\end{array}\right)_{\frac{1}{3}}\}$$
(68)
with $`\mathrm{q}=1/3`$. Since the fundamental group of $`E_3`$ can be embedded into $`H_{\mathrm{VII}_0}`$ with $`\mathrm{q}=3`$ as
$$\mathrm{\Gamma }_3=\{\left(\begin{array}{c}2\pi \\ g_1^2\\ g_1^3\end{array}\right)_{\frac{1}{3}},\left(\begin{array}{c}0\\ g_2^2\\ g_2^3\end{array}\right)_{\frac{1}{3}},\left(\begin{array}{c}0\\ R_{\pm 2\pi /3}\left(\begin{array}{c}g_2^2\\ g_2^3\end{array}\right)\end{array}\right)_{\frac{1}{3}}\},$$
(69)
the $`E_6`$ and $`E_3`$ models are dynamically equivalent. (The significance of the parameters $`g_i^j`$ is the same as in the case of Nil. In particular, we can take $`g_1{}_{}{}^{2}=g_1{}_{}{}^{3}=0`$, $`g_2{}_{}{}^{2}=1`$, and $`g_2{}_{}{}^{3}=0`$, in all cases above.)
As a result, we have only three dynamically inequivalent classes;
###### Proposition 8
The $`E_1`$, $`E_3`$, and $`E_4`$ models comprise a set of representatives for the dynamical equivalence in the two-surface orthogonal $`𝒰^2`$-symmetric models of $`E^3`$ type.
For the generic models of $`E^3`$ type with the shift functions, all the $`E_\lambda `$ models are dynamically inequivalent, since the map $`h`$ is not an isometry.
To end this subsection, we remark again that the $`𝒰^2`$-symmetric models of $`E^3`$ type best match with Bianchi VII<sub>0</sub>, not Bianchi I. However, if limited to the $`T^3`$ (or $`T^3/𝐙_2`$) model, which corresponds to $`E_1`$ ($`E_2`$), what it naturally corresponds is Bianchi I. To make later discussions easier, it may therefore be convenient to put forward the following convention.
Convention When speaking of the “$`E^3`$ models”, we understand them to have correspondence to Bianchi VII<sub>0</sub>. In particular, the “unbarred variables” for an $`E^3`$ model are those in the metric (50). However, when speaking of the “$`T^3`$ model”, we understand it to have correspondence to Bianchi I. Therefore the unbarred variables coincide with barred variables appearing in the canonical metric (113).
## 4 Translation and Reflection symmetries
The Gowdy equations (120) have several symmetries. First of all, from the fact that these equations do not have explicit $`x`$-dependence, they have the natural translation symmetry with respect to $`xxa`$, where $`a`$ is a real parameter. This means that if $`(\overline{P}(t,x),\overline{Q}(t,x))`$ is a solution for the equations, then $`(\overline{P}(t,x+a),\overline{Q}(t,x+a))`$ is also a solution. One more symmetry that naturally comes to our attention is the reflection symmetry with respect to $`xx`$, which is a result of the invariance of the equations under the transformation $`/x/x`$. This symmetry means that if $`(\overline{P}(t,x),\overline{Q}(t,x))`$ is a solution for the equations, then $`(\overline{P}(t,x),\overline{Q}(t,x))`$ is also a solution. Another feature that results from these symmetries is that (in the Hamiltonian picture) an initial data that is symmetric with respect to $`xx+a`$ or $`xx`$ maintains the same symmetry under the time evolution. This point will be discussed in the second subsection bellow.
For the usual Gowdy $`T^3`$ model, both $`xxa`$ and $`xx`$ transformations acting on $`(\overline{P}(t,x),\overline{Q}(t,x))`$ preserve the boundary conditions on $`\overline{P}`$ and $`\overline{Q}`$, so that they express true symmetries. For Nil and Sol models, however, these transformations do not preserve the boundary conditions, so that they do not express true symmetries. Nevertheless, we can find translation symmetries for Nil and Sol models in a generalized sense, and also find (generalized) reflection symmetries for them. We also find that the $`T^3`$ model admits larger amounts of reflection symmetries than mentioned above. In the first subsection bellow, we obtain these translation and reflection symmetries in a systematic way. Based on this result, in the second subsection we discuss the dynamical point of view of the symmetries, which reveals a manifestation of the influence of spatial topologies to dynamics.
### 4.1 Translation and Reflection Transformations
All symmetries concerning the Gowdy equations (120) should have their origin in the metric (113). For example, this metric is invariant under the diffeomorphism $`xxa`$, up to the redefinition of the metric functions $`(\overline{P}(x),\overline{Q}(x))(\overline{P}(x+a),\overline{Q}(x+a))`$ (with the similar one for $`\overline{\lambda }`$). (For simplicity, we omit writing the argument $`t`$ in the metric functions.) This redefinition accounts for the translation symmetry mentioned above. Recall that, if a metric is a solution for the vacuum Einstein equation, then so is the metric induced from a (spatial) diffeomorphism. Thus, in general, if such a diffeomorphism leaves the metric (113) invariant up to a redefinition of the metric functions, this redefinition defines a symmetry for the Gowdy equations.
Consider a diffeomorphism $`\varphi :(x,y,z)(\varphi _x(x,y,z),\varphi _y(x,y,z),\varphi _z(x,y,z))`$ with $`t`$ fixed. As remarked, $`\varphi `$ must preserve the characteristic form of the metric (113) so that a redefinition of the metric functions can make the metric invariant. In particular, to preserve the isothermal form $`e^{\overline{\lambda }/2}(\mathrm{d}t^2+\mathrm{d}x^2)`$ of the metric we must have $`\varphi _x(x,y,z)=\pm xa`$. Moreover, $`(y,z)(\varphi _y,\varphi _z)`$ must be a linear transformation on each $`y`$-$`z`$ plane. That is, we must have
$$\varphi :(x,𝐲)(\pm xa,L𝐲),$$
(70)
where $`L\mathrm{SL}^2`$ is a constant matrix and $`𝐲`$ is the column form vector of $`(y,z)`$. Here, $`\mathrm{SL}^2\{m\mathrm{GL}(2,𝐑)|det(m)=\pm 1\}`$. (We do not have to consider the translations in the $`y`$-$`z`$ plane of the type $`(x,y,z)(x,y+y_0,z+z_0)`$, which are isometries for the metric (113), since these give rise to no effect to the Gowdy equations.) We decompose the $`\varphi _x`$ part into $`xxa`$ and $`xx`$, and will call symmetries involved with the former (generalized) translation symmetries and call symmetries involved with the latter (generalized) reflection symmetries. We can also consider symmetries generated by diffeomorphisms involved with the identity $`xx`$, but since this type of symmetry is of little interest from the dynamical point of view discussed in the next subsection, we do not consider them here. We can find all possible translation and reflection symmetries defined above by finding $`\varphi _y(y,z)`$ and $`\varphi _z(y,z)`$ (or $`L`$) for each case, as follows.
We first consider the translation symmetries. As we will see in the next subsection, this symmetry implies that (in the Hamiltonian picture) a symmetric initial data with respect to a (generalized) translation maintains its symmetry under the time evolution. At the (locally) homogeneous limit, this implies that the diffeomorphism that gives rise to this symmetry should coincide with an isometry for the metric. Recall that there is a three-dimensional isometry group $`G`$ for each case of $`T^3`$, $`E^3`$, Nil, and Sol at the homogeneous limit. It contains a two-dimensional commutative isometry subgroup $`H𝐑^2`$ that is preserved as isometry group when the metric becomes $`𝒰^2`$-symmetric but inhomogeneous. Then, there is a one-parameter isometry subgroup $`K`$ which acts transversely to $`H`$. In the $`𝒰^2`$-symmetric but inhomogeneous cases, the actions of $`K`$ are served as the desired translations. As easily found from the multiplication rules for the Bianchi I, VII<sub>0</sub>, II, and VI<sub>0</sub> groups, they are, respectively, given by
$`T_a`$ $`:`$ $`(x,y,z)(xa,y,z),`$ (71)
$`E_a`$ $`:`$ $`(x,y,z)(xa,\mathrm{cos}\mathrm{q}ay+\mathrm{sin}\mathrm{q}az,\mathrm{sin}\mathrm{q}ay+\mathrm{cos}\mathrm{q}az),`$ (72)
$`N_a`$ $`:`$ $`(x,y,z)(xa,y,zay),`$ (73)
$`S_a`$ $`:`$ $`(x,y,z)(xa,e^{2qa}y,e^{2qa}z).`$ (74)
Since these leave the invariant one-forms for, respectively, Bianchi I, VII<sub>0</sub>, II, and VI<sub>0</sub> types invariant, the corresponding transformations (redefinitions) takes the simplest form for the unbarred metric functions:
$`T_a`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\overline{P}(x+a),\overline{Q}(x+a))\text{for metric functions in (}\text{113}\text{)},`$ (75)
$`E_a`$ $`:`$ $`(P(x),Q(x))(P(x+a),Q(x+a))\text{for metric functions in (}\text{50}\text{)},`$ (76)
$`N_a`$ $`:`$ $`(P(x),Q(x))(P(x+a),Q(x+a))\text{for metric functions in (}\text{18}\text{)},`$ (77)
$`S_a`$ $`:`$ $`(P(x),Q(x))(P(x+a),Q(x+a))\text{for metric functions in (}\text{45}\text{)}.`$ (78)
Here, we assume $`a=2\pi /n`$ for a positive integer $`n`$. Using the relations (51), (19) and (47), $`E_a`$, $`N_a`$ and $`S_a`$ expressed with the canonical metric (113) are found to be
$`E_a`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\mathrm{log}\mathrm{\Phi }(\overline{P}(x+a),\overline{Q}(x+a),\mathrm{cos}\mathrm{q}a,\mathrm{sin}\mathrm{q}a),`$ (79)
$`\mathrm{\Psi }(\overline{P}(x+a),\overline{Q}(x+a),\mathrm{cos}\mathrm{q}a,\mathrm{sin}\mathrm{q}a)),`$
$`N_a`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\overline{P}(x+a),\overline{Q}(x+a)+a),`$ (80)
$`S_a`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\overline{P}(x+a)2\mathrm{q}a,e^{2\mathrm{q}a}\overline{Q}(x+a)).`$ (81)
The functions $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are defined in Eq.(52). $`T_a`$, $`E_a`$, $`N_a`$ and $`S_a`$ preserve the boundary conditions for the $`𝒰^2`$-symmetric models of, respectively, $`T^3`$, $`E^3`$, Nil and Sol types, so they define the translation symmetry for each type. Note that the translations for each type form the infinite cyclic group $`𝐙`$ for a fixed $`a`$, or form the cyclic group $`𝐙_n`$ for a fixed $`a=2\pi /n`$ if the (periodic) boundary conditions are taken into account.
Next, we consider the reflection symmetries, by which we mean the symmetries generated by the reflection transformations defined as follows.
Definition Consider the map $`R:(x,𝐲)(x,L𝐲)`$ for a given $`L\mathrm{SL}^2`$. The symmetry transformation $`R_{}`$ induced from $`R`$, acting on the variables $`(P,Q)`$ of a given model, is called a reflection transformation, if the following three conditions are fulfilled;
For $`R`$ (and $`R_{}`$) to form the $`𝐙_2`$-group together with the identity $`\mathrm{id}`$,
$$R^2RR=\mathrm{id}.$$
(82)
As explained later,
$$\mathrm{Conj}(R)TRTR^1=T^1,$$
(83)
where $`T`$ is the map inducing the translation transformation of the given model. (Hence, $`T`$ is one of Eqs.(71) $``$ (74).)
$`R_{}`$ preserves the boundary conditions. That is, if variables $`(P,Q)`$ satisfy the boundary conditions appropriate for the model, $`R_{}(P,Q)`$ also satisfy the same boundary conditions.
First, from condition (i), $`L`$ must be one of
$$L_0\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),L_\pi \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\text{ or}L_{l_\theta }\left(\begin{array}{cc}\mathrm{cos}2\theta & \mathrm{sin}2\theta \\ \mathrm{sin}2\theta & \mathrm{cos}2\theta \end{array}\right),$$
(84)
where $`\theta [0,\pi )`$ is a real parameter. These act on the $`y`$-$`z`$ plane as, respectively, the identity, the rotation by angle $`\pi `$, and the reflection with respect to the line $`l_\theta `$ which passes through the origin with angle $`\theta `$. Thus we have obtained the following candidates for the reflection $`R`$:
$`R_0`$ $`:`$ $`(x,y,z)(x,y,z),`$ (85)
$`R_\pi `$ $`:`$ $`(x,y,z)(x,y,z),`$ (86)
$`R_{l_\theta }`$ $`:`$ $`(x,y,z)(x,y\mathrm{cos}2\theta +z\mathrm{sin}2\theta ,y\mathrm{sin}2\theta z\mathrm{cos}2\theta ).`$ (87)
Condition (ii) concerns a “compatibility” with the natural translation obtained earlier. That is, the reflection about a generic point $`x=a`$ should be induced from that about $`x=0`$ by a translation. We can equivalently require that the conjugation, $`\mathrm{Conj}(R)T`$, of the translation $`T`$ by a reflection $`R`$ reverse the original translation, as indicated by Eq.(83). For example, the conjugation of the translation $`T_a`$ for the $`T^3`$ model by the map $`R_0`$ is given by $`\mathrm{Conj}(R_0)T_a=R_0T_a(R_0)^1=R_0T_aR_0=T_a`$. This shows $`R_0`$ is compatible with $`T_a`$. On the other hand, say, the conjugation of $`N_a`$ by $`R_0`$ gives $`\mathrm{Conj}(R_0)N_a=R_0N_aR_0=T_{2a}N_a`$, which is not the reverse of $`N_a`$. All the compatible relations are given as follows:
$`\text{For }T^3\text{ model:}\mathrm{Conj}(R_0)T_a`$ $`=`$ $`T_a,`$ (88)
$`\mathrm{Conj}(R_\pi )T_a`$ $`=`$ $`T_a,`$ (89)
$`\mathrm{Conj}(R_{l_\theta })T_a`$ $`=`$ $`T_a(\theta [0,\pi )),`$ (90)
$`\text{For }E^3\text{ model:}\mathrm{Conj}(R_{l_\theta })E_a`$ $`=`$ $`E_a(\theta [0,\pi )),`$ (91)
$`\text{For Nil model:}\mathrm{Conj}(R_{l_\theta })N_a`$ $`=`$ $`N_a(\theta =0,{\displaystyle \frac{\pi }{2}}),`$ (92)
$`\text{For Sol model:}\mathrm{Conj}(R_{l_\theta })S_a`$ $`=`$ $`S_a(\theta ={\displaystyle \frac{\pi }{4}},{\displaystyle \frac{3}{4}}\pi ).`$ (93)
The transformations for the metric functions induced from these maps are given as follows:
$`R_0,R_\pi `$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\overline{P}(x),\overline{Q}(x)),`$ (94)
$`R_{l_\theta }`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\mathrm{log}\mathrm{\Phi }(\overline{P}(x),\overline{Q}(x),\mathrm{sin}2\theta ,\mathrm{cos}2\theta ),`$ (95)
$`\mathrm{\Psi }(\overline{P}(x),\overline{Q}(x),\mathrm{sin}2\theta ,\mathrm{cos}2\theta )),`$
where the functions $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are defined in Eq.(52). The explicit forms for Nil ($`\theta =0,\pi /2`$) and Sol ($`\theta =\pi /4,3\pi /4`$) models are
$`R_{l_0},R_{l_{\pi /2}}`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\overline{P}(x),\overline{Q}(x)),`$ (96)
$`R_{l_{\pi /4}},R_{l_{3\pi /4}}`$ $`:`$ $`(\overline{P}(x),\overline{Q}(x))(\overline{P}+\mathrm{log}(\overline{Q}^2+e^{2\overline{P}}),{\displaystyle \frac{\overline{Q}}{\overline{Q}^2+e^{2\overline{P}}}})|_{xx}.`$ (97)
Note that $`R_0=R_\pi `$, $`R_{l_0}=R_{l_{\pi /2}}`$, and $`R_{l_{\pi /4}}=R_{l_{3\pi /4}}`$.
As a final task we are left with checking condition (iii). For the $`T^3`$ model, we at once see that if $`\overline{P}(x)`$ and $`\overline{Q}(x)`$ are periodic functions, then their images by $`R_0`$ and $`R_{l_\theta }`$ are also periodic functions. Thus, $`R_0`$ and $`R_{l_\theta }`$ give reflections for the $`T^3`$ model. For the $`E^3`$, Nil and Sol models it is convenient to work with the unbarred variables. Interestingly the transformations (95) (for $`E^3`$), (96) (for Nil) and (97) (for Sol) are written with the unbarred variables in exactly the same form (except for the disappearance of bars);
For $`E^3`$:
$`R_{l_\theta }`$ $`:`$ $`(P(x),Q(x))(\mathrm{log}\mathrm{\Phi }(P(x),Q(x),\mathrm{sin}2\theta ,\mathrm{cos}2\theta ),`$ (98)
$`\mathrm{\Psi }(P(x),Q(x),\mathrm{sin}2\theta ,\mathrm{cos}2\theta )),`$
For Nil:
$$R_{l_0}:(P(x),Q(x))(P(x),Q(x)),$$
(99)
For Sol:
$$R_{l_{\pi /4}}:(P(x),Q(x))(P+\mathrm{log}(Q^2+e^{2P}),\frac{Q}{Q^2+e^{2P}})|_{xx}.$$
(100)
From these it is trivial to see that the periodicity of the unbarred variables are preserved for $`E^3`$, Nil and Sol. This confirms that these transformations do give reflections for the corresponding models. Table 1 bellow summarizes the translation and reflection transformations for each model.
Remark that we have inclusion relations. The $`T^3`$ model has the largest set of reflections. The $`E^3`$ models have part of that for the $`T^3`$ model. The Nil and Sol models have parts of that for the $`E^3`$ models.
### 4.2 Dynamical Interpretation
As well known , the Gowdy equations (120) admit a Hamiltonian formulation. Let $`\pi _P(x)`$ and $`\pi _Q(x)`$ be, respectively, conjugate momenta of $`P(x)`$ and $`Q(x)`$. The phase space $`𝒫`$ is spanned by the four functions $`(P(x),Q(x),\pi _P(x),\pi _Q(x))`$ (or similarly for variables with bars). Consider the set $`𝒪`$ of maps on $`𝒫`$, $`𝒪\{o|o:𝒫𝒫\}`$. The Hamiltonian flow $`\psi _\tau `$, which generates the time evolution by time $`\tau `$, is naturally regarded as an element in $`O`$, $`\psi _\tau 𝒪`$, for a fixed $`\tau `$. $`\psi _0`$ is supposed to be the identity.
What we are interested in here are maps (operators) which commute with $`\psi _\tau `$, so we define
$$𝒪_{\mathrm{com}}\{f𝒪|{}_{}{}^{}\tau 𝐑,f\psi _\tau =\psi _\tau f\}.$$
(101)
We call an $`f𝒪_{\mathrm{com}}`$ a symmetry operator. (If the operator $`f`$ was a smooth flow like $`\psi _\tau `$ we would be able to reformulate this commutativity to the vanishing of a Poisson bracket. However, since we consider discrete operators as $`f`$ this is not the case. As a result, we do not actually need symplectic structures if only the flow $`\psi _\tau `$ is properly defined.) For an initial data $`\gamma 𝒫`$, $`\psi _\tau (\gamma )`$ is a solution for the Einstein equation as a function of $`\tau `$. If there is a symmetry operator $`f`$, then the function of $`\tau `$, $`f(\psi _\tau (\gamma ))`$, is also a solution, since this equals to $`\psi _\tau (f(\gamma ))`$, which is the solution with the initial data $`f(\gamma )`$. Therefore a symmetry operator generates a (in general, distinct) solution from a solution.
Now, note that, while a symmetry transformation obtained in the previous subsection is defined in the configuration space spanned by $`P`$ and $`Q`$, it induces an operator on the phase space $`𝒫`$ by differentiating the configuration variables with respect to $`\tau `$. We use in particular the Hamiltonian equations
$$\dot{P}=\pi _P,\dot{Q}=e^{2P}\pi _Q.$$
(102)
(This relation is the same for all sets of barred or unbarred variables.) Thus, the symmetry transformations are also regarded as an operator on $`𝒫`$. This operator is clearly a symmetry operator. Conversely, if there is a symmetry operator, it gives a symmetry transformation, as seen from the previous paragraph. Thus, in effect, symmetry transformations and operators are equal entities. Note that, since we have exhausted all the symmetry transformations concerning translation and reflection, we now know all the symmetry operators for them.
To represent the symmetry operator obtained from a symmetry transformation we use the same symbol as that of the symmetry transformation. For example, the translation transformation $`T_a`$ for $`T^3`$ model defined in Eq.(75) naturally gives the operator, denoted also as $`T_a`$, that simply shifts the argument by a constant:
$$T_a:(\overline{P}(x),\overline{Q}(x),\overline{\pi }_P(x),\overline{\pi }_Q(x))(\overline{P}(x+a),\overline{Q}(x+a),\overline{\pi }_P(x+a),\overline{\pi }_Q(x+a)).$$
(103)
The translation symmetry operators $`E_a`$ (for $`E^3`$), $`N_a`$ (for Nil) and $`S_a`$ (for Sol) are also defined by similar maps if the phase space variables are taken as the “unbarred” ones. (cf. Eqs.(76), (77), and (78).)
An interesting consequence of the symmetry operators comes from considering data which are invariant with respect to this operator. For $`f𝒪_{\mathrm{com}}`$, let
$$𝒮(f)\{\gamma 𝒫|f(\gamma )=\gamma \}.$$
(104)
For example, if $`f`$ is a translation symmetry operator $`T_{2\pi /n}`$, the set $`𝒮(T_{2\pi /n})`$ is comprised of data in $`𝒫`$ which are translation-symmetric with period $`2\pi /n`$. If $`f`$ is a reflection symmetry operator $`R_{}`$, the set $`𝒮(R_{})`$ is comprised of reflection-symmetric data in $`𝒫`$. The remarkable fact is that these symmetries are preserved dynamically, since Eq.(101), together with $`f(\gamma )=\gamma `$, implies $`f(\psi _\tau (\gamma ))=\psi _\tau (\gamma )`$ for all $`\tau `$. In other words, any data $`\gamma 𝒮(f)`$ remains within $`𝒮(f)`$ under the dynamical flow, i.e., $`\psi _\tau (𝒮(f))𝒮(f)`$. Thus, $`𝒮(f)`$ defines an invariant subset of $`𝒫`$ for the flow $`\psi _\tau `$.
We interpret the invariant subset $`𝒮(f)`$ as a dynamical character of the model. Motivated by this interpretation, we compare $`𝒮(f)`$ to see if there are distinctions in dynamical characters of the models. To this, recall that if we think of the phase space $`𝒫`$ as being spanned by the unbarred variables, $`𝒫`$ for every model is a space of four functions which are all periodic functions. (We think that the barred and unbarred variables coincide for the $`T^3`$ model.) Therefore the phase spaces for the four models are naturally identified. In this view, the flow $`\psi _\tau `$ depends on the type of the model. Strictly speaking, however, since the “constraints” equations corresponding to Eq.(123) take different forms for the $`E^3`$, Nil and Sol models, this identification may be justified only approximately, but we neglect this point.
First, consider the translations. As already mentioned the translation operator for every model is obtained by simply shifting the spatial coordinate like Eq.(103) for the unbarred variables. Hence, all the invariant subsets $`𝒮(T_a)`$, $`𝒮(E_a)`$, $`𝒮(N_a)`$, and $`𝒮(S_a)`$ for $`a=2\pi /n`$, where $`n`$ is a positive integer, are spanned by periodic data with period $`2\pi /n`$. In this sense, we roughly represent
$$𝒮(T_a)𝒮(E_a)𝒮(N_a)𝒮(S_a).$$
(105)
Therefore we conclude that for the property of preservation of translation symmetry, there is no significant distinction between the four models.
Our main interest is therefore in the reflections. As we remarked, reflection operator
$$R_0:(\overline{P}(x),\overline{Q}(x),\overline{\pi }_P(x),\overline{\pi }_Q(x))(\overline{P}(x),\overline{Q}(x),\overline{\pi }_P(x),\overline{\pi }_Q(x))$$
(106)
does not preserve the boundary conditions for $`E^3`$, Nil and Sol models, so that it is not relevant for these three models. At first sight, the similar operator acting on unbarred variables
$$R_0^{}:(P(x),Q(x),\pi _P(x),\pi _Q(x))(P(x),Q(x),\pi _P(x),\pi _Q(x)),$$
(107)
which does preserve the boundary conditions for the three models, seems to be a natural reflection operator for them. Indeed, we can regard the spatial coordinate $`x`$ for the unbarred variables as a natural coordinate in view of the translation property shown above. That is, the reflection with respect to this coordinate seems like the most natural one. Nevertheless, the reflection symmetry with respect to this operator imposed on an initial data for $`E^3`$, Nil or Sol model is not preserved under the time evolution. This fact can be seen directly from the dynamical equations for the unbarred variables. For example, those for the Nil model are given by (See Eq.(19))
$`\ddot{P}`$ $``$ $`e^{2\tau }P^{\prime \prime }e^{2P}(\dot{Q}^2e^{2\tau }(Q^{}1)^2)=0,`$
$`\ddot{Q}`$ $``$ $`e^{2\tau }Q^{\prime \prime }+2(\dot{P}\dot{Q}e^{2\tau }P^{}(Q^{}1))=0,`$ (108)
which are apparently, in contrast to the $`T^3`$ model, not invariant under the transformation $`/x/x`$, due to the factor $`(Q^{}1)`$. Similarly, those for the $`E^3`$ and Sol models are not invariant under the same transformation. As a result, in contrast to translation, there are no invariant subsets for reflection in the $`E^3`$, Nil and Sol models which naturally correspond to the invariant subset $`𝒮(R_0)`$ for the $`T^3`$ model.
However, the $`T^3`$ model admits another one-parameter family of nontrivial reflection operators $`R_{l_\theta }`$. (See Eq.(95)) The momentum part is derived from Eq.(95) by differentiating with respect to $`\tau `$ and using Eqs.(102). The result is given by
$`R_{l_\theta }:(\overline{\pi }_P(x),\overline{\pi }_Q(x))`$ $``$ $`({\displaystyle \frac{((c+\overline{Q}s)^2e^{2\overline{P}}s^2)\overline{\pi }_P+2e^{2\overline{P}}s(c+\overline{Q}s)\overline{\pi }_Q}{(c+\overline{Q}s)^2+e^{2\overline{P}}s^2}},`$ (109)
$`2s(c+\overline{Q}s)\overline{\pi }_P((c+\overline{Q}s)^2e^{2\overline{P}}s^2)\overline{\pi }_Q\left)\right|_{xx},`$
where $`c\mathrm{cos}2\theta `$, $`s\mathrm{sin}2\theta `$.
As shown in the previous subsection, for Nil and Sol models only isolated values ($`\theta =0`$ or $`\pi /2`$ for Nil, and $`\pi /4`$ or $`3\pi /4`$ for Sol) of $`\theta `$ are permissible. As a result, the invariant subsets for reflection in the Nil and Sol models are much limited compared to the $`T^3`$ model. More precisely, we have the following.
###### Proposition 9
Let $`𝒮(R_T)`$, $`𝒮(R_E)`$, $`𝒮(R_N)`$, and $`𝒮(R_S)`$ be the unions of the all invariant subsets for reflection in, respectively, the $`T^3`$, $`E^3`$, Nil, and Sol models. Then, the inclusion relation
$$𝒮(R_T)𝒮(R_E)(𝒮(R_N),𝒮(R_S))$$
(110)
holds (if neglecting the constraints corresponding to Eq.(123)).
Proof: This can be seen from $`𝒮(R_T)=𝒮(R_0)(_\theta 𝒮(R_{l_\theta }))`$, $`𝒮(R_E)=_\theta 𝒮(R_{l_\theta })`$, $`𝒮(R_N)=𝒮(R_{l_0})`$, and $`𝒮(R_E)=𝒮(R_{l_{\pi /4}})`$.
It should be stressed that the relation (110) truly reflects the dynamical properties of the models. Note that, since the (periodic) boundary conditions imposed on initial data are the same for all models if the phase space $`𝒫`$ is spanned by the unbarred variables, a map $`f`$ which preserves the boundary conditions for a model always preserves the boundary conditions for other models, too. This means that there are the same amounts of “reflection symmetric initial data” for the four models. However, the time evolution does not always preserve the symmetry of such an initial data, as we have seen. To conclude, we have clarified how much “reflection” symmetric data exist for which the symmetry is preserved under the time evolution for each model, and found, in particular, the inclusion relation (110). This may be interpreted as a manifestation of the influence of topology to dynamics.
## 5 Summary and Comments
We have made the structures of $`𝒰^2`$-manifolds clear and classified the possible topologies of them. For convenience we have split the set of these manifolds into two kinds; the first kind, those with local Killing vectors which are not degenerate everywhere, and the second kind, those with ones which are degenerate somewhere. The local Killing vectors of any $`𝒰^2`$-manifold of the second kind are actually defined globally, so that all the $`𝒰^2`$-symmetric models of the second kind are contained in the usual Gowdy models. On the other hand, it is only $`T^3`$ that is contained in the Gowdy models among the varieties of the $`𝒰^2`$-manifolds of the first kind, so that in this paper we have basically restricted ourselves to the models of the first kind.
In the case that three-manifold $`M`$ is an $`𝒰^2`$-symmetric space of the first kind, $`M`$ is a $`T^2`$-bundle over the $`S^1`$ (Theorem 2), and according to the corresponding geometric structure, $`M`$ is naturally characterized by one of $`E^3`$, Nil, or Sol. In each case, the possible topologies of $`M`$ can be classified more precisely, but if we are interested in the dynamical properties of the corresponding spacetime models it is not necessary to consider all the models one by one. Restricted to representatives for this “dynamical equivalence class,” the number of models to consider is to great extent decreased. In particular, there is three representatives for $`E^3`$ (cf. Sec.3.3), and there is only one for Nil (cf. Sec.3.1). The representatives for Sol are parameterized by only one discrete parameter $`\mathrm{q}`$ (cf. Sec.3.2).
We have given two ways of representing the metric for each of $`E^3`$, Nil and Sol models. (We distinguish between the $`E^3`$ and $`T^3`$ models. See Convention in Sec.3.3.) Note that the metrics of all the $`𝒰^2`$-symmetric spacetimes can be represented locally in the same canonical form, but with distinct boundary conditions. In this view, however, global symmetries that arise from the spatial topology tend to be unclear. On the other hand, the other way of representing the metric makes the geometric structure the $`𝒰^2`$-manifold $`M`$ admits manifest. Geometrically it is obtained by “relaxing” the corresponding locally homogeneous metric. This type of metric has another advantage that it gives rise to a natural reduction of the model; that is, since the boundary conditions are always given by the periodicity for the spatial coordinate, the spatial manifold naturally reduces to the $`S^1`$, the base space of the bundle. Note that this fact makes the identifications of the phase spaces for the varieties of the models possible.
Finally, as an application of these fundamental facts we have given the translation and reflection operators which commute with the time-evolution, and have discussed their significance. Since, as stated above, the metrics of $`𝒰^2`$-symmetric models can always be locally represented in the same form, these models are considered to have the same local dynamical properties. However differences of topologies affect the global dynamical properties of the models through the differences of the boundary conditions. In this paper we have examined whether or not a global symmetry imposed on an initial data is preserved in time, or how much there exist global symmetries which are preserved in time, for each model of the first kind. As a result, we have indeed found that there are remarkable distinctions in the properties of reflection. The freedom of reflection symmetries is the largest as for the $`T^3`$ model, and we have the inclusion relation (110). In particular, naive (even type) reflection symmetric initial data for the $`E^3`$, Nil or Sol model do not evolve maintaining its symmetry in time, in contrast to the $`T^3`$ model.
In the following we make some comments. A first one concerns another point concerning the correlations between topology and dynamics. As for the dynamics of the Gowdy $`T^3`$ model, what has been being taken an interest most is whether the AVTD conjecture is correct or not that predicts a universal behavior of the approach toward the initial singularity. This conjecture has been basically supported both analytically and numerically . However, the subtlety has been pointed out that there is a measure-zero set of spatial points where the AVTD behavior is not achieved. Here we concern ourselves with these nongeneric points. It is known that at these points $`\overline{Q}^{}=0`$. Since this condition is local, it does not depend on the topology. However, note that while points of $`\overline{Q}^{}=0`$ are inevitable in the $`T^3`$ model, since $`\overline{Q}(x)`$ is a (smooth) periodic function, such points are not necessary in the Nil and Sol models. Hence we can naturally expect that topology affects the tendency of the appearance of the nongeneric points. The points where $`\overline{Q}^{}=0`$ correspond, from Eqs.(19) and (47), to the points where
Nil: $`Q^{}=1,`$ (111)
Sol: $`(\mathrm{log}|Q|)^{}=2\mathrm{q},`$ (112)
for the unbarred variables. The unbarred $`Q(x)`$ must be periodic, but this condition does not force the existence of points such that the above conditions are fulfilled. However, since the evolution of $`Q`$ freezes when approaching the initial singularity, we can naturally expect that nongeneric points are still generated if the maximal gradient of $`Q`$ at an initial surface is large enough and if, as a result, $`Q_{\mathrm{max}}^{}>1`$ (Nil), or $`(\mathrm{log}|Q|)_{\mathrm{max}}^{}>2\mathrm{q}`$ (Sol) at a time when the evolution of $`Q`$ freezes. Indeed, this property has been observed from numerical simulations the author performed. A detailed description for this point will be reported elsewhere. In Ref. it was suggested that the nongeneric points are not generated in the Nil and Sol models, since the conditions corresponding to Eqs.(112) and (112) were written in such different forms with different choice of metric functions that the conditions seem not to be fulfilled. This claim is not correct, as seen from above. Still, we may claim that the $`E^3`$ model is most likely to generate nongeneric points.
Next we comment on the locally $`\mathrm{U}(1)`$-symmetric models that are less symmetric than the $`𝒰^2`$-symmetric models. This model admits only one spatial local Killing vector. The spatial manifold $`M`$ is again assumed to be closed. First, from an analogous analysis to the proof of Lemma 1, the Killing orbits are found to be closed and are homeomorphic to $`S^1`$, if the local Killing vector does not degenerate everywhere (that is, if the model is the “first kind”). Hence $`M`$ is naturally a Seifert fiber space , and still admits a geometric structure, which is one of $`E^3`$, Nil, $`H^2\times 𝐑`$, $`\mathrm{SL}(2,𝐑)`$, $`S^2\times 𝐑`$, or $`S^3`$ . In the case that $`M`$ is the “second kind”, i.e., the local Killing vector vanishes on somewhere, the manifold $`M`$ is obtained from a Seifert fiber space by performing a finite number of particular kind of Dehn surgeries. The resulting manifold $`M`$ is found (from one in the series of famous theorems of Thurston concerning the $`H^3`$ structure) to again admit a geometric structure. Detailed accounts and applications will be presented elsewhere.
## Appendix A Vacuum Einstein equations for Gowdy Spacetimes
In this Appendix, we summarize the standard prescription of the reduced vacuum Einstein equations for Gowdy spacetimes, together with generalizations to Nil and Sol types.
We consider the two-surface orthogonal metric (12). This is used most frequently in the literature since the Einstein equations become to great extent simpler.
More explicitly, we write
$$\mathrm{d}s^2=e^{\overline{\gamma }/2}(\mathrm{d}t^2+\mathrm{d}x^2)+\overline{R}(e^{\overline{P}}(\mathrm{d}y+\overline{Q}\mathrm{d}z)^2+e^{\overline{P}}\mathrm{d}z^2).$$
(113)
where $`\overline{P}`$ and $`\overline{Q}`$ are functions of $`t`$ and $`x`$. Then, the vacuum Einstein equations for the metric functions are explicitly given as follows. The dynamical equations are
$`\overline{P}_{tt}`$ $``$ $`\overline{P}^{\prime \prime }e^{2\overline{P}}(\overline{Q}_t^2\overline{Q}^{}{}_{}{}^{2})+\overline{R}^1(\overline{R}_t\overline{P}_t\overline{R}^{}\overline{P}^{})=0,`$ (114)
$`\overline{Q}_{tt}`$ $``$ $`\overline{Q}^{\prime \prime }+2(\overline{P}_t\overline{Q}_t\overline{P}^{}\overline{Q}^{})+\overline{R}^1(\overline{R}_t\overline{Q}_t\overline{R}^{}\overline{Q}^{})=0,`$
and
$$\overline{R}_{tt}\overline{R}^{\prime \prime }=0,$$
(115)
where subscript $`t`$ represents the derivative with respect to $`t`$ and the primes represent the derivatives with respect to $`x`$. The remaining constraint equations takes a simple form in the null coordinates $`u=tx`$ and $`v=t+x`$;
$`\overline{R}(\overline{P}_u^2+e^{2\overline{P}}\overline{Q}_u^2)`$ $`=`$ $`2\overline{R}_{uu}+\overline{R}_u(\overline{\gamma }_u+\overline{R}^1\overline{R}_u),`$
$`\overline{R}(\overline{P}_v^2+e^{2\overline{P}}\overline{Q}_v^2)`$ $`=`$ $`2\overline{R}_{vv}+\overline{R}_v(\overline{\gamma }_v+\overline{R}^1\overline{R}_v),`$ (116)
where the subscripts stand for the derivatives thereof.
The last two equations can be solved for $`\overline{\gamma }_u`$ and $`\overline{\gamma }_v`$ at any spacetime point wherever $`\overline{R}_u`$ and $`\overline{R}_v`$ do not vanish. In this case the integrability condition for $`\overline{\gamma }`$, namely $`\overline{\gamma }_{uv}\overline{\gamma }_{vu}=0`$, is automatically satisfied if the dynamical equations (114) are satisfied. Hence we can integrate the equations and obtain $`\overline{\gamma }`$ by using solutions of the (unconstrained) dynamical equations (114). If there is a spacetime point such that $`\overline{R}_u=0`$ (or $`\overline{R}_v=0`$), Eqs.(A) cannot be solved for $`\overline{\gamma }_u`$ (or $`\overline{\gamma }_v`$), so that the constraint equation (A) constrains $`\overline{P}`$ and $`\overline{Q}`$ at that spacetime point. This condition is called the “matching condition” . Explicitly,
$`\overline{P}_u^2+e^{2\overline{P}}\overline{Q}_u^2`$ $`=`$ $`2\overline{R}^1\overline{R}_{uu}\mathrm{at}\overline{R}_u=0,`$
$`\overline{P}_v^2+e^{2\overline{P}}\overline{Q}_v^2`$ $`=`$ $`2\overline{R}^1\overline{R}_{vv}\mathrm{at}\overline{R}_v=0.`$ (117)
From the observation that the left hand sides of Eqs.(A) are positive semidefinite, we find that for $`R>0`$,
$`\overline{R}_{uu}0\mathrm{at}\overline{R}_u=0,`$
$`\overline{R}_{vv}0\mathrm{at}\overline{R}_v=0.`$ (118)
An implication of these inequalities is called the “corner theorem” . Here, a corner is a point in the $`t`$-$`x`$ plane such that $`\overline{R}_u=0`$ or $`\overline{R}_v=0`$, that is, a point where the gradient of a level curve of $`\overline{R}`$ becomes null. The implication of the inequalities would be clear if observing some level curves of $`\overline{R}`$ near a corner in the $`u`$-$`v`$ coordinates. It should be stressed that these inequalities do not depend on the boundary conditions for $`\overline{P}`$, $`\overline{Q}`$, and $`\overline{\gamma }`$, so hold for any spatial topology. (This theorem, however, tacitly assume the finiteness of $`\overline{\gamma }_u`$ and $`\overline{\gamma }_v`$ at the corner. If we allow for $`\overline{R}_u\overline{\gamma }_u`$ and $`\overline{R}_v\overline{\gamma }_v`$ to remain finite at the corner with diverging $`\overline{\gamma }_u`$ and $`\overline{\gamma }_v`$, we may have counterexamples to this theorem. )
For any $`𝒰^2`$-symmetric spacetime of the first kind, the area function $`\overline{R}`$ must be periodic (cf. Eq.(3)). Hence, corners appear only in pair, implying that $`\overline{R}`$ cannot have a corner. We can therefore choose $`\overline{R}`$ spatially constant
$$\overline{R}(t,x)=t,$$
(119)
(using the remaining freedom of the coordinate transformations of type $`uF(u),vG(v)`$). In this case, the reduced Einstein equations (114) and (A) become
$`\ddot{\overline{P}}`$ $``$ $`e^{2\tau }\overline{P}^{\prime \prime }e^{2\overline{P}}(\dot{\overline{Q}}^2e^{2\tau }\overline{Q}^{}{}_{}{}^{2})=0,`$ (120)
$`\ddot{\overline{Q}}`$ $``$ $`e^{2\tau }\overline{Q}^{\prime \prime }+2(\dot{\overline{P}}\dot{\overline{Q}}e^{2\tau }\overline{P}^{}\overline{Q}^{})=0,`$
with
$`\dot{\overline{\lambda }}`$ $`=`$ $`\dot{\overline{P}}^2+e^{2\tau }\overline{P}^{\prime \prime }+e^{2\overline{P}}(\dot{\overline{Q}}^2+e^{2\tau }\overline{Q}^{}{}_{}{}^{2}),`$ (121)
$`\overline{\lambda }^{}`$ $`=`$ $`2(\overline{P}^{}\dot{\overline{P}}+e^{2\overline{P}}\overline{Q}^{}\dot{\overline{Q}}).`$
Here, we have put
$$t=e^\tau ,e^{\overline{\gamma }/2}=e^{\overline{\lambda }/2+\tau /2}.$$
(122)
Dots in these equations represent derivatives with respect to $`\tau `$ (not $`t`$). $`\overline{\lambda }`$ should be periodic for any $`𝒰^2`$-symmetric spacetime of the first kind. This leads to the constraint for $`\overline{P}`$ and $`\overline{Q}`$,
$$_0^{2\pi }(\overline{P}^{}\dot{\overline{P}}+e^{2\overline{P}}\overline{Q}^{}\dot{\overline{Q}})dx=0,$$
(123)
because of the second equation of Eqs.(121).
The functions $`\overline{P}`$ and $`\overline{Q}`$ are not periodic functions for the Nil and Sol models. Appropriate boundary conditions are obtained from Eqs.(19) and (47). Specifically,
For Nil:
$$\overline{P}(x+2\pi )=\overline{P}(x),\overline{Q}(x+2\pi )=\overline{Q}(x)2\pi .$$
(124)
For Sol:
$$\overline{P}(x+2\pi )=\overline{P}(x)+4\pi \mathrm{q},\overline{Q}(x+2\pi )=e^{4\pi \mathrm{q}}\overline{Q}(x).$$
(125)
Remark The last boundary conditions (124) and (125) do not necessarily coincide with those obtained from Eqs.(3). This is because the choice of the covering group $`\mathrm{\Gamma }`$ (cf. Eq.(10)) is different. In obtaining Eqs.(3) we have fixed $`\mathrm{\Gamma }`$ first and found the appropriate metric, but for Eqs.(124) and (125) we try to fix the universal covering metric with appropriate $`\mathrm{\Gamma }`$ as long as possible. This is the essential point in the dynamical equivalences for the Nil and Sol models presented in Sec.3.
## Appendix B Compactification of Sol
In this Appendix we show calculations needed to perform the compactification of Sol. This is basically a review of that presented in Ref.. In this reference, the form of $`\mathrm{mcg}_+(T^2)\mathrm{SL}(2,𝐙)`$ was restricted to that of Eq.(9), and the parameter $`\mathrm{q}`$ (See bellow) was not introduced. In this review we show an explicit calculation with general form of the matrix and with $`\mathrm{q}`$, though no essential difference appears.
What we do is to embed the fundamental groups into Sol$`G_{\mathrm{VI}_0}`$ (the Bianchi $`\mathrm{VI}_0`$ group). For notational convenience, letting $`a=g_2`$, $`b=g_3`$, and $`c=g_1`$ in the representation (6),
$$\pi _1=a,b,c;[a,b]=1,cac^1=a^pb^r,cbc^1=a^qb^s,$$
(126)
where
$$A\left(\begin{array}{cc}p& q\\ r& s\end{array}\right)\mathrm{SL}(2,𝐙),\text{with}\left|\mathrm{Tr}A\right|>2.$$
(127)
We represent the $`\pi _1`$-generators with their components in $`G_{\mathrm{VI}_0}`$
$$a=\left(\begin{array}{c}a_1\\ a_2\\ a_3\end{array}\right),b=\left(\begin{array}{c}b_1\\ b_2\\ b_3\end{array}\right),c=\left(\begin{array}{c}c_1\\ c_2\\ c_3\end{array}\right).$$
(128)
To avoid confusions with powers, we use subscripts to distinguish each components as above. The multiplication rule for $`G_{\mathrm{VI}_0}`$ is given by
$$\left(\begin{array}{c}\alpha \\ \beta \\ \gamma \end{array}\right)\left(\begin{array}{c}x\\ y\\ z\end{array}\right)=\left(\begin{array}{c}\alpha +x\\ \beta +e^{\mathrm{q}\alpha }y\\ \gamma +e^{\mathrm{q}\alpha }z\end{array}\right),\left(\begin{array}{c}\alpha \\ \beta \\ \gamma \end{array}\right)^1=\left(\begin{array}{c}\alpha \\ e^{\mathrm{q}\alpha }\beta \\ e^{\mathrm{q}\alpha }\gamma \end{array}\right),$$
(129)
where $`\mathrm{q}>0`$ is a parameter (, which has nothing to do with the elements in the matrix $`A`$). Left invariant one-forms are then given by Eq.(44).
First, note the first components of the relations
$$cac^1=a^pb^r,cbc^1=a^qb^s.$$
(130)
Noting that the first component of the product of two elements of Sol is simply given by the sum of the first components of the elements, we obtain $`a_1=pa_1+rb_1`$ and $`b_1=qa_1+sb_1`$, that is,
$$\left(\begin{array}{cc}p1& r\\ q& s1\end{array}\right)\left(\begin{array}{c}a_1\\ b_1\end{array}\right)=\left(\begin{array}{c}0\\ 0\end{array}\right).$$
(131)
The determinant for the matrix appearing in the left hand side is given by $`detA\mathrm{Tr}A=1\mathrm{Tr}A`$, which is less than $`1`$ (when $`\mathrm{Tr}A>2`$) or greater than $`3`$ (when $`\mathrm{Tr}A<2`$), from the assumptions on the matrix $`A`$. Hence, the matrix has its inverse, implying
$$a_1=b_1=0.$$
(132)
Under this condition the relation $`[a,b]=1`$ becomes trivial. Moreover, we can at once calculate
$$a^p=\left(\begin{array}{c}0\\ pa_2\\ pa_3\end{array}\right),b^r=\left(\begin{array}{c}0\\ rb_2\\ rb_3\end{array}\right),cac^1=\left(\begin{array}{c}0\\ a_2e^{\mathrm{q}c_1}\\ a_3e^{\mathrm{q}c_1}\end{array}\right),\text{etc.},$$
(133)
so that the second and third components of the relations (130) are given by
$$A\left(\begin{array}{c}a_2\\ b_2\end{array}\right)=e^{\mathrm{q}c_1}\left(\begin{array}{c}a_2\\ b_2\end{array}\right),A\left(\begin{array}{c}a_3\\ b_3\end{array}\right)=e^{\mathrm{q}c_1}\left(\begin{array}{c}a_3\\ b_3\end{array}\right).$$
(134)
These reveal that $`\left(\begin{array}{c}a_2\\ b_2\end{array}\right)`$ and $`\left(\begin{array}{c}a_3\\ b_3\end{array}\right)`$ are eigenvectors of $`A`$, and $`e^{\pm \mathrm{q}c_1}`$ are the eigenvalues. The characteristic polynomial is $`\lambda ^2\mathrm{Tr}A\lambda +detA=0`$, or equivalently,
$$\lambda ^2\mathrm{Tr}A\lambda +1=0.$$
(135)
Since the eigenvalues $`e^{\pm \mathrm{q}c_1}`$ must be positive, the embedding is found to be possible only for the case $`\mathrm{Tr}A>2`$. However, even for the case $`\mathrm{Tr}A<2`$, if the map
$$h:(x,y,z)(x,y,z)$$
(136)
can be regarded as an isometry, we can perform the embedding by putting $`c=h\left(\begin{array}{c}c_1\\ c_2\\ c_3\end{array}\right)`$, where the column vector in the right hand side is an element of Sol, and $``$ represents the composition of maps. It is an easy task to retrace the calculation above. We will find in particular that the eigenvalues of $`A`$ must equal to $`e^{\pm \mathrm{q}c_1}`$ rather than $`e^{\pm \mathrm{q}c_1}`$.
Finally, the formal solutions for embedding are as follows.
(i) Case of $`\mathrm{Tr}A>2`$: Letting the eigensystems (the pairs of an eigenvalue and the corresponding normalized eigenvector) for the matrix $`A`$ be
$$\{e^{\mathrm{q}c_1},\left(\begin{array}{c}u_2\\ u_3\end{array}\right)\},\{e^{\mathrm{q}c_1},\left(\begin{array}{c}v_2\\ v_3\end{array}\right)\},$$
(137)
the solution is
$$a=\left(\begin{array}{c}0\\ \beta u_2\\ \gamma v_2\end{array}\right),b=\left(\begin{array}{c}0\\ \beta u_3\\ \gamma v_3\end{array}\right),c=\left(\begin{array}{c}c_1\\ c_2\\ c_3\end{array}\right),$$
(138)
where $`\beta `$, $`\gamma `$, $`c_2`$, and $`c_3`$ are parameters. Note that, while the product $`qc_1`$ is directly connected with the eigenvalues of $`A`$, we can freely specify one of $`q`$ or $`c_1`$.
(ii) Case of $`\mathrm{Tr}A<2`$: Letting the eigensystems be
$$\{e^{\mathrm{q}c_1},\left(\begin{array}{c}u_2\\ u_3\end{array}\right)\},\{e^{\mathrm{q}c_1},\left(\begin{array}{c}v_2\\ v_3\end{array}\right)\},$$
(139)
the solution is
$$a=\left(\begin{array}{c}\beta u_2\\ \gamma v_2\\ 0\end{array}\right),b=\left(\begin{array}{c}\beta u_3\\ \gamma v_3\\ 0\end{array}\right),c=h\left(\begin{array}{c}c_1\\ c_2\\ c_3\end{array}\right),$$
(140)
where $`\beta `$, $`\gamma `$, $`c_2`$ and $`c_3`$ are parameters. As in the case (i), one of $`q`$ or $`c_1`$ can be freely specified.
Remark We have seen that the fundamental group for $`\mathrm{Tr}A>2`$ can be embedded into $`G_{\mathrm{VI}_0}`$. Also, we have seen that the fundamental group for $`\mathrm{Tr}A<2`$ can be embedded into $`G_{\mathrm{VI}_0}^2`$, the group generated by $`G_{\mathrm{VI}_0}`$ and the $`𝐙_2`$-factor $`h`$. It is easy to see that the same groups can be embedded into the smaller group $`H_{\mathrm{VI}_0}G_{\mathrm{VI}_0}`$ which is given by making the first component of $`G_{\mathrm{VI}_0}`$ discrete like $`c_1𝐙`$, or into $`H_{\mathrm{VI}_0}^2`$, where $`H_{\mathrm{VI}_0}^2G_{\mathrm{VI}_0}^2`$ is the group generated by $`H_{\mathrm{VI}_0}`$ and the map $`h`$. Since $`H_{\mathrm{VI}_0}^2`$ is the isometry group of the universal covering spacetime (45) (with the appropriate choice of $`\mathrm{q}`$, as shown in Sec.3.2), the same embeddings are possible for this (inhomogeneous) spacetime.
## Acknowledgments
The author thanks Professor K. Tomita for explaining a subtle point in the Gowdy equations to him. He acknowledges financial support from the Japan Society for the Promotion of Science. |
warning/0003/physics0003027.html | ar5iv | text | # Magnetic Field Dependence of Ultracold Inelastic Collisions near a Feshbach Resonance
## Abstract
Inelastic collision rates for ultracold <sup>85</sup>Rb atoms in the F=2, m<sub>f</sub>=–2 state have been measured as a function of magnetic field. At 250 Gauss (G), the two- and three-body loss rates were measured to be K<sub>2</sub>=1.87$`\pm `$0.95$`\pm `$0.19 x 10<sup>-14</sup> cm<sup>3</sup>/s and K<sub>3</sub>=4.24$`{}_{0.29}{}^{+0.70}\pm `$0.85 x 10<sup>-25</sup> cm<sup>6</sup>/s respectively. As the magnetic field is decreased from 250 G towards a Feshbach resonance at 155 G, the inelastic rates decrease to a minimum and then increase dramatically, peaking at the Feshbach resonance. Both two- and three-body losses are important, and individual contributions have been compared with theory.
Feshbach resonances have recently been observed in a variety of cold atom interactions, including elastic scattering , radiative collisions and enhanced inelastic loss , photoassociation , and the loss of atoms from a Bose-Einstein condensate (BEC) . By changing the magnetic field through the resonance, elastic collision rates can be changed by orders of magnitude and even the sign of the atom-atom interaction can be reversed . The work in Ref. has received particular attention because the BEC loss rates were extraordinarily high and several proposals for exotic coherent loss processes have been put forward . However, even ordinary dipole relaxation and three-body recombination are expected to show dramatic enhancements by the Feshbach resonance , and the calculations of these ordinary enhancements have never been tested. This has left many outstanding questions as to the nature of inelastic losses near a Feshbach resonance. How large are the dipole relaxation and three-body recombination near the Feshbach resonance and how accurate are the calculations of these quantities? How much of the observed condensate losses in Ref. are due to these more traditional mechanisms and how much arise from processes unique to condensates? How severe are the“severe limitations” that inelastic loss puts on the use of Feshbach resonances to change the s-wave scattering length in a BEC? The nature of the inelastic losses near Feshbach resonances also has important implications for efforts to create BEC in <sup>85</sup>Rb, because these losses play a critical role in determining the success or failure of evaporative cooling. Thus it is imperative to better understand the nature of inelastic collisions between ground state atoms near a Feshbach resonance.
In this paper we present the study of the losses of very cold <sup>85</sup>Rb atoms from a magnetic trap as a function of density and magnetic field. There are two types of inelastic collisions that induce loss from a magnetic trap. The first is dipolar relaxation where two atoms collide and change spin states. The second process is three-body recombination, where three atoms collide and two of those atoms form a molecule. Measuring the losses as a function of density has allowed us to determine the two-body and three-body inelastic collision rates, while measuring the variation of these losses as a function of magnetic field has allowed us to find out how the Feshbach resonance affects them. In contrast to the work of Ref. we observe a pronounced dip in the inelastic losses near the Feshbach resonance. This will make it possible to create <sup>85</sup>Rb BECs with a positive scattering length.
To study these losses we needed a cold, dense <sup>85</sup>Rb sample. This was obtained through evaporative cooling with a double magneto-optic trap (MOT) system as described in Ref. . The first MOT repeatedly collected atoms from a background vapor and those atoms were transferred to another MOT in a low-pressure chamber. Once the desired load size was achieved, the MOTs were turned off and a baseball-type Ioffe-Pritchard magnetic trap was turned on, resulting in a trapped atom sample of about 3x10<sup>8</sup> F=2, m<sub>f</sub>=–2 <sup>85</sup>Rb atoms at 45 $`\mu K`$. Forced radio-frequency (rf) evaporation was used to increase the density of the atoms while decreasing the temperature. Because of the high ratio of inelastic/elastic collision rates in <sup>85</sup>Rb and the dependence of the ratio on magnetic field ($`B`$) and temperature ($`T`$), trap conditions must be carefully chosen to achieve efficient cooling. We evaporated to the desired temperature and density (typically 3x10<sup>11</sup> cm<sup>-3</sup> and 500-700 $`nK`$) at a final field of $`B`$=162 Gauss (G) .
We then adiabatically changed the DC magnetic field to various values and measured the density of the trapped atom cloud as a function of time. We observed the clouds with both nondestructive polarization-rotation imaging and destructive absorption imaging. In both cases, the cloud was imaged onto a CCD array to determine the spatial size and number. The nondestructive method allowed us to observe the time evolution of the number and spatial size of a single sample. The destructive method required us to prepare many samples and observe them after different delay times, but had the advantages of better signal to noise for a single image and larger dynamic range. A set of nondestructive imaging data is shown in Fig. 1. In addition to these techniques, we made a redundant check of the number by recapturing the atoms in the MOT and measuring their fluorescence .
We measured the time-dependence of the number of atoms and the spatial size, and from these measurements determined the density. We assigned a 10% systematic error to our density determination, based primarily on the error in our measurement of number. The value of $`B`$ was measured in the same way as in Ref. : the rf frequency at which the atoms in the center of the trap were resonantly spin-flipped was measured and the Breit-Rabi equation was then used to determine the magnitude of $`B`$. The magnetic field width of the clouds scaled as $`T^{1/2}`$ and was 0.39 G FWHM at 500 $`nK`$.
In each data set we observed the evolution of the sample while a significant fraction (20–35%) of the atoms was lost. The temperature of the sample also increased as a function of time. This heating rate scaled with the inelastic rates and therefore varied with $`B`$. Since the volume and temperature in a magnetic trap are directly related, the change in the volume (typically 50%) with time was fit to a polynomial—usually a straight line was sufficient within our precision. The number (N) as a function of time was then fit to the sum of three-body and two-body loss contributions given by
$$\dot{N}=K_2nNK_3n^2N\frac{N}{\tau }.$$
(1)
Here $`n=\frac{1}{N}n^2(𝐱)d^3x`$ is the density-weighted density and $`n^2=\frac{1}{N}n^3(𝐱)d^3x`$. The time evolution of the volume is contained in the computation of the density as a function of time. The background loss rate is given by $`\tau `$. It was independently determined by looking at low-density clouds for very long times. Typically, it was 450 seconds, with some slight dependence on $`B`$ ($`<`$20%), particularly near the Feshbach resonance.
The change in the combined inelastic rates as a function of $`B`$ around the Feshbach resonance is shown in Fig. 2(a). All of the inelastic data shown in Fig. 2(a) were taken with initial temperatures near 600 $`nK`$. The points less than 157 G were taken with initial densities within 10% of 1x10<sup>11</sup>cm<sup>-3</sup>, while the points higher than 157 G were taken with initial densities that were 2.7x10<sup>11</sup>cm<sup>-3</sup>. The decrease in initial density below 157 G was due to ramping through the high inelastic loss region after forming the sample at 162 G. Because of this density variation, we have used a different weighting, parameter $`\beta `$, in the sum of the two- and three-body rate constants for the two regions in Fig. 2(a).
Also shown in Fig. 2(a) is the elastic rate determined previously . This elastic rate was measured by forming atom samples similar to the ones in this work (but at much lower density), forcing them out of thermal equilibrium, and then observing their reequilibration time. The shape of the inelastic rate vs. $`B`$ roughly follows that of the elastic rate vs. $`B`$. In particular, the peak of the inelastic rate occurs at 155.4$`\pm `$0.5 G, identical to the position of the elastic peak at 155.2$`\pm `$0.4 G within the error. Also, just as is the case for the elastic rates, the inelastic loss rates around the Feshbach resonance vary by orders of magnitude. The peak in the inelastic rate is much less symmetric, however. Another interesting feature is that the loss rates not only increase near the elastic rate peak but decrease near its minimum. The field where the loss is minimum, 173.8$`\pm `$2.5 G, is higher than the minimum of the elastic rate at 166.8$`\pm `$0.3 G.
Even though the two- and three-body inelastic rates have different density dependencies, it is difficult to separate them. Figure 1 shows how a purely two-body or purely three-body loss curve fits equally well to a typical data set. An excellent signal-to-noise ratio or, equivalently, data from a large range of density are required to determine whether the loss is three-body, two-body, or a mixture of the two. To better determine density dependence, we decreased the initial density by up to a factor of 10 for several $`B`$ values. For fields with relatively high loss rates, the signal-to-noise ratio was adequate to distinguish between two- and three-body loss. However, where the rates were lower this was not the case.
Figure 2(b) shows the two-body inelastic rate determined from the density-varied data and Fig. 2(c) shows the same for the three-body rate. The character of the inelastic loss clearly changes as one goes from higher to lower field (right to left in Fig. 2). On the far right of the graph, at $`B`$ = 250 G, the inelastic rate is dominated by a three-body process. From $`B`$ = 250 G down to $`B`$ = 174 G the inelastic rate decreases to a minimum and then begins to increase again. Near the minimum, we could not determine the loss character, but at 162 G the losses are dominated by a two-body process at this density. At 158 G the two-body process is still dominant and rising rapidly as one goes toward lower field. However, by $``$157 G it has been overtaken by the three-body recombination that is the dominant process at 157-145 G. At lower fields, both two-body and three-body rates contribute significantly to the total loss at these densities.
Along with varying the density, the initial temperature was varied at $`B`$ = 145,156,160, and 250 G. There was no significant rate change in the loss rate for temperatures between 400 and 1000 $`nK`$ for the 160 and 250 G points. The combined loss rate increased by a factor of 8 at 156 G from 1 $`\mu K`$ to 400 $`nK`$, and by a factor of 2 at 145 G. This temperature dependence near the peak is expected in analogy to the temperature dependence of the elastic rates . The fact that the loss rates near the peak exhibit both a two- and three-body character plus a temperature dependence and considerable heating make interpreting the data challenging. This introduces additional uncertainty that is included in the error bars in Fig. 2(a-c).
A calculated dipolar loss rate is compared with the data in Fig. 2(b). This was calculated by finding the S-matrix between the trapped and untrapped states using carefully determined Rb-Rb molecular potentials . Keeping in mind the significance of error bars on a log plot, the agreement is reasonably good. In the region between $`B`$ = 160 and 167 G where the determination of K<sub>2</sub> is not complicated by three-body loss and temperature dependence, the agreement is particularly good.
Likewise, we show a prediction of the three-body recombination rate in Fig. 2(c) . It is predicted that the recombination rate scales as the s-wave scattering length to the fourth power (a$`{}_{}{}^{4}{}_{s}{}^{}`$) for both positive and negative a<sub>s</sub>, although with a smaller coefficient for the positive case . Since a<sub>s</sub> varies across the Feshbach resonance (see Fig. 2(a), upper plot), the recombination rate is expected to change . Temperature dependence has not been included in the prediction. Qualitatively, the main features of the predicted three-body recombination match the data. The three-body rate decreases as a<sub>s</sub> decreases, and it increases rapidly where a<sub>s</sub> diverges at $`B`$ = 155 G. From 155 G to 167 G, a<sub>s</sub> is positive and the three-body loss is much smaller than it is at $`B`$ fields with a comparable negative a<sub>s</sub>, as expected from theory. This overall level of agreement is reasonable given the difficulties and approximations in calculating three-body rates.
The marked dependence on $`B`$ of both the elastic and inelastic rates has important implications for the optimization of evaporative cooling to achieve BEC in <sup>85</sup>Rb. As the density increases with evaporative cooling, the absolute loss rate, which is already much greater than in <sup>87</sup>Rb and Na, becomes larger. However, the important ratio of elastic rate to inelastic rate is both temperature and density dependent and magnetic field dependent. The fact that the losses are dramatically lower at fields above the Feshbach resonance suggests that it should be possible to devise an evaporation path that will lead to the creation of an <sup>85</sup>Rb BEC .
The <sup>85</sup>Rb Feshbach resonance has a profound effect on the two- and three-body inelastic rates, changing them by orders of magnitude. Far from the resonance at 250 G, the two-body and three-body loss rates are measured to be K<sub>2</sub> = 1.87$`\pm `$0.95$`\pm `$0.19 x 10<sup>-14</sup> cm<sup>3</sup>/s and K<sub>3</sub>=4.24$`{}_{0.29}{}^{+0.70}\pm `$0.85 x 10<sup>-25</sup> cm<sup>6</sup>/s respectively. Here the first error is the statistical error and the second in the systematic uncertainty due to the number determination. Near the resonance, the two- and three-body rates change by orders of magnitude. The dependence of the inelastic rates on magnetic field is similar in structure to the dependence of the elastic rate: the maxima in both rates occur at the same field, while the minima are close but do not coincide. The total loss is a complicated mixture of both two- and three-body loss processes. They have different dependencies on field so both have field regions in which they dominate.
We are pleased to acknowledge useful discussions and with Eric Cornell, Jim Burke, Jr., Chris H. Greene, and Carl Williams. We also thank the latter three and their colleagues for providing us with loss predictions. This research has been supported by the NSF and ONR. One of us (S. L. Cornish) acknowledges the support of a Lindemann fellowship. |
warning/0003/hep-ph0003328.html | ar5iv | text | # Cosmology of Brane-Bulk Models in Five Dimensions
## I Introduction
Over the last two years there has been a lot of interest in models where our universe is a 3-brane (a hyper-surface) embedded in a higher dimensional bulk. The standard model particles are confined to the brane whereas gravity propagates in the bulk. Such models were conjectured early on as interesting possibilities and have recently been argued as plausible solutions of type I string theories . One of the attractive features of these models is the intriguing possibility that the fundamental scale, $`M`$, identified as the string scale, could be lower than the Planck scale, $`M_P\mathrm{}=(8\pi G_N)^{1/2}`$, by several orders of magnitude , perhaps even of TeV range . This last possibility may provide a new way to solve the hierarchy problem between the electroweak scale and the scale of gravity i.e. both scales being of same order, there is no hierarchy to worry about. The large value of the Planck scale in this picture owes its origin to the existence of very large hidden extra dimensions in nature. The price one has to pay is of course that now one has to understand why the extra dimension(s) is(are) so large. The relation between the fundamental scale, the $`M_P\mathrm{}`$ and the volume of the external space, $`V_n`$, is given by
$$M^{2+n}V_n=M_P\mathrm{}^2.$$
(1.1)
This picture leads to a modification of the inverse square law of gravity at small distances, $`rV_n^{1/n}`$ and can therefore be probed experimentally. If $`M`$ is in the TeV range, string theories become accessible to collider tests. All these make the idea phenomenologically quite attractive. It is therefore interesting to study the cosmology of these models.
One of the first things that one needs to know in order to study the cosmology of these models is the time dependence of the metric. In addition, the metric will also have a dependence on the bulk coordinate $`y`$, even when the bulk is totally empty, simply because the presence of a brane induces a nonzero curvature. It turns out that this fact leads to a variety of interesting consequences for the cosmology of such models . In particular, the new bulk-brane picture seems to drastically modify the standard time evolution law in the brane-confined universe. The five dimensional Einstein equations first studied in implies that the Hubble parameter $`H`$ is proportional to the density on the brane, $`\rho `$, instead of the usual $`H\sqrt{\rho }`$ of the standard big bang cosmology. Since the successes of the standard cosmology such as nucleosynthesis and the common understanding of subsequent evolution rely crucially on the assumption of $`H\sqrt{\rho }`$, a great deal of work has been devoted to understand how this standard behaviour of $`H`$ can be recovered. Some ideas proposed to solve this problem include cancellation of bulk and brane cosmological constants , consideration of a non zero thickness of the brane etc. It has also been noted that Friedmann equation could be recovered on the basis of fine tuning of cosmological constants even when the extra dimension is not stable and no matter in the bulk. However, in the process it has been found that a free parameter $`𝒞`$, a constant of integration, appears in this equation. It contributes to the evolution of the Hubble parameter in the form of an effective radiation term, jeopardizing the cosmic scale parameters. The same arbitrary parameter has been used to study cosmological evolution even when the bulk radius is stable . We reexamine this in this paper. Writing the equations governing the evolution in time, $`t`$ and the fifth coordinate, $`y`$ of the various parameters defining the metric (the analogs of the Friedman-Robertson-Walker (FRW) scale factor), we show that it is possible to get a solution for the metric as a function of the bulk coordinate without any arbitrary constant, if the bulk radius is stable. This is the main result of this brief note. We then point out some of the implications of this result.
This paper is arranged as follows: in section 2, we review Einstein equations in a diagonal metric, commonly used in the literature; in section 3, we look for solutions of Einstein equation keeping the bulk radius stable and obtain an explicit form for the metric which involves only the known FRW scale factor and no extra parameters. We recover the known behaviour for the Hubble parameter, i.e. $`H`$ depends linearly on the brane energy density. We also show that the equation involving the brane-space like components can be derived from the other equations; they reduce to the acceleration equation for the FRW scale factor. Finally, we reconsider some examples already studied in and rederive the exact solutions for the metric in the prescence of a cosmological constant in the bulk. In section 5, we make some remarks on the general nonfactorizable nature of this solutions.
## II Basic framework
The basic framework for our discussion is a five dimensional space-time where a flat (zero thickness) brane is localized at the position identified as $`y=0`$ along the fifth dimension. Although in this paper, we assume that the extra dimension is compact, this is not crucial for our conclusions and our results also apply when the extra dimension is noncompact.
Since we are interested in the brane cosmology, we start by adopting the cosmological principle of isotropy and homogeneity in the three space dimensions of the brane. The presence of the brane clearly breaks the isotropy along the fifth dimension and this is reflected in the explicit $`y`$ dependence of the metric tensor which we choose to have the following form:
$$ds^2=n^2(y,t)dt^2+a^2(y,t)\gamma _{ij}dx^idx^j+b^2(y,t)dy^2.$$
(2.1)
Here $`\gamma _{ij}=f(r)\delta _{ij}`$, with $`f^1(r)=1kr^2`$ being the usual Robertson-Walker curvature term, where $`k=1,0,1`$; $`t`$ and $`x^i`$, $`i=1,2,3`$ are the time and space-like coordinates along the brane respectively.
The five dimensional Einstein equations take the form
$$G_{AB}=R_{AB}\frac{1}{2}g_{AB}R=\kappa _5^2\left[\widehat{T}_{AB}+T_{\mu \nu }\delta _A^\mu \delta _B^\nu \delta (by)\right],$$
(2.2)
where $`\kappa _5^2=8\pi G_{(5)}=M^3`$ is the five dimensional coupling constant of gravity, $`R_{AB}`$ is the five dimensional Ricci tensor and $`R`$ the scalar curvature, $`A,B=0,1,2,3,4`$ and $`\mu ,\nu =0,1,2,3`$. In conformity with the usual practice, we identify the mass parameter $`M`$ with the string scale. In the last expression the various source terms have been explicitly separated. $`T_{\mu \nu }`$ and $`\widehat{T}_{AB}`$ represent the stress-energy-momentum tensors of the brane and bulk respectively. For the scenarios we are going to discuss from now on, it is sufficient to work in the perfect fluid approximation to those tensors
$`\widehat{T}^A_B`$ $`=`$ $`diag(\rho _B,P_B,P_B,P_B,P_T),`$ (2.3)
$`T_{}^{\mu }{}_{\nu }{}^{}`$ $`=`$ $`diag(\rho _b,p_b,p_b,p_b).`$ (2.4)
where $`\rho _B`$, $`P_B`$ and $`P_T`$ represent the densities and pressure on the bulk and, respectively, $`\rho _b`$ and $`p_b`$ are those on the brane. Notice that assuming $`\widehat{T}_{04}=0`$ avoids the complication of a matter flux along the fifth dimension.
So far, most of the attention paid to this model in literature appears to have focussed on the case where vacuum energy (a cosmological constant) is the only component of the bulk stress tensor. In general however, all the components of $`T_B^A`$ could be functions of time, if they are in the brane and could depend on both $`t`$ and $`y`$ if they describe the bulk . Therefore, in our analysis, we will keep the energy-momentum tensor dependent on $`t`$ (and $`y`$ for the bulk terms). It is of course straightforward to see that more branes could be considered by including their corresponding stress tensors in Eq. (2.2) and our discussion below easily generalises to this case.
In order to solve the Einstein equations on the presence of the delta-function type densities, we proceed as follows. First we observe that the brane divides the bulk into two different domains, where the only source is $`\widehat{T}_{AB}`$. We then solve the equations in each domain separately, and impose the boundary conditions at the brane to get the global solution. This also helps to define the metric on the brane itself. First, the metric tensor, $`g`$, clearly should be continuous, i.e. the solutions must satisfy
$$g_{AB}(y=0^{})=g_{AB}(y=0^+).$$
(2.5)
Next, as has already been noted, since $`G_{AB}`$ involves up to second derivatives on the metric tensor with respect of $`y`$, we must use them to match the delta function distributions . Technically speaking this means that the extrinsic curvature $`K_{AB}`$ in the Gauss-Codacci formulation should be discontinuous at the position of the branes. Then, by integrating Eq. (2.2) at both sides of the brane, we will get matching conditions for the first derivatives. This leads to the constraint
$$_0^{}^{0+}𝑑ybG_{\mu \nu }=\kappa _5^2T_{\mu \nu }.$$
(2.6)
To evaluate this integral we should assume that all other terms not involving second derivatives on $`y`$ are finite.
Typically, a parity symmetry $`P:yy`$ is assumed as in the Horava-Witten model . Physically, this symmetry could be seen as a residual effect of the broken isotropy along the fifth dimension, as in the $`S^1/Z_2`$ orbifold construction of ref.. We will assume this hereafter. In the presence of other branes that explicitly break this symmetry, our discussions have to be reconsidered depending on the number of branes.
Let us now proceed to the details. Using the above form of the bulk stress tensor, the non trivial components of the Einstein equations (away from the brane) are given as
$`G_{00}`$ $`=`$ $`3\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\right){\displaystyle \frac{n^2}{b^2}}\left[{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}\right)+{\displaystyle \frac{a^{\prime \prime }}{a}}\right]+k{\displaystyle \frac{n^2}{a^2}}\right\}=\kappa _5^2n^2\rho _B,`$ (2.7)
$`G_{ij}`$ $`=`$ $`{\displaystyle \frac{a^2}{b^2}}\left\{{\displaystyle \frac{a^{}}{a}}\left(2{\displaystyle \frac{n^{}}{n}}+{\displaystyle \frac{a^{}}{a}}\right){\displaystyle \frac{b^{}}{b}}\left({\displaystyle \frac{n^{}}{n}}+2{\displaystyle \frac{a^{}}{a}}\right)+2{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{n^{\prime \prime }}{n}}\right\}\gamma _{ij}+`$ (2.9)
$`{\displaystyle \frac{a^2}{n^2}}\left\{{\displaystyle \frac{\dot{a}}{a}}\left(2{\displaystyle \frac{\dot{n}}{n}}{\displaystyle \frac{\dot{a}}{a}}\right)+{\displaystyle \frac{\dot{b}}{b}}\left({\displaystyle \frac{\dot{n}}{n}}2{\displaystyle \frac{\dot{a}}{a}}\right)2{\displaystyle \frac{\ddot{a}}{a}}{\displaystyle \frac{\ddot{b}}{b}}\right\}\gamma _{ij}k\gamma _{ij}=\kappa _5^2a^2P_B\gamma _{ij},`$
$`G_{04}`$ $`=`$ $`3\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{a}^{}}{a}}\right)=0,`$ (2.10)
$`G_{44}`$ $`=`$ $`3\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^2}{n^2}}\left[{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{n}}{n}}\right)+{\displaystyle \frac{\ddot{a}}{a}}\right]k{\displaystyle \frac{b^2}{a^2}}\right\}=\kappa _5^2b^2P_T;`$ (2.11)
where primes (dots) are used to denote derivatives with respect to $`y`$ ($`t`$). This system of equations is supplemented by the Bianchi identity $`\widehat{T}_{B;A}^A=0`$, which translates into the conservation laws
$`\dot{\rho }_B+3{\displaystyle \frac{\dot{a}}{a}}\left(\rho _B+P_B\right)+{\displaystyle \frac{\dot{b}}{b}}\left(\rho _B+P_T\right)`$ $`=`$ $`0`$ (2.12)
$`P_T^{}+3{\displaystyle \frac{a^{}}{a}}\left(P_TP_B\right)+{\displaystyle \frac{n^{}}{n}}\left(P_T+\rho _B\right)=0`$ (2.13)
Next, by using Eq. (2.6), the following boundary conditions are easily obtained
$`{\displaystyle \frac{\mathrm{\Delta }a^{}}{ab}}|_0`$ $`=`$ $`{\displaystyle \frac{\kappa _5^2}{3}}\rho _b,`$ (2.14)
$`{\displaystyle \frac{\mathrm{\Delta }n^{}}{nb}}|_0`$ $`=`$ $`{\displaystyle \frac{\kappa _5^2}{3}}(3p_b+2\rho _b);`$ (2.15)
where the left hand side of the above equations has to be evaluated at the position of the brane, and the function $`\mathrm{\Delta }a^{}(0):=a^{}(0^+)a^{}(0^{})`$ give the size of the jump of the derivative of $`a(y)`$. The same applies for $`\mathrm{\Delta }n^{}`$. Since we are assuming $`P`$ parity, the jump on the above equations could be expressed in terms of the limiting value on one side of the brane, for instance by $`\mathrm{\Delta }a^{}(0)=2a^{}(0^+)`$, and a similar relation for $`\mathrm{\Delta }n^{}(0)`$.
Let us notice that if we use the above boundary conditions, we may evaluate equation (2.10) on the brane to get the conservation equation
$$\dot{\rho }_b+3\frac{\dot{a}}{a}(p_b+\rho _b)=0.$$
(2.16)
This is a general result that is independent of the bulk content. In the subsequent discussion, we will assume that the bulk is stable and use the freedom of coordinate redefinition to set $`b=1`$. This condition will simplify the Einstein equations making it easier to extract its physical meaning, as we will see later.
## III Master equations of brane cosmology
Once we assume that the fifth dimension is stable, we can follow standard procedure as in four dimensional cosmology to reduce the equations given in the last section to a minimal set. First, notice that the conservation laws (2.12), (2.13) and (2.16), will have $`\dot{b}=0`$. Taking $`b=1`$, we notice that the $`G_{00}`$ component (Eq. 2.7) of the Einstein equations is already the equivalent of the Friedmann equation, with the Hubble parameter defined as a function of $`y`$. However, the presence of the brane will require that the term with second derivative be regularized by extracting the divergent part. The Friedmann equation then has the form
$$H^2(y):=\left(\frac{\dot{a}}{a}\right)^2=\frac{\kappa _5^2}{3}n^2\rho _B+n^2\left[\left(\frac{a^{}}{a}\right)^2+\frac{a_R^{\prime \prime }}{a}\right]k\frac{n^2}{a^2};$$
(3.1)
where $`a_R^{\prime \prime }`$ stands for the regular part of the function. Clearly, we recognize an expression similar to that given early . However, let us stress that there is no unknown constant of integration as in . Also, we may identify the regular term as the contribution of the Weyl tensor of the bulk . Since this expression is continuous, thanks to parity symmetry, we may evaluate it on the brane to get the effective Friedmann equation of our universe, where the time component of the metric tensor is chosen to be $`n_0=1`$:
$$H_0^2=\frac{\kappa _5^2}{3}\rho _{B0}+\left(\frac{\kappa _5^2}{6}\rho _b\right)^2+\left(\frac{a_R^{\prime \prime }}{a}\right)_0\frac{k}{a_0^2},$$
(3.2)
where the subindex $`0`$ stand for the evaluation at $`y=0`$. As already known, this expression has the squared dependence on the brane density. On the other hand, it also has a dependence on the metric outside (i.e. in the bulk), through $`a_R^{\prime \prime }`$, which however could be evaluated, once we solve for $`a`$ as a function of $`y`$. Particularly, for simple cases as an empty bulk, the $`y`$ dependence can be explicitly extracted and this term evaluated.
Next, let us consider the equation involving $`G_{ij}`$, (2.9). Again, we separate the singular and the regular parts of the second derivative term and by introducing (3.1) we reduce this equation to a form equivalent to the acceleration equation given by
$$\frac{\ddot{a}}{a}=\frac{\kappa _5^2}{6}\left(3P_B+\rho _B\right)n^2+n^2\left(\frac{a^{}}{a}\frac{n^{}}{n}\right)+H\frac{\dot{n}}{n}+\frac{n^2}{2}\left[\frac{a_R^{\prime \prime }}{a}+\frac{n_R^{\prime \prime }}{n}\right].$$
(3.3)
Note that, as in the standard FRW cosmology, this is not an independent equation. Indeed, it can be derived by taking the time derivative of (3.1) and combining that expression with the energy conservation law (2.12) and $`G_{04}`$ equation (2.10). This derivation holds regardless of whether the bulk radius is constant or changing with time. As a result, this equation does not provide any extra information, but it will be useful in what follows.
We now turn to the equation involving $`G_{44}`$, (2.11). In conjunction with the equation involving $`G_{04}`$, this represents the new ingredient of the brane cosmology and can be a window to understand the $`y`$ dependence of the cosmological parameters. Notice that our procedure is in contrast to that used in previous works where the Friedmann equation has been obtained from $`G_{44}`$ component. If we substitute Eqs. (3.1) and (3.3) in the equation for the $`G_{44}`$ component (Eq. (2.9)), we find a simple equation that governs the behaviour of $`a(y,t)`$ and $`n(y,t)`$ on the bulk (into each one of the domains). It is given by
$$3\frac{a^{\prime \prime }}{a}+\frac{n^{\prime \prime }}{n}=\frac{\kappa _5^2}{3}\left(3P_B2P_T\rho _B\right).$$
(3.4)
This expression is supplemented by $`G_{04}`$ which trivialy reduces to
$$n(y,t)=\lambda (t)\dot{a}(y,t),$$
(3.5)
where $`\lambda (t)`$ is an arbitrary function of time. Using the freedom of fixing the gauge on the coordinate system, we can set $`n_0=1`$ in which case we get $`\lambda =\dot{a}_0^1`$. However in most of the results this choice is not necessary at all.
Let us emphasize that Eqs. (3.1), (3.4) and (3.5) form the set of master equations, in the sense that they determine the $`t`$ and $`y`$ dependence of the metric. We study the solutions of these equations in the next section.
## IV Exact solutions on the bulk
As a simple application of our master equations let us reconsider the cases already studied in the literature. Let us assume that the stress tensor of the bulk gets contribution only from a cosmological constant, $`\widehat{T}^A{}_{B}{}^{}=\mathrm{\Lambda }_B\delta _B^A`$. Then, Eq. (3.4) reduces into
$$3\frac{a^{\prime \prime }}{a}+\frac{n^{\prime \prime }}{n}=\frac{2}{3}\kappa _5^2\mathrm{\Lambda }_B.$$
(4.1)
This equation, together with the scaling equation (3.5), can be solved in a straightforward manner. For $`\mathrm{\Lambda }_B=0`$ we get a linear solution just as in
$$a(y,t)=A|y|+B;n=\lambda \left(\dot{A}|y|+\dot{B}\right).$$
(4.2)
Here we have already imposed $`P`$ (parity) symmetry on the solution. Next, by using the boundary conditions (2.14) and (2.15) we get the final result
$$a(y,t)=a_0\left(1\frac{\kappa _5^2}{6}\rho _b|y|\right);\text{and}n(y,t)=n_0\left(1+\frac{\kappa _5^2}{6}(3p_b+2\rho _b)|y|\right).$$
(4.3)
Notice that in the expression for $`n`$ requiring consistency with the scaling (3.5) leads to the conservation law
$$\dot{\rho }_b+3H_0(p_b+\rho _b)=0.$$
(4.4)
Note that $`a_R^{\prime \prime }=0`$ in this case, and there is no obvious way to get the correct Friedmann equation (3.2) (i.e. linear rather than squared dependence of $`H^2`$ on $`\rho _b`$).
Next, we assume that $`\mathrm{\Lambda }_B`$ is non zero. Again (4.1) is easy to solve, and after using the boundary conditions, we get for $`\mathrm{\Lambda }_B<0`$
$`a(y,t)`$ $`=`$ $`a_0\left(\mathrm{cosh}(\mu |y|){\displaystyle \frac{\kappa _5^2}{6\mu }}\rho _b\mathrm{sinh}(\mu |y|)\right),`$ (4.5)
$`n(y,t)`$ $`=`$ $`n_0\left(\mathrm{cosh}(\mu |y|)+{\displaystyle \frac{\kappa _5^2}{6\mu }}(3p_b+2\rho _b)\mathrm{sinh}(\mu |y|)\right);`$ (4.6)
where $`\mu ^2=\kappa _5^2\mathrm{\Lambda }_B/6`$. We recognize those solution presented by Binetruy et al. in when they take their integration constant as zero. Our finding is that this is the only allowed solution with a stable extra dimension. In this case, the condition (3.5) reduces to (4.4). For $`\mathrm{\Lambda }_B>0`$ the solutions have a similar form, with the hyperbolic functions replaced by $`\mathrm{cos}`$ and $`\mathrm{sin}`$ respectively.
If we also assume a brane tension, $`\mathrm{\Lambda }_b`$, the Friedmann equation on the brane becomes
$$H_0^2=\frac{\kappa _5^4}{18}\mathrm{\Lambda }_b\rho _b+\frac{\kappa _5^4}{36}\rho _b^2\frac{k}{a_0^2}+\frac{\kappa _5^4}{36}\mathrm{\Lambda }_b^2+\frac{\kappa _5^2}{6}\mathrm{\Lambda }_B.$$
(4.7)
We emphasize that this result shows that $`𝒞`$, the unknown constant found in is actually zero. A fine tunning among the cosmological constants
$$\mathrm{\Lambda }_B=\frac{\kappa _5^2}{6}\mathrm{\Lambda }_b^2$$
(4.8)
may linearize the Friedmann equation in the limit where $`\rho _b\mathrm{\Lambda }_b`$. As it is clear, this will work only if the vacuum energy of the bulk is negative. Moreover, to get the right expression of the Friedmann equation as in standard cosmology, an extra fine tunning
$$\mathrm{\Lambda }_b=6\frac{\kappa _4^2}{\kappa _5^4}$$
(4.9)
is required, where $`\kappa _4^2=M_P\mathrm{}^2`$. With those new assumptions the solution (4.6) reduces into
$`a=a_0\left(e^{\mu |y|}{\displaystyle \frac{\rho _b}{\mathrm{\Lambda }_b}}\mathrm{sinh}(\mu |y|)\right)`$ (4.10)
$`n=n_0\left(e^{\mu |y|}+{\displaystyle \frac{3p_b+2\rho _b}{\mathrm{\Lambda }_b}}\mathrm{sinh}(\mu |y|)\right)`$ (4.11)
Clearly, if we neglect the contribution of the brane densities, we identify the Randall-Sundrum static solutions
$$a=a_0e^{\mu |y|}\text{and}n=n_0e^{\mu |y|}.$$
(4.12)
So far we have only been concerned with finding solutions of the five dimensional cosmological equations with an empty bulk. If we wanted to include a bulk field, e.g. to provide a new picture of inflation, one would have a modified form of the Equation 3.4:
$$3\frac{a^{\prime \prime }}{a}+\frac{n^{\prime \prime }}{n}=2\frac{\kappa _5^2}{3}\left(V(\varphi )+\mathrm{\Lambda }_B\right).$$
(4.13)
Clearly, if the scalar field evolves very slowly, e.g in the inflation epoch, the solution will be approximately of the form (4.6), but now with
$$\mu ^2=\frac{\kappa _5^2}{6}\left(V(\varphi )+\mathrm{\Lambda }_B\right).$$
(4.14)
By setting the fine tunning (4.8) on (3.2) we get
$$H_0^2=\frac{\kappa _5^2}{6}V(\varphi (0))\frac{k}{a_0^2}+\frac{\kappa _5^4}{18}\mathrm{\Lambda }_b\rho _b+\frac{\kappa _5^4}{36}\rho _b^2.$$
(4.15)
As expected, the inflaton potential contributes linearly to last equation, and the effective potential is just the value of the bulk potential on the brane.
## V Non factorization of the metric parameters
Before closing the present discussion, let us comment briefly on the nature of the exact solutions. First we point out that all the solutions presented on the previous section are not factorizable. This is in fact a general property of this class of models as we will show now. For this purpose, let us start with case when the bulk radius $`b`$ is time dependent and we will see that factorization of the metric parameters $`n(y,t)`$ and $`a(y,t)`$ will not be consistent with a stable $`b`$. For this purpose, we start with the identity
$$\frac{d}{dt}\left(\frac{a^{}}{a}\right)=\frac{\dot{a}^{}}{a}\frac{\dot{a}}{a}\frac{a^{}}{a};$$
(5.1)
which, together with the $`G_{04}`$ equation leads to
$$\frac{d}{dt}\left(\frac{a^{}}{a}\right)=H\left(\frac{n^{}}{n}\frac{a^{}}{a}\right)+\frac{\dot{b}}{b}\frac{a^{}}{a}.$$
(5.2)
Note that this expression reduces to Eq. (2.16) on the brane.
Now, let us assume that $`a`$ and $`n`$ are factorizable i.e.
$$a=a_0(t)\beta _1(y)n=n_0(t)\beta _2(y).$$
(5.3)
One can then easily see that
$$\frac{d}{dt}\left(\frac{a^{}}{a}\right)=0.$$
(5.4)
This can be used to rewrite Eq. (5.2) as
$$\frac{\dot{b}}{b}+H_0\left(\frac{\beta _1\beta _2^{}}{\beta _1^{}\beta _2}1\right)=0.$$
(5.5)
This equation can be integrated in a straightforward manner and leads to
$$b=a_0^{\beta _3(y)};$$
(5.6)
where $`\beta _3`$ is given in terms of $`\beta _{1,2}`$ by the expression between parenthesis in (5.5). From this, we conclude that the factorization ansatz works when the bulk radius, $`b`$ is time dependent since clearly $`a_0`$ grows with time. A stable $`b`$ would then mean that $`\beta _3`$ must equal zero, which is possible only if $`\beta _1=\beta _2`$. Furthermore, since Eq. (5.4) holds everywhere, we may evaluate it on the brane and then conclude using Eq. (2.12) that $`\dot{\rho }_b=0`$. From this, it follows that factorizability of $`n(y,t)`$ and $`a(y,t)`$ proposed above implies that only a cosmological constant may be present in the brane. Since we require a true time dependent density in the brane to describe realistic cosmology e.g. the transition from a radiation to matter dominated universe, the exact solution to the metric can not be factorizable.
Parenthetically, let us note that if we had chosen a more general form for the exact solutions
$$a=a_0(t)\beta _1(y,t)n=n_0(t)\beta _2(y,t);$$
(5.7)
where the functions $`\beta _{1,2}`$ satisfy the boundary condition $`\beta _1(0,t)=\beta _2(0,t)=1`$, no such restriction on brane energy density would emerge. This is indeed the form of the solutions discussed above. For this case, the five dimensional scalar curvature can be written as
$$R=\beta _2^2R_{(4)}+\mathrm{};$$
(5.8)
where $`R_{(4)}`$ is the four dimensional scalar curvature formed by $`a_0`$ and $`n_0`$, the dots represent extra terms.
At this point, we note a puzzling feature with regard to the true definition of the Newton’s constant. One can define the Newton’s constant in two ways: one by identifying the coefficient in front of the energy density in Friedmann equation and another way is by integrating over the fifth dimension in the action integral. We call the first definition a “local” one whereas the second one we call “global”. Clearly in the presence of matter, the first one gives a constant $`G_N`$ whereas the second gives a time-dependent $`G_N`$. We can check this easily as follows. Consider the gravity action in five dimensions
$$S=d^4x𝑑y\sqrt{g}\frac{1}{2\kappa _5^2}R,$$
(5.9)
Using the relation
$$\sqrt{g}=\sqrt{g_4}\beta _1^3\beta _2$$
(5.10)
(with $`g_4`$ the determinant of the corresponding four dimensional metric) and Eq. (5.8), we obtain
$$\frac{1}{\kappa _4^2}=\frac{1}{\kappa _5^2}𝑑y\beta _1^3\beta _2^1$$
(5.11)
adopting the “global” definition of Newton’s constant. The integral in the last equation is taken over the whole fifth dimension. It is clear from the above expression, that this leads, in general to a time dependent Newton’s constant whereas if we used Friedmann equation, we would get it to be time independent. We wish to note that when one uses the static solution, a lá Randall and Sundrum, there is no time dependence due to the absence of matter and hence no puzzle.
## VI conclusions
In summary, we have analyzed Einstein equations for cosmology in five dimensions within a brane-bulk picture. Restricting to the case of a stable bulk radius, we extract the generalized Friedmann equation which does not contain any integration constant, but a term that involves a regular part of the second derivative over the spatial component of the metric along the brane coordinates. We then evaluate this second derivative using the complete set of the master equations for five dimensional cosmology and find that for the case of a stable bulk radius, there is no arbitrary constant in the Friedmann equation in the brane. This makes it easier to interpret the brane cosmology as the standard big bang picture. An advantage of our analysis is that we do not need to make any explicit assumptions regarding the bulk content. In this sense it generalizes the results presented before in the literature where only a cosmological constant was assumed to be present in the bulk to obtain the Friedmann equation. Finally we make some remarks on the general nonfactorizable nature of the exact solutions.
Acknowledgements. The work of RNM is supported by a grant from the National Science Foundation under grant number PHY-9802551. The work of APL is supported in part by CONACyT (México). The work of CP is supported by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP). We wish to thank C. Van de Bruck and S. Pastor for discussions. |
warning/0003/astro-ph0003223.html | ar5iv | text | # The Tip of the Red Giant Branch and Distance of the Magellanic Clouds: results from the DENIS survey
## 1 Introduction
In the evolution of stars the position of the tip of the red giant branch (TRGB) marks the starting point of helium burning in the core. It is one of the strongest characteristics of the life of stars seen in theoretical models, together with the main sequence turn–off point, the red giant and the asymptotic giant clump. It has been used successfully for several decades (Sandage san (1971)) to estimate the distance of resolved galaxies (e.g., Lee, Freedman & Madore lfm (1993)). The TRGB magnitude depends only very weakly on age and metallicity, and yields comparable precision as classical distance indicators such as Cepheids and RR–Lyra variables.
Cioni et al. (2000a ) prepared the DENIS Catalogue towards the Magellanic Clouds (DCMC), as part of the Deep Near Infrared Southern Sky Survey performed with the 1m ESO telescope (Epchtein et al. epal (1997)). The catalogue contains about $`\mathrm{1\hspace{0.17em}300\hspace{0.17em}000}`$ and $`\mathrm{300\hspace{0.17em}000}`$ sources toward the LMC and the SMC, respectively; $`70`$% of them are real members of the Clouds and consist mainly of red giant branch (RGB) stars and asymptotic giant branch (AGB) stars, and $`30`$% are galactic foreground objects. This is a very large and homogeneous statistical sample that allows a highly accurate determination of the TRGB magnitude at the corresponding wavelengths. Among other things, this yields an important new determination of the distance modulus of the LMC. This distance modulus is one of the main stepping stones in the cosmological distance ladder, yet has remained somewhat uncertain and controversial (e.g., Mould et al. mould (2000)).
Section 2 describes how the data were selected from the DCMC catalogue to avoid crowding effects, and how we have calculated bolometric corrections. Section 3 discusses the luminosity function (LF) and the subtraction of the foreground component. Section 4 discusses the TRGB determination and gives comparisons with previous measurements. Section 5 discusses the implications for the distances to the Magellanic Clouds. Concluding remarks are given in Section 6. The Appendix provides a detailed description of the new method that we have used to quantify the TRGB magnitude, as well a discussion of the formal and systematic errors in the analysis.
## 2 The Sample
### 2.1 The Data
The DCMC covers a surface area of $`19.87\times 16`$ square degrees centered on $`(\alpha ,\delta )=(5^h27^m20^s`$,$`69\mathrm{°}00\mathrm{}00\mathrm{})`$ toward the LMC and $`14.7\times 10`$ square degrees centered on $`(\alpha ,\delta )=(1^h02^m40^s`$,$`73\mathrm{°}00\mathrm{}00\mathrm{})`$ toward the SMC (J2000 coordinates). We extracted all the sources detected simultaneously in the three DENIS photometric wave bands: $`I`$ ($`0.8\mu `$m), $`J`$ ($`1.25\mu `$m) and $`K_S`$ ($`2.15\mu `$m). We excluded sources that were detected in all three wave bands but at different times (this can happen because DENIS strips overlap). The selection of sources that are present in all three wave bands strongly reduces possible crowding effects that affect mostly the $`I`$ band. We removed sources affected, even slightly, by image defects (null image flag) and sources with bright neighbours or bad pixels, sources that were originally blended, or sources with at least one saturated pixel (null extraction flag). This increases the level of confidence on the resulting sample. The main final sample for the present analysis contains $`\mathrm{33\hspace{0.17em}117}`$ sources toward the SMC and $`\mathrm{118\hspace{0.17em}234}`$ sources toward the LMC. This constitutes about $`10`$% of all the sources listed for each Cloud in the DCMC.
To estimate the contribution of the foreground component we also considered the data in offset fields outside the spatial limits of the DCMC<sup>1</sup><sup>1</sup>1These data are not part of the DCMC catalogue but are available on request from the first author., covering the same range in right ascension and from a maximum of $`\delta =57^{}`$ to a minimum of $`\delta =87^{}`$ (the full declination range of a DENIS strip). These data were reduced together and the same selection criteria, on the basis of the detection wave bands and the flags, were applied as to the data constituting the DCMC. The total sample (DCMC plus extension in declination) contains $`\mathrm{92\hspace{0.17em}162}`$ and $`\mathrm{184\hspace{0.17em}129}`$ sources in the RA ranges for the SMC and the LMC, respectively.
The distribution of the formal photometric errors in each wave band is shown in Fig. 1. At the brighter magnitudes (those of interest for the TRGB determination), the random errors in the sample are not dominated by the formal photometric errors, but by random errors in the photometric zero–points for the individual strips. The dispersions ($`1\sigma `$) of these zero-point variations are $`0.07`$ mag in the $`I`$ band, $`0.13`$ mag in the $`J`$ band and $`0.16`$ mag in the $`K_S`$ band. Note that the formal error with which the TRGB magnitude can be determined is not limited to the size of these zero-point variations, but instead can be quite small (the formal error is proportional to $`1/\sqrt{N}`$, where $`N`$ is the number of stars in the sample).
The $`I`$, $`J`$ and $`K_S`$ magnitudes in the present paper are all in the photometric system associated with the DENIS passbands. These magnitudes are not identical to the classical Cousins $`I`$ and CTIO $`J`$ and $`K`$ magnitudes, although they are close (differences are $`0.1`$ magnitudes). The final transformation equations for the passbands will not be available until the survey is completed, but a preliminary analysis is presented by Fouqué et al. (fou (1999)). Note that our determinations of the distance moduli for the LMC and the SMC (Section 5) are based on bolometric magnitudes derived from the data, which are fully corrected for the specifics of the DENIS passbands.
### 2.2 Bolometric correction
We have calculated the apparent bolometric magnitude ($`m_{bol}`$) for all the sources selected according to the criteria described in Section 2.1, and with $`(JK_S)0.4`$. We have chosen to use only the $`J`$ and $`K_S`$ bands to derive $`m_{bol}`$ (see below). Sources with $`(JK_S)<0.4`$ do not influence the position of the TRGB (see Fig. 5 below), and have too low a percentage of flux in the near-infrared (NIR) to give a reliable measure of $`m_{bol}`$ with these criteria. We used two different bolometric corrections, depending on the $`(JK_S)`$ colour. For sources with $`(JK_S)<1.25`$, we simply use a blackbody fit on the $`(JK_S)`$ colour; such sources are mostly RGB or early AGB (E–AGB) stars in our sample. Sources with larger values of $`(JK_S)`$ are mostly thermally pulsing AGB (TP–AGB) stars, some of which are losing mass and are surrounded by a circumstellar envelope. For them we used the results of individual modelling of galactic carbon (C) stars by Groenewegen et al. (groen (1999)), combined with a series of models of increasing dust opacity where the central star has a spectral type $`M5`$ and the dust grains are composed of silicates (Groenewegen, private communication).
In both cases, blackbody fit and spectral models, our method to infer $`m_{bol}`$ is different from what is usually performed in the literature. We do not make any attempt to transform a magnitude, i.e. an integrated flux over the passband, into a flux density at a reference wavelength, in order to suppress one step which already makes an assumption on the spectral distribution of the source. We only use the integrated flux measured over the $`J`$ and $`K_S`$ DENIS passbands. Theoretical spectral distributions, i.e. blackbodies with temperatures ranging from $`10,000`$ to $`300`$ K and the models from Groenewegen and collaborators, were multiplied with the DENIS passbands (which includes a mean atmosphere at la Silla observatory) to derive the percentage of the total flux which is measured in each DENIS passband as a function of the DENIS colours. Then, for each selected DCMC source, $`m_{bol}`$ is calculated by interpolating in the theoretical grids the percentage of flux measured in the $`J`$ and $`K_S`$ bands from the observed $`(JK_S)`$ colour. We have used here the same zero point as in Montegriffo et al. (monte (1998)). More details are provided in Loup et al. (lou (2000)).
We have compared our results with the bolometric corrections $`BC_K`$ inferred by Montegriffo et al. As can be seen in their Fig. 3, their bolometric correction is valid only for sources with $`0.2<(JK_S)<0.7`$, with a typical spread around the fit of $`0.1`$ magnitude. For sources with $`0.4<(JK_S)<0.5`$ our blackbody fit agrees with their bolometric corrections to within the errors. On the other hand, for some sources with $`(JK_S)>0.5`$, they underestimate $`m_{bol}`$ by $`0.5`$ to $`2`$ magnitudes compared to our calculations. This is not surprising and can be inferred already from their Fig. 3; it does not indicate a shortcoming in our approach. We also compared our results with what one obtains by making blackbody fits using both the ($`IJ`$) and ($`JK_S`$) colours. For sources with $`0.4<(JK_S)<1.25`$ it does not produce any systematic effect; there is merely a spread of typically $`0.1`$ magnitude between both calculations, consistent with the formal errors. Inclusion of the $`I`$ band would produce a systematic effect for bluer sources than those selected here, but those are not relevant for the TRGB determination. We therefore decided to use only the $`J`$ and $`K_S`$ band data in our calculations of $`m_{\mathrm{bol}}`$, to minimize the effects of the interstellar reddening which are much more pronounced in the $`I`$ band than in $`J`$ and $`K_S`$.
There are both random and systematic errors in our estimates of $`m_{\mathrm{bol}}`$. The random errors come from two sources, namely from the observational uncertainties in the observed $`J`$ and $`K_S`$ band magnitudes, and from the corresponding uncertainties in the $`(JK_S)`$ color. We have calculated the resulting random errors in the $`m_{\mathrm{bol}}`$ estimates through propagation of these errors. There are also two sources of systematic error in the $`m_{\mathrm{bol}}`$ estimates. The first one derives from uncertainties in the dust extinction correction. Our treatment of dust extinction is discussed in Section 2.3; Appendix A.3.4 discusses how the uncertainties in this correction introduce a small systematic error on the TRGB magnitude determination. The second source of systematic error comes from the difference between the real spectral energy distribution of the star and the one we assume to estimate $`m_{bol}`$. For blackbody fits, we did not make any attempt to estimate this error because we lack information for that purpose (we would need spectra and/or UBVRIJHKL photometry on a sample of stars). For the AGB star models from Groenewegen and collaborators, we can estimate part of this error. The $`(JK_S)`$ colour does not provide enough information to fully constrain the set of model parameters, i.e. $`(JK_S)`$ does not give a unique solution, especially when the chemical type of the star is unknown. With the models available in this work, we have estimated this systematic model error to be 5% on the interpolated percentage of flux. This is of course a lower limit as there can be some objects whose spectral energy distribution differs from all the ones produced in the models. On the other hand, for most stars near the TRGB the blackbody fit is the relevant model, and for these the systematic errors could be smaller. In the end we have included in our final error budget a systematic error of $`\pm 0.05`$ mag in our $`m_{\mathrm{bol}}`$ estimates due to uncertainties in the underlying spectral model, but it should be noted that this estimate is not very rigorous.
In our analysis of the TRGB magnitude we have propagated the random and systematic errors on $`m_{bol}`$ separately. However, for illustrative purposes we show in Fig. 1d the combined error. The surprising shape of the error on $`mbol`$ as a function of $`m_{\mathrm{bol}}`$ should not be taken as real. It is an artifact coming from the fact that a systematic model error was included in the figure only for TP-AGB stars. The great majority of the brightest stars are TP-AGB stars for which we use AGB models. Going towards fainter stars, the $`(JK_S)`$ colour decreases and we mostly use blackbody fits, for which we have not included a systematic model error in the figure. The error on $`m_{bol}`$ thus seems to decrease around the TRGB.
### 2.3 Dust Extinction
The contribution of the internal reddening for the Magellanic Clouds is on average only $`E(BV)=0.06`$ while the foreground reddening can be very high in the outskirts of the Clouds. We have not attempted to correct our sample for extinction on a star by star basis. Instead we correct all data for one overall extinction. We adopt $`E(BV)=0.15\pm 0.05`$ as the average of known measurements (Westerlund west (1997)) for both Clouds. Adopting the extinction law by Glass (glass (1999)) for the DENIS pass bands \[$`A_V`$ : $`A_I`$ : $`A_J`$ : $`A_{K_S}`$ = $`1`$ : $`0.592`$ : $`0.256`$ : $`0.089`$\] and $`R_v=3.1`$ we obtain $`A_I=0.27`$, $`A_J=0.11`$ and $`A_{K_S}=0.04`$. Our approach to correct for dust extinction is a simple approximation to what is in reality a very complicated issue (e.g., Zaritsky zar (1999)). We discuss the effect of uncertainties in the dust extinction on our results in Sections 4.4 and 5. While this is an important issue in the $`I`$ band, the bolometric magnitudes that we use to determine the distance modulus are impacted only at a very low level.
## 3 The Luminosity Function
The luminosity function (LF) of a stellar population is a powerful tool to probe evolutionary events and their time scales. Major characteristics of a stellar population are associated to bumps, discontinuities and slope variations in the differential star counts as a function of magnitude. However, for a proper interpretation of observed luminosity functions several important issues should be taken into account. These include the completeness of the sample of data, the foreground contamination with respect to the analyzed population, the photometric accuracy and the size of the sampling bins. The total number of objects involved plays an important role to make the statistics significant.
In most previous studies of the luminosity functions of stellar populations in clusters or galaxies, in either the optical or the NIR, limited statistics have been the main problem. The DENIS (Cioni et al. 2000a ) and 2MASS (Nikolaev & Weinberg Nik (2000)) samples provide the first truly large statistical sample in the NIR of the Magellanic Cloud system. This wavelength domain is the most suitable to study late evolutionary stages such as the RGB and the AGB. In the present paper we restrict the discussion of the luminosity function mostly to the TRGB; a more general discussion is given elsewhere (Cioni, Messineo & Habing 2000).
### 3.1 The contribution of the Galaxy
For the removal of foreground contamination we considered two offset fields around each cloud. The range of right ascension (RA) is the same for both the cloud and the offset fields; it is the same of the DCMC catalogue (Section 2.1). For the LMC the north field has $`58^{}>\delta >60^{}`$ and the south field has $`80^{}>\delta >86^{}`$; for the SMC the north field has $`60^{}>\delta >66^{}`$ and the south field has $`80^{}>\delta >86^{}`$. The LMC region itself was limited to the declination range $`62^{}>\delta >76^{}`$, and the SMC region to $`69^{}>\delta >77^{}`$. Fig. 2 shows the distribution versus declination of the sources in the sample, using bins of $`0.1`$ degrees. The foreground contribution clearly decreases toward more negative declinations, due to the difference in Galactic latitude. The difference in number between the foreground contribution around the LMC and around the SMC is consistent with the fact that the LMC is observed closer to the galactic plane than the SMC is. The structure of the LMC is clearly wider than the one of the SMC and this may contribute to create the strong declination trend around the LMC.
For each field and photometric band we constructed a histogram of the observed magnitudes (thin solid curves in the $`N(m)`$ panels of Fig. 3). For the two different offset fields at each right ascension range the data were combined into one histogram. This offset–field histogram (thin dashed curves) was then scaled to fit the corresponding LMC or SMC field histogram at bright magnitudes, for which almost all the stars belong to the foreground. Subtraction yields the foreground–subtracted magnitude distribution for each of the Clouds (heavy solid curves). For comparison we also extracted from the catalogue an extended sample consisting of those stars detected in the $`I`$ and $`J`$ bands (irrespective of whether or not they were detected in $`K_S`$). This sample (heavy dashed curves) is complete to fainter magnitudes than the main sample, and therefore illustrates the completeness limit of the main sample.
### 3.2 The shape
The resulting statistics of the subtracted LF are impressive, despite the restricted source selection. We proceed with a description of the major characteristics of the LF. The maximum corresponds to giants that lie on the upper part of the RGB. The decrease at fainter magnitudes is due to the selections applied to the data and to the decrease in sensitivity of the observations (Cioni et al. 2000a ). Features like the horizontal branch or the red clump are too faint to be detected by DENIS. Towards brighter magnitudes we encounter a strong kink in the profile, which we associate with the position of the TRGB discontinuity. Brightward of the kink follows a bump of objects which we discuss below. At very bright magnitudes the LF has a weak tail which is composed of stars of luminosity type I and II (Frogel & Blanco fb (1983)), but the LF at these bright magnitudes could be influenced by small residuals due to inaccurate foreground subtraction.
To explain the bump brightward of the TRGB discontinuity we cross–identified (Loup loup (2000)) the DCMC sources with the sources in some of the Blanco fields in the LMC (Blanco bbmc (1980)). In the $`(K_S,JK_S)`$ diagram there are two regions populated only by oxygen rich AGB stars (O–rich) and by carbon rich AGB stars (C–rich), respectively. O–rich stars are concentrated around $`K_S=11.5`$ and have a constant color $`(JK_S)=1.2`$, and C–rich stars are concentrated around $`(JK_S)=1.7`$ and around $`K_S=10.5`$ (see Fig. 5b). These TP–AGB stars cause the bump visible in the LF. This bump should not be confused with the AGB bump caused by E–AGB stars (Gallart gal (1998)). Fig. 4 shows an enlargement of Figs. 3c and 3k (continuous line). The dashed line refers to O–rich AGB stars and the dotted line to C–rich AGB stars selected in the $`(K_S,JK_S)`$ diagram. In the case of the SMC we selected regions with slightly bluer color and fainter magnitude to match the two groups of AGB stars in the $`(K_S,JK_S)`$ diagram, cf. Fig. 5d. Fig. 4 also plots the LF (thick line) that results when we cross–identify our sample with the spectroscopically confirmed carbon stars by Rebeirot et al. (reb (1983)) in the SMC. We found $`1451`$ sources out of $`1707`$ and we attribute the missing cross–identifications to the selection criteria that we applied to the DCMC data to obtain the sample for the present paper. It is interesting to note that at higher luminosities the distribution of the confirmed C–rich stars matches the distribution of C–rich stars selected only on the basis of $`K_S`$ and $`(JK_S)`$. At the fainter luminosities C–rich AGB stars cannot be discriminated from O–rich AGB stars only on the basis of $`(JK_S)`$ and $`K_S`$ because they overlap with the RGB, principally constituted by O–rich stars.
## 4 The tip of the RGB
### 4.1 Theory
Theoretically stars climb the RGB with an expanding convective envelope and an hydrogen burning shell, while increasing the core–Helium content, the central temperature, the central density, and the luminosity. Low–mass stars ($`0.81.0M_{}<M<22.3M_{}`$) develop an electron–degenerate core, which causes an explosive start (Helium–flash) of the core–Helium burning when the core mass reaches $`0.45M_{}`$, almost independently from the initial mass and composition of the star (Chiosi et al. cbb (1992)); intermediate mass stars ($`22.3M_{}<M<89M_{}`$) are not affected by degeneracy at this stage and initiate helium burning quietly, when a suitable temperature and density are reached. The RGB transition phase between the two behaviors occurs when the population is at least $`0.6`$ Gyr old and lasts roughly for $`0.2`$ Gyr, determining an abrupt event in the population life time (Sweigart et al. sgr (1990)). The Helium–flash is followed by a sudden decrease in the luminosity because of the expansion of the central region of the star and because of the extinction of the hydrogen–burning shell, the major nuclear energy supply. The star reaches its maximum luminosity and radius (in the RGB phase) at the TRGB, which also marks the end of the phase itself (Iben iben (1967)). Low–mass stars with the same metallicity accumulate along the RGB up to a TRGB luminosity of about $`2500L_{}`$ (Westerlund west (1997)); the resulting RGB is quite extended. Stars with masses just above the transition mass (which discriminates between low and intermediate masses) have a TRGB luminosity as low as $`200L_{}`$ (Sweigart et al. sgro (1989), sgr (1990)) and the RGB is almost non–existent. Both low and intermediate mass stars that finish burning their Helium in the core evolve on the AGB phase. They are in the so called E–AGB when Helium is burning in a thick shell and in the so called TP–AGB when both the Hydrogen and the Helium shells are active. The luminosity increases because of the increase in mass of the degenerate carbon core. The AGB evolution is characterized by a strong mass loss process that ends the phase when the outer envelope is completely lost. The maximum AGB luminosity defines the tip of the AGB (TAGB), with core mass $`M_{\mathrm{core}}=1.4M_{}`$ and magnitude $`M_{bol}=7.1`$ mag (Paczynski pac (1970)).
### 4.2 Detections
In the observed diagrams $`(I,IJ)`$ and $`(K_S,JK_S)`$ the RGB is clearly visible (Fig. 5). The beginning of the RGB phase is below the detection limits and the spread at the fainter magnitudes is due to the photometric errors. The TRGB is clearly defined at the brightest point of this branch as an outstanding roughly horizontal feature. Dashed horizontal lines in the figure indicate the values of the TRGB discontinuity that we derive below for these data. The plume of objects brighter than the TRGB is composed of AGB stars experiencing the TP phase. From these diagrams the foreground contribution has not been subtracted but the contamination of these to the RGB/AGB is negligible (Cioni et al. chl (1998), 2000a ) if only the very central region of each cloud is selected; Fig. 5 contains sources with $`67^{}>\delta >69^{}`$ toward the LMC and $`72^{}>\delta >74^{}`$ toward the SMC. Stars populating the RGB up to the TRGB are low–mass stars older than $`0.6`$ Gyr. TP–AGB stars on the other hand, which lie above the TRGB, can be either low–mass stars or intermediate mass–stars. For $`M_{bol}<6`$ mag they all originate from main–sequence stars with $`M<3M_{}`$ (Westerlund west (1997)), which corresponds to a minimum age of $`0.2`$ Gyr. TP–AGB stars that are low–mass stars should be older than $`1`$ Gyr (Vassiliadis and Wood vw (1993)). Note that the thickness of the RGB ($`0.3`$ mag) is larger than the photometric errors involved ($`0.1`$ mag) and this indicates a spread in either metallicity or extinction within each cloud.
### 4.3 Method
The algorithm that we have used for the determination of the position of the magnitude $`m_{\mathrm{TRGB}}`$ of the TRGB is described in great detail in Appendix A. The TRGB discontinuity causes a peak in both the first derivative $`N^{}(m)\mathrm{d}N(m)/\mathrm{d}m`$ and the second derivative $`N^{\prime \prime }(m)\mathrm{d}^2N(m)/\mathrm{d}m^2`$ of the observed stellar magnitude distribution $`N(m)`$. Previous authors have generally used $`N^{}(m)`$ to estimate $`m_{\mathrm{TRGB}}`$ (e.g., Madore & Freedman mado (1995)). Based on extensive tests and simulations we found that for our dataset $`N^{\prime \prime }(m)`$ provides a better handle on $`m_{\mathrm{TRGB}}`$ (cf. Appendix A.1). We therefore adopted the following approach. First, we use a Savitzky-Golay filter (e.g., Press et al. press (1992)) to estimate $`N^{\prime \prime }(m)`$. We then search for a peak in $`N^{\prime \prime }(m)`$, and fit a Gaussian to it to obtain the quantities $`m_{2g}`$ and $`\sigma _{2g}`$ that are the mean and dispersion of the best-fitting Gaussian, respectively. The magnitude $`m_{\mathrm{TRGB}}`$ is then estimated as $`m_{2g}+\mathrm{\Delta }m_{2g}(\sigma _{2g})`$, where $`\mathrm{\Delta }m_{2g}(\sigma _{2g})`$ is a small correction (Fig. 8b) derived from a phenomenological model described in Section A.1. The formal errors on the $`m_{\mathrm{TRGB}}`$ determinations are inferred from extensive Monte-Carlo simulations, as described in Section A.2. The possible influence of systematic errors is discussed in Section A.3. There is no evidence for any possible systematic errors due to possible incompleteness in the sample, or inaccuracies in the foreground subtraction. Systematic errors due to uncertainties in the phenomenological model on which the corrections $`\mathrm{\Delta }m_{2g}(\sigma _{2g})`$ are based can be up to $`\pm 0.02`$ magnitudes. Extinction variations within the Clouds do not cause systematic errors in either the estimate of $`m_{\mathrm{TRGB}}`$ or its formal error. However, any error in the assumed average extinction for the sample does obviously translate directly into an error in $`m_{\mathrm{TRGB}}`$.
Fig. 3 summarizes the results of the analysis. The second and fourth row of the panels show the estimates of $`N^{\prime \prime }(m)`$. The Gaussian fit to the peak is overplotted, and its center $`m_{2g}`$ is indicated by a vertical dotted line. The corresponding estimate $`m_{\mathrm{TRGB}}`$ is indicated by a vertical dotted line in the panel for $`N(m)`$. Table 1 lists the results. It includes both the observed value for $`m_{\mathrm{TRGB}}`$, as well as the value obtained after correction for extinction with $`E(BV)=0.15`$. Formal errors are listed as well, and are typically $`0.03`$$`0.04`$ magnitudes. The last column of the Table lists the amount by which the extinction-corrected $`m_{\mathrm{TRGB}}`$ would change if the assumed $`E(BV)`$ were increased by $`+0.05`$ (a shift of $`0.05`$ in the assumed $`E(BV)`$ would produce the opposite shift in $`m_{\mathrm{TRGB}}`$).
When applying comparable methods to resolvable galaxies in the Local Group (e.g., Soria et al soria (1996); Sakai et al. sakai (1996)) one of the major sources of contamination on the TRGB determination is the presence of a relative strong AGB population. The Magellanic Clouds also have a strong AGB population, but in our case this does not confuse the determination of $`m_{\mathrm{TRGB}}`$. This is due to the large statistics available, and above all to the fact that TP–AGB stars are definitely more luminous than the TRGB. E–AGB stars overlap with the RGB stars but there is no reason to assume, according to models, that they accumulate at the TRGB. Probably they distribute rather constantly and due to the very short evolutionary time scale we do not expect them to exceed more than $`10`$% of the RGB population.
### 4.4 Discussion
The absolute magnitude of the TRGB generally depends on the metallicity and the age of the stellar population and therefore need not to be the same for the LMC and the SMC. Nonetheless, if we assume that such differences in TRGB absolute magnitude are small or negligible, and if we assume that the extinction towards the LMC and the SMC have been correctly estimated, then one may subtract for each photometric band the inferred $`m_{\mathrm{TRGB}}`$(LMC) from the inferred $`m_{\mathrm{TRGB}}`$(SMC) to obtain an estimate of the difference $`\mathrm{\Delta }(mM)_{\mathrm{SMC}}(mM)_{\mathrm{LMC}}`$ between the distance moduli of the SMC and the LMC. This yields the following results: $`0.41\pm 0.04`$ ($`I`$ band), $`0.56\pm 0.04`$ ($`J`$ band), $`0.64\pm 0.08`$ ($`K_S`$ band) and $`0.46\pm 0.05`$ ($`m_{\mathrm{bol}}`$). The dispersion among these four numbers is $`0.09`$, which is somewhat larger than the formal errors. Averaging the four determinations yields $`\mathrm{\Delta }=0.52\pm 0.04`$, where the error is the formal error in the mean. This is not inconsistent with determinations found in the literature, which generally fall in the range $`\mathrm{\Delta }=0.4`$$`0.5`$ (Westerlund west (1997)).
Upon taking a closer look at the values of $`\mathrm{\Delta }`$ for the different bands one sees that the values in $`J`$ and $`K_S`$ exceed those in $`I`$ by $`0.15`$ mag or more. It is quite possible that this is due to differences in the metallicity and age of the LMC and the SMC, which affect the TRGB absolute magnitude $`M_{\mathrm{TRGB}}`$ differently in different bands. In the $`I`$ band $`M_{\mathrm{TRGB}}`$ is reasonably insensitive to metallicity and age. Lee et al. (lfm (1993)) showed that $`M_{\mathrm{TRGB}}(I)`$ changes by less than $`0.1`$ mag for $`2.2<[Fe/H]<0.7`$ dex and for ages between $`2`$ and $`17`$ Gyr. For the $`K`$ band, Ferraro et al. (ferr (1999)) derived an empirical relation between $`M_{\mathrm{TRGB}}(K)`$ and the metallicity in galactic globular clusters. For metallicities in the range of the Magellanic Clouds the variation of $`M_{\mathrm{TRGB}}(K)`$ is about $`0.2`$ mag; however, this relation might not be valid for intermediate age populations. From the theoretical isochrones by Girardi et al. (gir (2000)) the spread of $`M_{\mathrm{TRGB}}(K)`$ is about $`0.3`$ mag for ages greater than $`2`$ Gyr and constant metallicity. This spread is somewhat less for the $`J`$ band but it remains higher than the one derived for the $`I`$ band. The fact that $`M_{\mathrm{TRGB}}`$ is modestly sensitive to variations in metallicity and age for the $`J`$ and $`K`$ bands implies that the values of $`\mathrm{\Delta }`$ derived in these bands may not be unbiased estimates of the true difference in distance modulus between the SMC and the LMC. The $`I`$ band value should be better in this respect, but on the other hand, that value is more sensitive to possible differences in the dust extinction between the Clouds. So the best estimate of $`\mathrm{\Delta }`$ is probably obtained using $`m_{\mathrm{bol}}`$, as discussed further in Section 5.
For the LMC there are several observed TRGB magnitude determinations in the literature that can be compared to our results. Reid et al. (rmt (1987)) obtained $`m_{\mathrm{TRGB}}(I)=14.53\pm 0.05`$, after extinction-correction with an assumed $`A_I=0.07`$. Romaniello et al. (rscp (1999)) obtained $`m_{\mathrm{TRGB}}(I)=14.50\pm 0.25`$ for the field around SN1987A. They corrected each star individually for extinction, but found a mode of $`E(BV)=0.20`$ for their sample (corresponding to $`A_I=0.30`$). Sakai et al. (szk (1999)) obtained $`m_{\mathrm{TRGB}}(I)=14.54\pm 0.04`$. They also corrected each star individually for extinction, but restricted their sample to low-extinction regions with $`A_V<0.2`$ (corresponding to $`A_I<0.10`$). The observed value of $`m_{\mathrm{TRGB}}(I)`$ for our sample, $`14.54\pm 0.03`$, is nicely consistent with all these determinations. However, when we apply an extinction correction of $`A_I=0.27`$, as appropriate for an assumed $`E(BV)=0.15`$ (Section 2.3), our corrected value falls significantly below the previous determinations. This may mean that our assumed extinction is an overestimate. Support from this comes from a recent study by Zaritsky (zar (1999)). He demonstrates that the average extinction towards cool stars is much lower than for the hotter stars which have typically been used to estimate the extinction towards the LMC (the latter generally reside in star-forming regions which are more dusty, among other things). The analysis of Zaritsky (cf. his Fig. 12) suggests that the mode of the distribution of $`A_V`$ for stars with temperatures appropriate for the RGB is as low as $`A_V0.1`$ (corresponding to $`A_I0.05`$), but with a long tail towards higher extinctions. Either way, it is clear that any proper interpretation of the TRGB magnitude in the $`I`$ band requires an accurate understanding of the effects of dust extinction. We have not (yet) performed such an extinction analysis for our sample, and therefore refrain from drawing conclusions from our $`I`$ band results. However, our results are not inconsistent with observations by previous authors, provided that the extinction is actually as low as suggested by Zaritsky.
The best way to circumvent any dependence of the results on uncertainties in the dust extinction is to go far into the near IR. There is one very recent determination of $`m_{\mathrm{TRGB}}`$ in the $`K_S`$ band that can be compared to our results. Nikolaev & Weinberg (Nik (2000)) used data from the 2MASS survey to derive $`m_{\mathrm{TRGB}}(K_S)=12.3\pm 0.1`$ for the LMC, without correcting for extinction. For a proper comparison of this value to our results we must correct for possible differences in the photometric magnitude systems used by 2MASS and DENIS. Neither system is identical to the standard CTIO $`K`$ magnitude system, but both are quite close. Nikolaev & Weinberg quote that their $`K_S`$ magnitude system agrees with the standard $`K`$ to within $`0.05`$ mag. For the DENIS system the final transformation equations will not be available until the survey is completed, but the analysis of Fouqué et al. (fou (1999)) yields an absolute flux zero-point (in Jy) for the DENIS $`K_S`$ system that differs from the CTIO $`K`$-band by $`0.08`$ mag. Based on this, we do not expect the $`K_S`$ magnitudes of 2MASS and DENIS to differ by much more than $`0.1`$ magnitudes. To determine the actual difference, we compare in Fig. 6 our LMC $`K_S`$ histogram to that presented by Nikolaev & Weinberg (using identical binning). The 2MASS histogram was shifted horizontally to obtain the best agreement. From this we obtain $`K_S(`$DENIS$`)=K_S(`$2MASS$`)0.11\pm 0.02`$. With this photometric correction the histograms are in good agreement. The slight differences at $`K_S<11`$ magnitudes are probably due to differences in foreground subtraction. At faint magnitudes the DENIS data become incomplete at brighter magnitudes than the 2MASS data. However, tests discussed in Appendices A.3.2 and A.3.3 show that our determinations of $`m_{\mathrm{TRGB}}`$ are not influenced significantly either by possible incompleteness near the TRGB or by possible uncertainties in the foreground subtraction. Upon correction of the Nikolaev & Weinberg $`m_{\mathrm{TRGB}}`$ determination to the DENIS $`K_S`$ magnitude system one obtains $`m_{\mathrm{TRGB}}(K_S)=12.19\pm 0.1`$. Somewhat surprisingly, this exceeds our determination $`m_{\mathrm{TRGB}}(K_S)=11.98\pm 0.04`$ by as much as $`0.21`$ magnitudes. Given that the histograms themselves are in good agreement (Fig 6), we are forced to conclude that this must be due to differences in how $`m_{\mathrm{TRGB}}`$ is defined and determined. While we search for a peak in $`N^{\prime \prime }(m)`$ and then add a correction term that is based on a model, Nikolaev & Weinberg just determine the peak in the first derivative $`N^{}(m)`$. As discussed in Section A.1 (see Fig. 8) this generally yields on overestimate of the actual TRGB magnitude. Since Nikolaev & Weinberg do not describe their analysis technique in detail, it is difficult to estimate the size of this bias in their result. However, Monte-Carlo simulations that we discuss in Section A.4 indicate that it could be $`0.15\pm 0.06`$, which would explain the observed discrepancy. Note that the same effect may also affect some of the $`I`$ band comparisons listed above, although for those the influence of extinction probably plays the more significant role.
## 5 Distance to the Magellanic Clouds
To estimate the distance modulus of the Magellanic Clouds we can use the observed magnitude of the TRGB in either $`I`$, $`J`$, $`K_S`$ or $`m_{\mathrm{bol}}`$. As discussed in Section 4.4, $`I`$ has the disadvantage of being sensitive to uncertain extinction corrections, while $`J`$ and $`K_S`$ have the disadvantage of being sensitive to the assumed metallicity and age. The most accurate information on the distance is therefore provided by $`m_{bol}`$, which is not particularly sensitive to either dust extinction (cf. Table 1) or metallicity and age. To quantify the latter we use the stellar evolutionary model calculations of Salaris & Cassisi (saca (1998)). They quantified the dependence of $`M_{\mathrm{TRGB}}(\mathrm{bol})`$ on the total metallicity ($`[M/H]`$) of a population, and found that
$$M_{\mathrm{TRGB}}(\mathrm{bol})=3.9490.178[M/H]+0.008[M/H]^2,$$
(1)
valid for $`2.35<[M/H]<0.28`$ and for ages larger than a few Gyr.
We determined $`[M/H]`$ by qualitatively fitting isochrones (Girardi et al. gir (2000)) to the color–magnitude diagram $`(K_S,JK_S)`$. We obtain $`Z=0.004\pm 0.002`$ for the LMC, in agreement with the value derived by Nikolaev & Weinberg, and $`Z=0.003\pm 0.001`$ for the SMC. For $`Z_{}=0.02`$ this corresponds to $`[M/H]=0.70`$ and $`[M/H]=0.82`$ for the LMC and the SMC, respectively. This in turn yields $`M_{\mathrm{TRGB}}(\mathrm{bol})=3.82`$ for the LMC and $`M_{\mathrm{TRGB}}(\mathrm{bol})=3.80`$ for the SMC. When combined with the results in Table 1 we obtain for the LMC that $`(mM)=18.55\pm 0.04`$ (formal) $`\pm 0.08`$ (systematic), and for the SMC that $`(mM)=18.99\pm 0.03`$ (formal) $`\pm 0.08`$ (systematic). The corresponding distances are $`51`$ and $`63`$ kpc to the LMC and the SMC respectively.
The systematic errors that we quote in our results are the sum in quadrature of the following possible (identified) sources of error: (i) $`\pm 0.02`$ mag due to uncertainties in the phenomenological model on which the corrections $`\mathrm{\Delta }m_{2g}(\sigma _{2g})`$ are based (cf. Section A.3.1); (ii) $`\pm 0.03`$ mag to account for the fact that our assumed average dust extinction of $`E(BV)=0.15`$ could plausibly be in error by $`0.05`$ (cf. Table 1); (iii) $`\pm 0.04`$ mag, reflecting the uncertainties in $`M_{\mathrm{TRGB}}(\mathrm{bol})`$ due to uncertainties in $`[M/H]`$; (iv) $`\pm 0.04`$ mag, reflecting the uncertainty in $`M_{\mathrm{TRGB}}(\mathrm{bol})`$ at fixed $`[M/H]`$ suggested by comparison of the predictions of different stellar evolution models (Salaris & Cassisi saca (1998); their Fig. 1); (v) $`\pm 0.05`$ mag, being an estimate of the possible systematic error in our calculation of bolometric magnitudes due to uncertainties in the underlying spectral model (see Section 2.2).
There have been many previous determinations of the distance modulus of the LMC, and these have varied widely, from about $`18.0`$ to $`18.7`$. Based on a collection of many determinations, the HST Key Project Team adopted $`(mM)=18.50\pm 0.13`$ (Mould et al. mould (2000)). Our determination is in excellent agreement with this value, and actually has a smaller error. The TRGB method itself has been used previously by several other authors to study the distance modulus of the LMC, and our results are consistent with all of these. Reid, Mould & Thompson (rmt (1987)) were the first to apply this technique to the LMC (by studying the Shapley Constellation III using photographic plates), and obtained $`(mM)=18.42\pm 0.15`$. Romaniello et al. (rscp (1999)) obtained $`(mM)=18.69\pm 0.25`$ from a field around SN1987A in the LMC using HST/WFPC2 data. Sakai et al. (szk (1999)) obtained $`18.59\pm 0.09`$ from an area of $`4\times 2.7`$ square degrees (north of the LMC bar) studied as part of the Magellanic Cloud Photometric Survey (Zaritsky, Harris & Thompson zht (1997)) using the Las Campanas 1m telescope. Nikolaev & Weinberg (Nik (2000)) obtained $`(mM)=18.50\pm 0.12`$ from the subset of 2MASS data that covers the LMC. For the SMC we are not aware of (recent) TRGB distance modulus measurements, but our result is consistent with the value $`(mM)=18.90\pm 0.10`$ quoted by Westerlund (west (1997)) from a combination of measurements available in the literature from a variety of techniques.
## 6 Conclusions
We have determined the position of the TRGB for both Magellanic Clouds using the large statistical sample offered by the DCMC (Cioni et al. 2000a ). We have presented a new algorithm for the determination of the TRGB magnitude, which we describe in detail in the Appendix and test extensively using Monte-Carlo simulations. We note that any method that searches for a peak in the first derivative (used by most authors) or the second derivative (used by us) of the observed luminosity function does not yield an unbiased estimate for the actual magnitude of the TRGB discontinuity. We stress the importance of correcting for this bias, which is not generally done. Our analysis shows that when large enough statistics are available, contamination by AGB stars does not provide a significant limitation to the accuracy of the TRGB magnitude determination.
In our analysis we have adopted global values for the extinction of the Magellanic Clouds and we have derived the metallicity from an isochrone fit to the giant population to obtain a representative value for each cloud as a whole. In reality, extinction and metallicity are likely to vary within each cloud. Clearly, the production of a detailed extinction map together with precise measurements of the metallicity is a requirement for a detailed analysis of variations in structure between different locations within the Clouds, either on the plane of the sky or along the line of sight. However, such variations do not influence our distance determinations, which should be accurate in a globally averaged sense. Uncertainties in the average dust extinction or metallicity for each cloud are included in the systematic error budget of our final estimates.
We combine our apparent bolometric TRGB magnitude determinations with theoretical predictions to derive the distance modulus of the Clouds. We obtain $`(mM)=18.55\pm 0.04`$ (formal) $`\pm 0.08`$ (systematic) for the Large Magellanic Cloud (LMC), and $`(mM)=18.99\pm 0.03`$ (formal) $`\pm 0.08`$ (systematic) for the Small Magellanic Cloud (SMC). These results are consistent with many previous studies, including a recent compilation by Mould et al. (mould (2000)). However, only very few previous studies have yielded determinations of similar accuracy as those presented here. This re-confirms the TRGB method to be a high quality method for distance determination of resolved stellar populations, and stresses the power of large statistical samples in the NIR such as those provided by the DENIS survey.
## Appendix A Determination of the TRGB magnitude: methodology and error analysis
### A.1 The nature of the TRGB discontinuity
We wish to determine the magnitude $`m_{\mathrm{TRGB}}`$ of the TRGB discontinuity from an observed magnitude distribution $`f_{\mathrm{obs}}(m)`$. In general, the observed distribution will be the convolution of the intrinsic magnitude distribution of the stars, $`f_{\mathrm{int}}(m)`$, with some broadening function $`E(m)`$:
$$f_{\mathrm{obs}}(m)=_{\mathrm{}}^{\mathrm{}}f_{\mathrm{int}}(m^{})E(mm^{})dm^{}.$$
(2)
The function $`E(m)`$ characterizes the probability that a star with magnitude $`m_0`$ is observed to have magnitude $`m_{\mathrm{obs}}=m_0+m`$. The shape of $`E(m)`$ is generally determined by the properties of the observational errors, but other effects (such as differences in extinction or distance among the stars in the sample) can contribute as well.
To gain an understanding of the issues involved in the determination of $`m_{\mathrm{TRGB}}`$ we start by considering a simple model. We assume that $`E(m)`$ is a Gaussian of dispersion $`\sigma `$:
$$E(m)=\frac{1}{\sqrt{2\pi }\sigma }e^{\frac{(m/\sigma )^2}{2}}.$$
(3)
We approximate $`f_{\mathrm{int}}(m)`$ by expanding it into a first-order Taylor expansion near the position of the discontinuity, which yields
$$f_{\mathrm{int}}(m)=\{\begin{array}{cc}f_0+a_1(mm_{\mathrm{TRGB}}),\hfill & \text{if }m<m_{\mathrm{TRGB}}\text{ ;}\hfill \\ f_0+\mathrm{\Delta }f+a_2(mm_{\mathrm{TRGB}}),\hfill & \text{if }m>m_{\mathrm{TRGB}}\text{ .}\hfill \end{array}$$
(4)
The parameters $`a_1`$ and $`a_2`$ measure the slope of $`f_{\mathrm{int}}`$ for magnitudes that are brighter and fainter than $`m_{\mathrm{TRGB}}`$, respectively. At brighter magnitudes the sample is dominated by AGB stars, while at fainter magnitudes both AGB and RGB stars contribute. The parameter $`\mathrm{\Delta }f`$ measures the size of the discontinuity; the ratio $`\mathrm{\Delta }f/f_0`$ is an estimate of the ratio of the number of RGB to AGB stars at the magnitude of the RGB tip.
We fitted the model defined by Eqs. (2)–(4) to the observed (foreground-subtracted) $`J`$ band magnitude histogram for the LMC, which is shown as a connected heavy dashed curve in Fig. 7a. The heavy solid curve shows the model distribution $`f_{\mathrm{obs}}`$ that provides the best fit. The fit is acceptable. The parameters for this model are: $`f_0=0.091`$, $`\mathrm{\Delta }f=0.250`$ (both in units in which the normalization of $`f`$ is arbitrary), $`a_1=0.108`$, $`a_2=0.928`$, $`m_{\mathrm{TRGB}}=13.16`$ and $`\sigma =0.126`$. The long-dashed curve shows the underlying distribution $`f_{\mathrm{int}}(m)`$ for this model. For these $`J`$ band data we know that the magnitude errors are dominated by photometric zero–point variations between the scan-strips that constitute the LMC sample (Cioni et al. 2000a ). These variations have a dispersion of $`0.13`$ (which significantly exceeds the formal photometric errors near the TRGB magnitude, cf. Fig. 1). In view of this, the value $`\sigma =0.126`$ inferred from the model fit is very reasonable.
Model fitting can be used as a general tool to estimate $`m_{\mathrm{TRGB}}`$ from an observed magnitude distribution. However, this technique is error-prone, since one is essentially solving a deconvolution problem in which neither the exact shape of the intrinsic magnitude distribution $`f_{\mathrm{int}}(m)`$ nor that of the kernel $`E(m)`$ is well known a priori. A more robust approach is to locate a feature in the observed distribution $`f_{\mathrm{obs}}(m)`$ that is a direct consequence of the discontinuity at $`m_{\mathrm{TRGB}}`$. Since a discontinuity corresponds (by definition) to an infinitely steep gradient, one obvious approach is to search for a maximum in the first derivative $`f_{\mathrm{obs}}^{}\mathrm{d}f_{\mathrm{obs}}/\mathrm{d}m`$. This approach has been used in several previous studies of TRGB magnitude determinations (e.g., Lee, Freedman & Madore lfm (1993)). For a model with $`a_1=a_2a`$ one can show that one expects simply $`f_{\mathrm{obs}}^{}(m)=a+\mathrm{\Delta }fE(mm_{\mathrm{TRGB}})`$, i.e., the first derivative is a Gaussian centered at $`m_{\mathrm{TRGB}}`$ plus a constant. However, the above analysis shows that $`a_1a_2`$. So while the derivative $`f_{\mathrm{obs}}^{}`$ generally does have a maximum near $`m_{\mathrm{TRGB}}`$, the structure of the first derivative is generally more complicated than a Gaussian. The heavy curve in Fig. 7b shows $`f_{\mathrm{obs}}^{}(m)`$ for the model with the parameters determined from the $`J`$ band data.
The magnitude distribution of stars on the AGB is very different from that on the RGB. While the former is approximately constant and in fact even slightly increasing to brighter magnitudes ($`a_1<0`$), the latter increases very sharply to fainter magnitudes ($`a_2>0`$). Hence, not only $`f_{\mathrm{int}}`$, but also its derivative is discontinuous at $`m_{\mathrm{TRGB}}`$. This corresponds to an infinitely steep gradient in the first derivative (see the long dashed curves in Fig. 7), which can be identified by searching for a maximum in $`f_{\mathrm{obs}}^{\prime \prime }\mathrm{d}^2f_{\mathrm{obs}}/\mathrm{d}m^2`$. For a model with $`\mathrm{\Delta }f=0`$ one can show that one expects simply that $`f_{\mathrm{obs}}^{\prime \prime }(m)=(a_2a_1)E(mm_{\mathrm{TRGB}})`$, i.e., the second derivative is a Gaussian centered at $`m_{\mathrm{TRGB}}`$. While the above discussion shows that the best fit to the data is obtained for $`\mathrm{\Delta }f0`$, the value of $`\mathrm{\Delta }f`$ is close enough to zero to ensure that $`f_{\mathrm{obs}}^{\prime \prime }(m)`$ is always modestly well approximated by a Gaussian (especially near its peak). Fig. 7c shows $`f_{\mathrm{obs}}^{\prime \prime }`$ for the model with the parameters determined from the $`J`$ band data.
While the discontinuity in $`f_{\mathrm{int}}`$ causes both a maximum in $`f_{\mathrm{obs}}^{}`$ at a position $`m_1`$ and a maximum in $`f_{\mathrm{obs}}^{\prime \prime }`$ at a position $`m_2`$, it is important to realize that neither provides a unbiased estimate of $`m_{\mathrm{TRGB}}`$. Fig. 8a shows for the model derived from the $`J`$ band data the differences $`\mathrm{\Delta }m_1m_1m_{\mathrm{TRGB}}`$ and $`\mathrm{\Delta }m_2m_2m_{\mathrm{TRGB}}`$ as function of $`\sigma `$. In absolute value, the differences increase monotonically with $`\sigma `$. The value of $`m_1`$ always provides an overestimate of $`m_{\mathrm{TRGB}}`$ while $`m_2`$ always provides an underestimate. It is important to realize that in practice, because of finite statistics, one must always apply a certain amount of smoothing to real data to obtain an adequate estimate of either $`f_{\mathrm{obs}}^{}`$ or $`f_{\mathrm{obs}}^{\prime \prime }`$. This smoothing usually takes the form of binning (e.g., Lee, Freedman & Madore lfm (1993))) or kernel smoothing (e.g., Sakai, Madore & Freedman sakai (1996)). When assessing the size of the bias terms in Fig. 8a for any particular application, the value of $`\sigma `$ along the abscissa should therefore not be taken merely as the average photometric error for the data, but should include the effect of the additional smoothing that was applied to obtain the estimate of either $`m_1`$ or $`m_2`$. While photometric errors of a few hundredths of a magnitude are often routinely achieved, the additional smoothing or binning applied during data processing is often as large as 0.1 to 0.2 magnitudes. According to Fig. 8a, this can induce systematic biases in the estimate of $`m_{\mathrm{TRGB}}`$ that are of the same order. So while this is not typically done (e.g., Sakai, Zaritsky & Kennicutt szk (1999); Nikolaev & Weinberg Nik (2000)), we do believe that such systematic biases should be calculated and corrected for.
Previous authors have generally searched for the magnitude of the TRGB by determining the position of the peak in $`f_{\mathrm{obs}}^{}`$. As far as we know, no one has yet used $`f_{\mathrm{obs}}^{\prime \prime }`$. This is presumably for the obvious reason that it is more difficult to determine the second derivative from noisy data than the first derivative. However, the situation for the DCMC catalogue differs considerably from that for most other studies. First, we have a very large number of stars, so that it is actually not a problem to accurately determine $`f_{\mathrm{obs}}^{\prime \prime }`$. Second, the random errors in the sample are relatively large. This is not because of photometric errors (which are small, cf. Fig. 1) but because of photometric zero–point variations between the scan-strips that constitute the sample. The effect of the size of the errors on the properties of $`f_{\mathrm{obs}}^{}`$ and $`f_{\mathrm{obs}}^{\prime \prime }`$ are illustrated by the dotted curves in Fig. 7, which show predictions for the same model as before, but for values of $`\sigma `$ of $`0.05`$, $`0.10`$, $`0.15`$ and $`0.20`$, respectively. We have found that the values of $`\sigma `$ appropriate for our analysis are such that the peak in $`f_{\mathrm{obs}}^{}(m)`$ is generally not the most easily recognizable feature in the data. After extensive testing we concluded that for our data $`f_{\mathrm{obs}}^{\prime \prime }(m)`$ provides a better handle on $`m_{\mathrm{TRGB}}`$ than does $`f_{\mathrm{obs}}^{}(m)`$.
In practice, we estimate the properties of the peak in $`f_{\mathrm{obs}}^{\prime \prime }(m)`$ by performing a Gaussian fit. This yields $`m_{2g}`$, the center of the best-fitting Gaussian, and $`\sigma _{2g}`$, the dispersion of the best-fitting Gaussian (in general, the value of $`\sigma _{2g}`$ is roughly of the same order as $`\sigma `$, and $`\mathrm{\Delta }m_{2g}`$ is roughly of the same order as $`\mathrm{\Delta }m_2`$). For given $`f_{\mathrm{int}}`$, both $`m_{2g}`$ and $`\sigma _{2g}`$ are unique monotonic functions of $`\sigma `$. So one can view $`\mathrm{\Delta }m_{2g}m_{2g}m_{\mathrm{TRGB}}`$ to be a function of $`\sigma _{2g}`$. The solid curve in Fig. 8b shows this function for the $`f_{\mathrm{int}}`$ parameterization derived from the $`J`$ band data.
### A.2 Implementation and formal errors
To implement our strategy we bin the observed stellar magnitudes for the region of the sky of interest into a histogram, using a fixed bin size $`b`$. As described in Section 3.1, we do the same for observations of an offset field, and subtract an appropriately scaled version of the offset field histogram from the main field histogram to obtain a foreground-subtracted histogram $`N(m)`$. We then apply a Savitzky-Golay filter (e.g., Press et al. press (1992)) to estimate the second derivative $`\mathrm{d}^2N(m)/\mathrm{d}m^2`$ at the position of each bin. This yields for bin number $`i`$
$$[\mathrm{d}^2N/\mathrm{d}m^2]_i=\underset{j=J}{\overset{J}{}}c_j[N(m)]_{i+j},$$
(5)
where the $`c_j`$ are Savitzky-Golay coefficients for the chosen value of $`J`$ and the desired derivative order $`L=2`$. The filter fits a polynomial of order $`M`$ to the data points $`[N(m)]_j`$ with $`j=iJ,\mathrm{},i+J`$, and then evaluates the $`L^{\mathrm{th}}`$ derivative of the polynomial at bin $`i`$ to estimate $`[\mathrm{d}^2N/\mathrm{d}m^2]_i`$. Once a histogram approximation to $`[\mathrm{d}^2N/\mathrm{d}m^2]`$ has been calculated, we search for a peak and fit a Gaussian in the region around the peak to obtain $`m_{2g}`$ and $`\sigma _{2g}`$ (the mean and dispersion of the best-fitting Gaussian). From these values we estimate the magnitude $`m_{\mathrm{TRGB}}`$ as
$$m_{\mathrm{TRGB}}=m_{2g}\mathrm{\Delta }m_{2g}(\sigma _{2g}),$$
(6)
where the correction term $`\mathrm{\Delta }m_{2g}(\sigma _{2g})`$ is taken from Fig. 8b. To summarize, $`m_{\mathrm{TRGB}}`$ is estimated as the position where the second derivative of the observed histogram has its maximum, plus a small correction that is based on a model for the underlying magnitude distribution $`f_{\mathrm{int}}`$.
We performed extensive Monte-Carlo simulations to assess the accuracy of the $`m_{\mathrm{TRGB}}`$ estimates produced by this algorithm. In these simulations Cloud stars are drawn from the magnitude distribution $`f_{\mathrm{int}}`$ given by Eq. (4), using as before the parameters determined from the $`J`$ band data. Foreground stars are drawn from a smooth magnitude distribution that matches that inferred from our data, both for the main field and a hypothetical offset field. To each stellar magnitude an error is added that is drawn from a Gaussian with dispersion $`\sigma `$. The numbers of stars in the simulations were chosen to match those in our datasets. In each simulation, the magnitudes thus generated are analyzed in exactly the same way as the real data to obtain $`m_{2g}`$ and $`\sigma _{2g}`$, and from these (using Eq. 6) an estimate $`\stackrel{~}{m}_{\mathrm{TRGB}}`$. This procedure is then repeated many times in Monte-Carlo fashion, and for the resulting ensemble we calculated the mean $`\stackrel{~}{m}_{\mathrm{TRGB}}`$ and dispersion $`\sigma _{m,\mathrm{TRGB}}`$ of the $`\stackrel{~}{m}_{\mathrm{TRGB}}`$ estimates, as well as the mean $`\sigma _{2g}`$ of the $`\sigma _{2g}`$. In the simulations we experimented with the choice of the algorithm parameters $`b`$, $`J`$, and $`M`$. We found that accurate results were obtained with, e.g., $`J=3`$, $`M=2`$ and a binsize $`b=0.07`$ magnitudes. These parameters were therefore generally adopted for the further analysis (with the exception of the SMC $`K_S`$ band data, for which we used the slightly larger bin size $`b=0.10`$ magnitudes). The Savitzky-Golay coefficients for this choice of parameters are $`c_j=\overline{c}_j/b^2`$, with $`\overline{c}_0=0.0476`$, $`\overline{c}_1=\overline{c}_1=0.0357`$, $`\overline{c}_2=\overline{c}_2=0`$, $`\overline{c}_3=\overline{c}_3=0.0595`$. With these parameters we found that $`|\stackrel{~}{m}_{\mathrm{TRGB}}m_{\mathrm{TRGB}}|<0.01`$ magnitudes, independent of the assumed $`\sigma `$. Hence, the algorithm produces unbiased estimates of $`m_{\mathrm{TRGB}}`$. This result was found to be rather insensitive to the precise choice of the algorithm parameters; different parameters generally yielded similar results for $`m_{\mathrm{TRGB}}`$. The formal error on a determination of $`m_{\mathrm{TRGB}}`$ from real data is obtained as follows: (i) we run simulations with the appropriate numbers of stars, for a range of $`\sigma `$ values; (ii) we identify the value of $`\sigma `$ that yields a value of $`\sigma _{2g}`$ that equals the value of $`\sigma _{2g}`$ inferred from the data; (iii) the corresponding value of $`\sigma _{m,\mathrm{TRGB}}`$ is the formal error that was sought. The errors thus inferred are listed in Table 1; typical values are $`0.02`$$`0.05`$ magnitudes.
### A.3 Assessment of systematic errors
The Monte-Carlo simulations provide accurate estimates of the formal errors in the $`m_{\mathrm{TRGB}}`$ determinations due to the combined effects of the finite number of stars and the properties of our adopted algorithm. However, they provide no insight into possible systematic errors. We have performed a number of additional tests to assess the influence of possible sources of systematic errors.
#### A.3.1 Accuracy of the correction term $`\mathrm{\Delta }m_{2g}`$
Our estimates for $`m_{\mathrm{TRGB}}`$ are obtained from Eq. (6), in which we add to the observed magnitude $`m_{2g}`$ of the $`f_{\mathrm{obs}}^{\prime \prime }(m)`$ peak a correction $`\mathrm{\Delta }m_{2g}`$ that is derived from a model. Any error in the model will change the correction $`\mathrm{\Delta }m_{2g}`$, which in turn yields a systematic error in the derived $`m_{\mathrm{TRGB}}`$. It is therefore important to understand the accuracy of the model.
There are two main parameters in fitting the model defined by Eqs. (2)–(4) to an observed histogram, namely the ‘step-size’ $`\mathrm{\Delta }f`$ of the function $`f_{\mathrm{int}}(m)`$, and the dispersion $`\sigma `$ of the convolution kernel $`E(m)`$. These parameters are highly correlated. If (as compared to the best fit model) $`\mathrm{\Delta }f`$ is increased, then an appropriate simultaneous increase in $`\sigma `$ will yield a predicted profile $`f_{\mathrm{obs}}(m)`$ that is only slightly altered. From experiments with our Monte-Carlo simulations we conclude that for all $`0.18\mathrm{\Delta }f0.38`$ one can still obtain an acceptable fit to the observed $`J`$ band magnitude histogram. At the lower end of this range we require $`\sigma =0.105`$ and at the high end $`\sigma =0.169`$, neither of which seems entirely implausible for the $`J`$ band data. The dashed curves in Fig. 8b show the correction factors $`\mathrm{\Delta }m_{2g}(\sigma _{2g})`$ for these models. These can be compared to the solid curve, which pertains to the model with $`\mathrm{\Delta }f=0.25`$ shown in Fig. 7. A typical value of $`\sigma _{2g}`$ for our data is $`0.11`$. Fig. 8b shows that for this $`\sigma _{2g}`$ the systematic error in $`\mathrm{\Delta }m_{2g}`$ (and hence $`m_{\mathrm{TRGB}}`$) due to uncertainties in $`\mathrm{\Delta }f`$ is approximately $`0.02`$ magnitudes.
The correction term $`\mathrm{\Delta }m_{2g}(\sigma _{2g})`$ that we have applied to all our data was derived from LMC data in the $`J`$ band. This would not be adequate if the shape of $`f_{\mathrm{int}}(m)`$ differs significantly among the $`I`$, $`J`$ and $`K_S`$ bands, or among the LMC and the SMC. However, visual inspection of Fig. 3 does not strongly suggest that this is the case: the shape of the observed magnitude histograms near the TRGB is similar in all cases. Quantitative analysis supports this, and demonstrated that values of $`0.18\mathrm{\Delta }f0.38`$ are adequate for all our data.
#### A.3.2 Incompleteness
In our main sample we have only included stars that were confidently detected in all three photometric bands. Fig. 3 shows that for this sample incompleteness starts to be an issue at brightnesses that are only a few tens of a magnitude fainter than the inferred $`m_{\mathrm{TRGB}}`$. One may wonder whether this could have had a systematic influence on the $`m_{\mathrm{TRGB}}`$ determinations. To assess this we applied our algorithm also to a different (extended) sample consisting of those stars that were detected in the $`I`$ and $`J`$ bands (irrespective of whether or not they were detected in $`K_S`$), which is complete to much fainter magnitudes than the main sample (heavy dashed curves in Fig. 3). The RMS difference between the $`m_{\mathrm{TRGB}}`$ estimates from the main and the extended sample (for those cases where both are available) was found to be $`0.04`$, which can be attributed entirely to the formal errors in these estimates. We therefore conclude that there is no evidence for systematic errors due to possible incompleteness.
#### A.3.3 Foreground subtraction
Our method for foreground subtraction (see Section 3.1) is based on an empirical scaling of the magnitude histogram for an offset field. To assess the effect of possible uncertainties in the foreground subtraction we have, as a test, done our analysis also without any foreground subtraction (i.e., using the thin solid curves in the $`N(m)`$ panels of Fig. 3). Even this very extreme assumption was found to change the inferred $`m_{\mathrm{TRGB}}`$ values only at the level of $`0.02`$, which can be attributed entirely to the formal errors in the estimates. We therefore conclude that there is no evidence for systematic errors due to uncertainties in the foreground subtraction.
#### A.3.4 Extinction
Extinction enters into our analysis in various ways. For the I, J and $`K_S`$ data we have performed our analysis on data that were not corrected for extinction. Instead, we apply an average extinction correction to the inferred $`m_{\mathrm{TRGB}}`$ values after the analysis. Obviously, any error in the assumed average extinction for the sample translates directly into an error in $`m_{\mathrm{TRGB}}`$. Table 1 lists for each band the shift in $`m_{\mathrm{TRGB}}`$ that would be introduced by a shift of $`+0.05`$ in the assumed $`E(BV)`$ (a shift of $`0.05`$ in the assumed $`E(BV)`$ would produce the opposite shift in $`m_{\mathrm{TRGB}}`$). It should be noted that our analysis does not assume that the extinction is constant over the region of sky under study. If there are variations in extinction then this causes an additional broadening of the convolution kernel $`E(m)`$ beyond what is predicted by observational errors alone. The width of the convolution kernel is not assumed to be known in our analysis, but is calibrated indirectly through our determination of $`\sigma _{2g}`$ (the dispersion of the $`f_{\mathrm{obs}}^{\prime \prime }(m)`$ peak). Hence, any arbitrary amount of extinction variations within the Clouds will neither invalidate our results, nor increase the formal errors.
In our calculation of the bolometric magnitudes $`m_{\mathrm{bol}}`$ of the individual stars in our sample from the observed $`J`$ and $`K_S`$ magnitudes we do correct for extinction. The effect of a change in the assumed $`E(BV)`$ affects the inferred $`m_{\mathrm{TRGB}}`$ values in a complicated way, because both the magnitudes and the colors of individual stars are affected. We therefore performed our entire analysis of the $`m_{\mathrm{bol}}`$ histograms for three separate assumed values of $`E(BV)`$, namely $`0.10`$, $`0.15`$ and $`0.20`$. From these analyses we conclude that an increase in $`E(BV)`$ of $`+0.05`$ decreases the inferred bolometric $`m_{\mathrm{TRGB}}`$ by $`0.03`$ (a shift of $`0.05`$ in the assumed $`E(BV)`$ would produce the opposite shift in $`m_{\mathrm{TRGB}}`$). As for the $`I`$, $`J`$ and $`K_S`$ data, extinction variations within the Clouds will not invalidate the results or increase the formal errors.
### A.4 Comparison to other methods
Most authors have searched for the magnitude $`m_1`$ of the peak in the first derivative $`f_{\mathrm{obs}}^{}`$ to estimate the magnitude $`m_{\mathrm{TRGB}}`$ of the TRGB discontinuity. While this is a perfectly good approach, it is important to realize that this by itself does not yield an unbiased estimate of $`m_{\mathrm{TRGB}}`$. This was pointed out previously by Madore & Freedman (mado (1995); see their Fig. 3). However, they were not overly concerned with this, since their aim was to test the limitations on determining $`m_{\mathrm{TRGB}}`$ to better than $`\pm 0.2`$ mag. As a result, it has not been common practice to estimate the bias $`\mathrm{\Delta }m_1`$ intrinsic to $`m_1`$ and correct for it. Fig. 8a also shows that for small values of $`\sigma `$ one has $`|\mathrm{\Delta }m_1|<|\mathrm{\Delta }m_2|`$, so the application of a correction may seem less important for methods based on the first derivative than for those based on the second derivative. On the other hand, it has now become possible to determine $`m_1`$ with formal errors of order $`0.1`$ mag or less (e.g., Sakai, Zaritsky & Kennicutt szk (1999); Nikolaev & Weinberg Nik (2000)), so it is important to correct for systematic biases even if one uses the first derivative, as we will illustrate.
To estimate quantitatively the size of possible biases in the results of previous authors one must do Monte-Carlo simulations for their exact observational setup and analysis procedure, which is beyond the scope of the present paper. However, as an illustration it is useful to consider the result of Nikolaev & Weinberg (Nik (2000)), who find from 2MASS data for the LMC that $`m_{\mathrm{TRGB}}(K_S)=12.3\pm 0.1`$. This corresponds to $`m_{\mathrm{TRGB}}(K_S)=12.19\pm 0.1`$ in the DENIS photometric system, which conflicts significantly with our result $`m_{\mathrm{TRGB}}(K_S)=11.98\pm 0.04`$ (see Section 4.4). Nikolaev & Weinberg derived their result from an analysis of the derivative of the observed magnitude distribution; the latter is shown and listed as a histogram with $`0.2`$ mag. bins in their Fig. 9 and Table 1. If they used the Sobel edge detection filter suggested by Madore & Freedman (mado (1995)) on this histogram, then Monte-Carlo simulations that we have done (similar to those in Section A.2) indicate that their estimate of $`m_1`$ could overestimate $`m_{\mathrm{TRGB}}`$ by as much as $`0.15\pm 0.06`$. If we correct their result for this bias, then we obtain $`m_{\mathrm{TRGB}}(K_S)=12.04\pm 0.12`$ for their data, in good agreement with our result. Romaniello et al. (rscp (1999)) use a bin size as large as $`0.25`$ mag in their analysis, and their estimate of the TRGB magnitude is therefore likely to be biased upward even more.
Our method differs from that employed by Sakai, Madore & Freedman (sakai (1996)) in that they employ kernel smoothing and estimate $`f_{\mathrm{obs}}^{}`$ as a continuous function, while we employ histograms. Sakai et al. quote as an advantage of their technique that it avoids the arbitrary choice of bin size and histogram starting point. While this is true, we have not found any evidence that this makes a significant quantitative difference. Our Monte-Carlo simulations indicate that our results obtained from histograms are unbiased to better than 0.01 mag., and we have found this to be true for all histogram starting points and a large range of reasonable bin sizes. However, we should point out that for this to be the case it is important to apply appropriate corrections for systematic biases (which applies equally to histograms estimates and kernel smoothing estimates).
A final issue worth mentioning is the estimation of the formal error in $`m_{\mathrm{TRGB}}`$. We have done this through Monte-Carlo simulations, which is probably the most robust way to do this. By contrast, Sakai, Zaritsky & Kennicutt (szk (1999)) quote as the formal error the FWHM of the observed peak in $`f_{\mathrm{obs}}^{}`$. It should be noted that this is not actually accurate (it is probably conservative). Recall from Section A.1 that for the simplified case in which $`a_1=a_2a`$ in Eq. (4), one has $`f_{\mathrm{obs}}^{}(m)=a+\mathrm{\Delta }fE(mm_{\mathrm{TRGB}})`$. Hence, the dispersion of the peak in $`f_{\mathrm{obs}}^{}(m)`$ measures the random error in the individual stellar magnitude measurements (plus whatever smoothing was applied to the data). This dispersion is independent of the number of stars in the sample ($`N`$), and therefore cannot be a measure of the formal error in $`m_{\mathrm{TRGB}}`$. The true formal error (i.e., the dispersion among the results obtained from different randomly drawn samples) scales with the number of stars as $`1/\sqrt{N}`$. |
warning/0003/hep-ph0003137.html | ar5iv | text | # Conformal Phase Transition, 𝛽-Function, and Infrared Dynamics in QCD
## 1 Introduction
The infrared dynamics of QCD can be viewed only dimly via presently available tools. Because of that, there may be surprises, as it has already happened with our understanding of nonperturbative dynamics in $`N`$=1 supersymmetric QCD . In particular, those studies showed that the conventional wisdom which has accepted that the asymptotic freedom implies confinement is not always true.
Recently, there has been considerable interest in the existence of a nontrivial conformal dynamics in 3+1 dimensional non-supersymmetric vector like gauge theories, with a relatively large number of fermion flavors $`N_f`$ . The roots of this problem go back to a work of Banks and Zaks who were first to discuss the consequences of the existence of an infrared-stable fixed point $`\alpha =\alpha ^{}`$ for $`N_f>N_f^{}`$ in vector-like gauge theories. The value $`N_f^{}`$ depends on the gauge group: in the case of SU(3) gauge group, $`N_f^{}=8`$ in the two-loop approximation.
A new insight in this problem has been, on the one hand, connected with using the results of the analysis of the Schwinger-Dyson (SD) equations describing chiral symmetry breaking in QCD (for a review, see Refs.) and, on the other hand, with the discovery of the conformal window in $`N=1`$ supersymmetric QCD .
In particular, Appelquist, Terning, and Wijewardhana suggested that, in the case of the gauge group SU($`N_c`$), the critical value $`N_f^{cr}4N_c`$ separates a phase with no confinement and chiral symmetry breaking ($`N_f>N_f^{cr}`$) and a phase with confinement and with chiral symmetry breaking ($`N_f<N_f^{cr}`$). The basic point for this suggestion was the observation that at $`N_f>N_f^{cr}`$ the value of the infrared fixed point $`\alpha ^{}`$ is smaller than a critical value $`\alpha _{cr}\frac{2N_c}{N_c^21}\frac{\pi }{3}`$, presumably needed to generate the chiral condensate .
The authors of Ref. considered only the case when the running coupling constant $`\alpha (\mu )`$ is less than the fixed point $`\alpha ^{}`$. In this case the dynamics is asymptotically free (at short distances) both at $`N_f<N_f^{cr}`$ and $`N_f^{cr}<N_f<N_f^{}\frac{11N_c}{2}`$.
Yamawaki and the author analyzed the dynamics in the whole ($`\alpha ,N_f`$) plane and suggested the ($`\alpha ,N_f`$)-phase diagram of the SU($`N_c`$) theory (see Fig. 1 below).<sup>1</sup><sup>1</sup>1This phase diagram is essentially different from the original Banks-Zaks diagram . For details, see Sec.VII in Ref.. In particular, it was pointed out that one can get an interesting non-asymptotically free dynamics when the bare coupling constant $`\alpha ^{(0)}`$ is larger than $`\alpha ^{}`$, though not very large.
The dynamics with $`\alpha ^{(0)}>\alpha ^{}`$ admits a continuum limit and is interesting in itself. Also, its better understanding can be important for establishing the conformal window in lattice computer simulations of the SU($`N_c`$) theory with such large values of $`N_f`$. In order to illustrate this, let us consider the following example. For $`N_c=3`$ and $`N_f=16`$, the value of the infrared fixed point $`\alpha ^{}`$ is small: $`\alpha ^{}`$0.04 (see below). To reach the asymptotically free phase, one needs to take the bare coupling $`\alpha ^{(0)}`$ less than this value of $`\alpha ^{}`$. However, because of large finite size effects, the lattice computer simulations of the SU(3) theory with such a small $`\alpha ^{(0)}`$ would be unreliable. Therefore, in this case, it is necessary to consider the dynamics with $`\alpha (\mu )>\alpha ^{}`$.
The existence of the phase transition with respect to $`N_f`$ in QCD raises the following question: what are the infrared properties of the QCD $`\beta `$ function for different $`N_f`$ and how the $`\beta `$ function structure reflects the existence of this phase transition? I will discuss those issues at the end of my talk but first I will discuss the dynamics in the conformal window of QCD in detail. In particular, I will consider the spectrum of low energy excitations in that dynamics in the presence of a bare fermion mass . We will see that in this case, unlike the familiar QCD with a small $`N_f`$ ($`N_f`$=2 or 3), glueballs are much lighter than bound states composed of fermions, if the value of the infrared fixed point is not too large. Another characteristic point is a strong (and simple) dependence of the masses of all the colorless bound states on the bare fermion mass, even if the latter is tiny.
This talk is based on papers .
## 2 Dynamics in the Conformal Window in QCD like Theories
I begin by recalling the basic facts concerning the two-loop $`\beta `$ function in an SU($`N_c`$) theory. The $`\beta `$ function is
$$\beta =b\alpha ^2c\alpha ^3$$
(1)
with
$`b`$ $`=`$ $`{\displaystyle \frac{1}{6\pi }}(11N_c2N_f),`$ (2)
$`c`$ $`=`$ $`{\displaystyle \frac{1}{24\pi ^2}}(34N_c^210N_cN_f3{\displaystyle \frac{N_c^21}{N_c}}N_f).`$ (3)
While these two coefficients are invariant under change of a renormalization scheme, the higher-order coefficients are scheme dependent. Actually, there is a renormalization scheme in which the two-loop $`\beta `$ function is (perturbatively) exact . We will use such a renormalization scheme.
If $`b>0`$ ($`N_f<N_f^{}\frac{11N_c}{2}`$) and $`c<0`$, the $`\beta `$ function has a zero, corresponding to a infrared-stable fixed point, at
$$\alpha =\alpha ^{}=\frac{b}{c}.$$
(4)
When $`N_f`$ is close to $`N_f^{}`$, the value of $`\alpha ^{}`$ is small. For example, from Eqs.(2), (3), and (4), one gets $`\alpha ^{}`$ 0.04, 0.14, 0.28, and 0.47 for $`N_c`$=3 and $`N_f`$=16, 15, 14, and 13, respectively.
The value of $`\alpha ^{}`$ becomes equal to $`\alpha _{cr}=\frac{2N_c}{N_c^21}\frac{\pi }{3}`$ at $`N_f`$ close to $`N_f4N_c`$. And the fixed point disappears at the value $`N_f=N_f^{}`$, when the coefficient $`c`$ becomes positive ($`N_f^{}`$ is $`N_f^{}8.05`$ for $`N_c`$=3).
The $`\beta `$ function (1) leads to the following solution for the running coupling:
$$b\mathrm{log}\left(\frac{q}{\mu }\right)=\frac{1}{\alpha (q)}\frac{1}{\alpha (\mu )}\frac{1}{\alpha ^{}}\mathrm{log}\left(\frac{\alpha (q)(\alpha (\mu )\alpha ^{})}{\alpha (\mu )(\alpha (q)\alpha ^{})}\right).$$
(5)
We emphasize that this solution is valid both for $`\alpha (\mu )<\alpha ^{}`$ and $`\alpha (\mu )>\alpha ^{}`$.
Let us first consider the case with $`\alpha (\mu )<\alpha ^{}`$. It is convenient to introduce the parameter
$$\mathrm{\Lambda }\mu \mathrm{exp}\left[\frac{1}{b\alpha ^{}}\mathrm{log}\left(\frac{\alpha ^{}\alpha (\mu )}{\alpha (\mu )}\right)\frac{1}{b\alpha (\mu )}\right].$$
(6)
Then, Eqs. (5) and (6) imply that
$$\frac{1}{\alpha (q)}=b\mathrm{log}\left(\frac{q}{\mathrm{\Lambda }}\right)+\frac{1}{\alpha ^{}}\mathrm{log}\left(\frac{\alpha (q)}{\alpha ^{}\alpha (q)}\right).$$
(7)
Taking $`q=\mathrm{\Lambda }`$, we find that
$$\frac{\alpha ^{}}{1+e^1}0.73\alpha ^{}<\alpha (\mathrm{\Lambda })<\alpha ^{}.$$
(8)
One may think that $`\mathrm{\Lambda }`$ plays here the same role as $`\mathrm{\Lambda }_{QCD}`$ in the confinement phase. However, as we will see, its physical meaning is somewhat different.
Eq. (7) implies that
$$\alpha (q)\frac{1}{b\mathrm{log}\frac{q}{\mathrm{\Lambda }}}$$
(9)
for $`q>>\mathrm{\Lambda }`$ (the usual behavior in asymptotically free theories), and
$$\alpha (q)\frac{\alpha ^{}}{1+e^1(\frac{q}{\mathrm{\Lambda }})^{b\alpha ^{}}}$$
(10)
for $`q<<\mathrm{\Lambda }`$, governed by the infrared fixed point $`\alpha ^{}`$.
Let us turn to a less familiar case with $`\alpha (\mu )>\alpha ^{}`$. One still can use Eq.(5). Introduce now the parameter $`\stackrel{~}{\mathrm{\Lambda }}`$ as
$$\stackrel{~}{\mathrm{\Lambda }}\mu \mathrm{exp}\left[\frac{1}{b\alpha ^{}}\mathrm{log}\left(\frac{\alpha (\mu )\alpha ^{}}{\alpha (\mu )}\right)\frac{1}{b\alpha (\mu )}\right]$$
(11)
(compare with Eq.(6)). Then, Eqs.(5) and (11) imply
$$\frac{1}{\alpha (q)}=b\mathrm{log}\frac{q}{\stackrel{~}{\mathrm{\Lambda }}}+\frac{1}{\alpha ^{}}\mathrm{log}\left(\frac{\alpha (q)}{\alpha (q)\alpha ^{}}\right).$$
(12)
What is the meaning of $`\stackrel{~}{\mathrm{\Lambda }}`$? It is a Landau pole at which $`\alpha (q)|_{q=\stackrel{~}{\mathrm{\Lambda }}}=\mathrm{}`$. Indeed, taking $`q=\stackrel{~}{\mathrm{\Lambda }}`$ in Eq.(12), one gets
$$\frac{1}{\alpha (\stackrel{~}{\mathrm{\Lambda }})}=\frac{1}{\alpha ^{}}\mathrm{log}\frac{\alpha (\stackrel{~}{\mathrm{\Lambda }})}{\alpha (\stackrel{~}{\mathrm{\Lambda }})\alpha ^{}}.$$
(13)
The only solution of this equation is $`\alpha (\stackrel{~}{\mathrm{\Lambda }})=\mathrm{}`$.
The presence of the Landau pole implies that the dynamics is not asymptotically free. To get a more insight in this dynamics, let us introduce an ultraviolet cutoff $`M`$ with the bare coupling constant $`\alpha ^{(0)}\alpha (q)|_{q=M}`$. Now all momenta $`q`$ satisfy $`qM`$.
Eq.(12) implies that at finite $`\alpha ^{(0)}=\alpha (M)`$, the cutoff $`M`$ is less than $`\stackrel{~}{\mathrm{\Lambda }}`$, with $`\alpha (\stackrel{~}{\mathrm{\Lambda }})=\mathrm{}`$. Therefore the Landau pole is unreachable in the theory with cut off $`M`$ and with $`\alpha ^{(0)}<\mathrm{}`$. Still one can of course use $`\stackrel{~}{\mathrm{\Lambda }}`$ (11) for a convenient parametrization of the running coupling $`\alpha (q)`$ (see Eq.(12)). However, one should remember that momenta $`q`$ satisfy $`qM<\stackrel{~}{\mathrm{\Lambda }}`$.
Eq.(12) implies that
$$\alpha ^2(q)\frac{\alpha ^{}}{2b\mathrm{log}\frac{\stackrel{~}{\mathrm{\Lambda }}}{q}}$$
(14)
for $`\alpha (q)>>\alpha ^{}`$, and
$$\alpha (q)\frac{\alpha ^{}}{1e^1(\frac{q}{\stackrel{~}{\mathrm{\Lambda }}})^{b\alpha ^{}}}$$
(15)
when $`\alpha (q)`$ is close to $`\alpha ^{}`$, i.e. when $`\alpha (q)\alpha ^{}<<\alpha ^{}`$. Thus, now $`\alpha (q)`$ approaches the fixed point $`\alpha ^{}`$ from above (compare with Eq.(10)). And, in general, Eq.(11) implies that $`\alpha (q)`$ monotonically decrease with $`q`$, from $`\alpha (q)=\alpha ^{(0)}`$ at $`q=M`$ to $`\alpha (q)=\alpha ^{}`$ at $`q=0`$.
Does a meaningful continuum limit exist in this case? The answer is of course ”yes”. As it follows from Eq.(12), when $`M`$ (and therefore $`\stackrel{~}{\mathrm{\Lambda }}`$) goes to infinity, and the bare coupling $`\alpha ^{(0)}>\alpha ^{}`$ is arbitrary but fixed, $`\alpha (q)`$ is equal to the fixed value, $`\alpha (q)=\alpha ^{}`$, for all $`q<\mathrm{}`$. Therefore it is a non-trivial conformal field theory.
So far we considered the solution for $`\alpha (q)`$ connected with the perturbative (and perturbatively exact in the ’t Hooft renormalization scheme ) $`\beta `$ function (1). However, unlike ultraviolet stable fixed points, defining dynamics at high momenta, infrared-stable fixed points (defining dynamics at low momenta) are very sensitive to nonperturbative dynamics leading to the generation of particle masses. For example, if fermions acquire a dynamical mass, they decouple from the infrared dynamics, and therefore the perturbative infrared fixed point (4) will disappear.
The phase diagram in the ($`\alpha ^{(0)},N_f`$)-plane in this theory was suggested in Ref.. It is shown in Fig. 1. The left-hand portion of the curve in this figure coincides with the line of the infrared-stable fixed points $`\alpha ^{}(N_f)`$ in Eq.(4). It separates two symmetric phases, $`S_1`$ and $`S_2`$, with $`\alpha ^{(0)}<\alpha ^{}`$ and $`\alpha ^{(0)}>\alpha `$, respectively. Its lower end is $`N_f=N_f^{cr}`$ (with $`N_f^{cr}4N_c`$ if $`\alpha _{cr}\frac{2N_c}{N_c^21}\frac{\pi }{3}`$): at $`N_f^{}<N_f<N_f^{cr}`$ the infrared fixed point is washed out by generating a dynamical fermion mass.
The horizontal, $`N_f=N_f^{cr}`$, line describes a phase transition between the symmetric phase $`S_1`$ and the phase with confinement and chiral symmetry breaking. As it was suggested in Refs., based on a similarity of this phase transition with that in quenched $`QED_4`$ and in $`QED_3`$ , there is the following scaling law for $`m_{dyn}^2`$:
$$m_{dyn}^2\mathrm{\Lambda }_{cr}^2\mathrm{exp}\left(\frac{C}{\sqrt{\frac{\alpha ^{}(N_f)}{\alpha _{cr}}1}}\right)$$
(16)
where the constant $`C`$ is of order one and $`\mathrm{\Lambda }_{cr}`$ is a scale at which the running coupling is of order $`\alpha _{cr}`$.
It is a continuous phase transition with an essential singularity at $`N_f=N_f^{cr}`$. The characteristic point of this phase transition is that the critical line $`N_f=N_f^{cr}`$ separates phases with essentially different spectra of low energy excitations and the different structure of the equation for the divergence of the dilatation current (i.e. with essentially different realizations of the conformal symmetry) . It was called the conformal phase transition in Ref..
At present it is still unclear whether the phase transition on the line $`N_f=N_f^{cr}`$ is indeed a continuous phase transition with an essential singularity or it is a first order phase transition . However, anyway, the two properties (the abrupt change of the spectrum of excitations and the different structure of the equation for the divergence of the dilatation current in those two phases) have to take place.
At last, the right-hand portion of the curve on the diagram occurs because at large enough values of the bare coupling, spontaneous chiral symmetry breaking takes place for any number $`N_f`$ of fermion flavors. This portion describes a phase transition called a bulk phase transition in the literature, and it is presumably a first order phase transition. <sup>2</sup><sup>2</sup>2The fact that spontaneous chiral symmetry breaking takes place for any number of fermion flavors, if $`\alpha ^{(0)}`$ is large enough, is valid at least for lattice theories with Kogut-Susskind fermions. Notice however that since the bulk phase transition is a lattice artifact, the form of this portion of the curve can depend on the type of fermions used in simulations (for details, see Ref.). The vertical line ends above $`N_f`$=0 since in pure gluodynamics there is apparently no phase transition between weak-coupling and strong-coupling phases.
Up to now we have considered the case of a chiral invariant action. But how will the dynamics change if a bare fermion mass term is added in the action? This question is in particular relevant for lattice computer simulations: for studying a chiral phase transition on a finite lattice, it is necessary to introduce a bare fermion mass. We will show that adding even an arbitrary small bare fermion mass results in a dramatic changing the dynamics both in the $`S_1`$ and $`S_2`$ phases.
Recall that in the case of confinement SU($`N_c`$) theories, with a small, $`N_f<N_f^{cr}`$, number of fermion flavors, the role of a bare fermion mass $`m^{(0)}`$ is minor if $`m^{(0)}<<\mathrm{\Lambda }_{QCD}`$ (where $`\mathrm{\Lambda }_{QCD}`$ is a confinement scale). The only relevant consequence is that massless Nambu-Goldstone pseudoscalars get a small mass (the PCAC dynamics).
The reason for that is the fact that the scale $`\mathrm{\Lambda }_{QCD}`$, connected with a scale anomaly, describes the breakdown of the conformal symmetry connected both with perturbative and nonperturbative dynamics: the running coupling and the formation of bound state. Certainly, a small bare mass $`m^{(0)}<<\mathrm{\Lambda }_{QCD}`$ is irrelevant for the dynamics of those bound states.
Now let us turn to the phase $`S_1`$ and $`S_2`$, with $`N_f>N_f^{cr}`$. At finite $`\mathrm{\Lambda }`$ in $`S_1`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ in $`S_2`$, there is still conformal anomaly: because of the running of the effective coupling constant, the conformal symmetry is broken. It is restored only if $`\mathrm{\Lambda }0`$ in $`S_1`$ and $`\stackrel{~}{\mathrm{\Lambda }}>M\mathrm{}`$ in $`S_2`$. However, the essential difference with respect to confinement theories is that both $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ have nothing with the dynamics forming bound states: since at $`N_f>N_f^{cr}`$ the effective coupling is relatively weak, it is impossible to form bound states from $`\mathrm{𝑚𝑎𝑠𝑠𝑙𝑒𝑠𝑠}`$ fermions and gluons (recall that the $`S_1`$ and $`S_2`$ phases are chiral invariant).
Therefore the absence of a mass for fermions and gluons is a key point for not creating bound states in those phases. The situation changes dramatically if a bare fermion mass is introduced: indeed, even weak gauge, Coulomb-like, interactions can easily produce bound states composed of massive constituents, as it happens, for example, in QED, where electron-positron (positronium) bound states are present.
To be concrete, let us first consider the case when all fermions have the same bare mass $`m^{(0)}`$. It leads to a mass function $`m(q^2)B(q^2)/A(q^2)`$ in the fermion propagator $`G(q)=(\widehat{q}A(q^2)B(q^2))^1`$. The current fermion mass $`m`$ is given by the relation
$$m(q^2)|_{q^2=m^2}=m.$$
(17)
For the clearest exposition, let us consider a particular theory with a finite cutoff $`M`$ and the bare coupling constant $`\alpha ^{(0)}=\alpha (q)|_{q=M}`$ being not far away from the fixed point $`\alpha ^{}`$. Then, the mass function is changing in the ”walking” regime with $`\alpha (q^2)\alpha ^{}`$. It is
$$m(q^2)m^{(0)}\left(\frac{M}{q}\right)^{\gamma _m}$$
(18)
where the anomalous dimension $`\gamma _m1(1\frac{\alpha ^{}}{\alpha _{cr}})^{1/2}`$ . Eqs.(17) and (18) imply that
$$mm^{(0)}\left(\frac{M}{m^{(0)}}\right)^{\frac{\gamma _m}{1+\gamma _m}}.$$
(19)
There are two main consequences of the presence of the bare mass:
(a) bound states, composed of fermions, occur in the spectrum of the theory. The mass of a n-body bound state is $`M^{(n)}nm`$;
(b) At momenta $`q<m`$, fermions and their bound states decouple. There is a pure SU($`N_c`$) Yang-Mills theory with confinement. Its spectrum contains glueballs.
To estimate glueball masses, notice that at momenta $`q<m`$, the running of the coupling is defined by the parameter $`\overline{b}`$ of the Yang-Mills theory,
$$\overline{b}=\frac{11}{6\pi }N_c.$$
(20)
Therefore the glueball masses $`M_{gl}`$ are of order
$$\mathrm{\Lambda }_{YM}m\mathrm{exp}(\frac{1}{\overline{b}\alpha ^{}}).$$
(21)
For $`N_c=3`$, we find from Eqs.(2), (3), and (20) that $`\mathrm{exp}(\frac{1}{\overline{b}\alpha ^{}})`$ is $`6\times 10^7`$, $`2\times 10^2`$, $`10^1`$, and $`3\times 10^1`$ for $`N_f`$=16, 15, 14, and 13, respectively. Therefore at $`N_f`$=16, 15 and 14, the glueball masses are essentially lighter than the masses of the bound states composed of fermions. The situation is similar to that in confinement QCD with heavy quarks, $`m>>\mathrm{\Lambda }_{QCD}`$. However, there is now a new important point: in the conformal window, any value of $`m^{(0)}`$ (and therefore $`m`$) is ”heavy”: the fermion mass $`m`$ sets a new scale in the theory, and the confinement scale $`\mathrm{\Lambda }_{YM}`$ (21) is less, and rather often much less, than this scale $`m`$.
This leads to a spectacular ”experimental” signature of the conformal window in lattice computer simulations: glueball masses rapidly, as $`(m^{(0)})^{\frac{1}{1+\gamma _m}}`$, decrease with the bare fermion mass $`m^{(0)}`$ for all values of $`m^{(0)}`$ less than cutoff $`M`$.
Few comments are in order:
(1) The phases $`S_1`$ and $`S_2`$ have essentially the same long distance dynamics. They are distinguished only by their dynamics at short distances: while the dynamics of the phase $`S_1`$ is asymptotically free, that of the phase $`S_2`$ is not. In particular, when all fermions are massive (with the current mass $`m`$), the continuum limit $`M\mathrm{}`$ of the $`S_2`$-theory is a non-asymptotically free confinement theory. Its spectrum includes colorless bound states composed of fermions and gluons. For $`q<m`$ the running coupling $`\alpha (q)`$ is the same as in pure SU($`N_c`$) Yang-Mills theory, and for all $`q>m`$ $`\alpha (q)`$ is very close to $`\alpha ^{}`$ (”walking”, actually, ”standing” dynamics). For those values $`N_f`$ for which $`\alpha ^{}`$ is small (as $`N_f`$=16, 15 and 14 at $`N_c`$=3), glueballs are much lighter than the bound states composed of fermions. Notice that, unlike the case with $`m=0`$, there exists an S-matrix in this theory.
(2) In order to get the clearest exposition, we assumed such estimates as $`N_f^{cr}4N_c`$ for $`N_f^{cr}`$ and $`\gamma _m=1\sqrt{1\frac{\alpha ^{}}{\alpha _{cr}}}`$ for the anomalous dimension $`\gamma _m`$. While the latter should be reasonable for $`\alpha ^{}<\alpha _{cr}`$ (and especially for $`\alpha ^{}<<\alpha _{cr}`$) , the former is based on the assumption that $`\alpha _{cr}\frac{2N_c}{N_c^21}\frac{\pi }{3}`$ which, though seems reasonable, might be crude for some values of $`N_c`$. It is clear however that the dynamical picture presented in this paper is essentially independent of those assumptions.
(3) So far we have considered the case when all fermions have the same bare mass $`m^{(0)}`$. The generalization to the case when different fermions may have different bare masses is evident.
(4) Lattice computer simulations of the SU(3) theory with a relatively large number of $`N_f`$ indeed indicate on the existence of a symmetric phase.
However, the value of the critical number $`N_f^{cr}`$ is different in different simulations: it varies from $`N_f^{cr}=6`$ through $`N_f^{cr}=12`$ .
I hope that the signature of the conformal window considered in this talk can be useful to settle this important issue.
## 3 Padé-Summation Approach to QCD $`\beta `$-function Infrared Properties
How does the structure of the QCD $`\beta `$ function reflect the existence of the phase transition with respect to $`N_f`$? This problem has been addressed in the work .
The impetus for that work was the structure of the $`\beta `$ function of the $`SU(N_c)`$ SUSY gluodynamics which is known exactly if no matter fields are present :
$$\beta (x)=\frac{3N_cx^2}{4}\left[\frac{1}{1N_cx/2}\right];x\frac{\alpha }{\pi }.$$
(22)
This structure shows the following two noticeable features:
a) Eq.(22) implies that, besides the conventional asymptotically free phase with $`x<2/N_c`$, there exists a strong ultraviolet phase in which the coupling $`x`$ is greater than the $`\beta `$-function pole at $`x=2/N_c`$ . In that phase, the value $`x=\mathrm{}`$ is an ultraviolet fixed point. As in the case of the infrared fixed point considered in the previous section, these two phases share common infrared properties. Because of that, the $`\beta `$-function pole is called an infrared attractor. Unlike theories with an infrared fixed point, theories with an infrared attractor correspond to the confinement phase.
In Ref. , we addressed whether Padé-summations of the $`\overline{MS}`$ QCD $`\beta `$ function for a given number of flavors exhibit an infrared fixed point, or alternatively, an infrared attractor. The main results are the following. Below an approximant-dependent flavor threshold $`(6N_f8)`$, the Padé-summation $`\beta `$ functions incorporating $`[2|1],[1|2],[2|2],[1|3]`$, and $`[3|1]`$ approximants whose Maclaurin expansions match known higher-than-one-loop contributions to the $`\beta `$-function series always exhibit a positive pole prior to the occurrence of their first positive zero, precluding any identification of this first positive zero as an infrared fixed point. This result is shown to be true regardless of the magnitude of the presently-unknown five-loop $`\beta `$-function contribution explicitly appearing within Padé-summation $`\beta `$ functions incorporating $`[2|2],[1|3]`$, and $`[3|1]`$ approximants. Like in the case of supersymmetric gluodynamics, the pole in question suggests the occurrence of dynamics in which both a strong and an asymptotically free phase share a common infrared attractor. As $`N_f`$ increases above an approximant-dependent flavor threshold, Padé-summation $`\beta `$ functions exhibit dynamics controlled by an infrared fixed point. This fixed point decreases in magnitude with increasing flavor number.
Thus utilizing Padé-summation QCD $`\beta `$ functions, we obtain a good degree of agreement with infrared properties predicted via the ’t Hooft renormalization scheme in which the $`\beta `$ function is truncated subsequent to two-loop order. It is noticeable that the infrared structure of the $`\beta `$ function we obtained in QCD with a small number of flavors is similar to that of the $`\beta `$ function in SUSY gluodynamics: there are strong arguments in the literature in the support of essentially the same mechanism of confinement in those theories.
## 4 Acknowledgments
I am grateful to the organizers of the TMU-Yale Symposium, in particular, Hisakazu Minakata and Noriaki Kitazawa, for their warm hospitality. My special thanks to Koichi Yamawaki for his hospitality during my stay at Nagoya University. This work was supported by the Grant-in-Aid of Japan Society for the Promotion of Science No. 11695030. |
warning/0003/quant-ph0003016.html | ar5iv | text | # Classical and quantum mechanics on information spaces with applications to cognitive, psychological, social and anomalous phenomena
## 1 Introduction
We develop classical and quantum formalisms on information spaces. Basic objects of this model are so called transformers of information; basic processes are information processes. Our main aim is a description of classical and quantum dynamics of information states.
This information dynamics may be fruitful in the study of cognitive, psychological and social processes. Here flows of information are more important than flows of matter. We think that it would be possible to explain some aspects of the process of thinking and psychological, social and anomalous phenomena on the basis of our model. Thus the readers who are only interested in applications to cognitive sciences, sociology and psychology may consider our model as only a new apparatus to investigate these phenomena.
Our model of information reality can be considered as an attempt to extend the standard model of physical reality. We interpret material objects as a particular class of transformers of information (which are characterized by stable or slowly changing information states). On the other hand, our model might be used for the description of information flows which are not directly related to flows of matter. These are conscious, social (and even anomalous) information processes.
Different models of information and cognitive reality have been discussed by many scientists in relation to foundations of quantum physics -, cognitive sciences and psychology - and anomalous phenomena -.
We use a new mathematical apparatus to describe information reality (”the world of ideas”). Many authors discuss the idea that such ”ideal” objects as ideas, consciousness, information can play an important role to provide the right picture of physical reality. However, typically they use (with some modifications) the standard mathematical model based on the description of physical reality by real numbers. In particular, many of them discuss a ”conscious field”, but they try to describe this object as a new field on the standard real space-time. We think that some of cognitive processes could not be described by using the real model of physical reality. There is simply no place for such phenomena in this model. The real model was created to describe a particular class of physical phenomena (material objects). This model does not play an exceptional role. We need not try to input all physical phenomena into this real model of reality. There can be other models of physical reality. We propose to describe physical reality by using information spaces (see Appendix 1).
From our viewpoint real spaces (Newton’s absolute space or spaces of general relativity) give only a particular class of information spaces. These real information spaces are characterized by the special system for the coding of information and the special distance on the space of vectors of information. Any natural number $`m>1`$ can be chosen as the basis of the coding system. Each $`x[0,1]`$ can be presented in the form:
$$x=a_0a_1\mathrm{}a_n\mathrm{},$$
(1)
where $`a_j=1,\mathrm{},m1,`$ are digits. We denote the set of all sequences of the form (1) by the symbol $`X_m.`$ For example, let us fix $`m=10.`$ One of the main properties of the real cording system is the identification of the form:
$$10\mathrm{}0\mathrm{}=09\mathrm{}9\mathrm{};010\mathrm{}.0\mathrm{}=009\mathrm{}9\mathrm{};\mathrm{}$$
(2)
In fact, this identification is closely connected with the order structure on the real line $`𝐑`$ (and the metric related to this order structure). For each $`x`$, there exist ”right” and ”left” neighborhoods; there exist arbitrary small right and left shifts. The identification (2) is connected with the description of left neighborhoods.
Example 1.1. Let $`x=10\mathrm{}0\mathrm{}.`$ Then $`x`$ can be approximated from the left hand side with an arbitrary precision by numbers of the form $`y=09\mathrm{}90\mathrm{}.`$
The following description of right neighborhoods will be very important in our further considerations.
($`AS`$) Let $`x=a_0\mathrm{}a_m\mathrm{}.`$ Then the numbers (vectors of information) which are close to the $`x`$ from the right hand side have the form $`y=b_0\mathrm{}b_m\mathrm{},`$ where $`a_0=b_0,\mathrm{},a_m=b_m`$ for sufficiently large $`m.`$
This nearness has a natural information interpretation: ($`AS`$) implies the ability to form associations for cognitive systems which use this nearness to compare vectors of information. By ($`AS`$) two communications (two ideas in a model of human thinking, - ) which have the same codes for sufficiently large number of first (the most important) positions in cording sequences are identified by a comparator of a cognitive system.
Numbers (vectors of information) which are close to $`x`$ from the left hand side could not be characterized in the same way (see Example 1.1, there $`x`$ and $`y`$ are very close but their codes differ strongly).
Conclusion. The system of real numbers has been created as a coding system for information which the consciousness receives from reality. The main properties of this coding system are the order structure on the set of information vectors<sup>2</sup><sup>2</sup>2Of course, the idea about an order structure is a consequence of properties of the special system which is used for observations of reality. and the restricted ability (see ($`AS`$) ) to form associations.
Finally, we pay attention to the ”universal coding property” of the real system: any natural number $`m>1`$ can be used as the basis of this system. Thus it is assumed that any information process can be equivalently described by using, for example, 2-bits coding or 1997-bits coding.
All these properties of the real coding system were incorporated in every physical model <sup>3</sup><sup>3</sup>3From this point of view there is no large difference between Newton’s absolute space and real manifolds used in general relativity..
I do not think that all information processes have an order structure. On the other hand, the scale of coding system $`m>1`$ may play an important role in a description of an information process.
Let us ”modify” the real coding system. We eliminate the identification (2). Since now, there is no order structure on the set $`X_m`$ of information vectors. We consider on $`X_m`$ the nearness defined by ($`AS`$)<sup>4</sup><sup>4</sup>4Thus here all information is considered from the viewpoint of associations.. This nearness can be described by a metric. The corresponding (complete) metric space is isomorphic to the ring of so called $`m`$-adic integers $`𝐙_m`$ (see and section 2). Therefore it is natural to use $`m`$-adic numbers for a description of information processes. Mathematically it is convenient to use prime numbers $`m=p>1`$ (see ). We arrive to the domain of an extended mathematical formalism, $`p`$-adic analysis.
To use $`p`$-adic numbers in physics is not a new idea (see - ). A new idea is to use them for a description of information (in particular, cognitive ) processes. On the other hand, apparatus which has been developed in $`p`$-adic quantum physics may be fruitfully used in our model.
We develop a quantum formalism for information systems. The mathematical basis for this formalism has been presented in , , , , . In this paper we apply the $`p`$-adic quantum formalism to information systems. As in ordinary quantum mechanics over the reals, the problem of an interpretation plays the important role in information quantum mechanics. Of course, all difficulties of an interpretation of the ordinary quantum theory (see, for example, , – ) are reproduced in the information quantum theory. There are many viewpoints on an interpretation of the quantum theory (which may be very different). However, they are mainly based on the following two general interpretations of a quantum state: (1) an individual (or orthodox Copenhagen) interpretation by which a quantum state provides the complete description an individual quantum system; (2) an ensemble (or statistical) interpretation by which a quantum state provides the description of a statistical ensemble of quantum systems. In fact, by analysing the process of measurement for information quantum systems we understood that we have to follow the ensemble interpretation. This analysis might be also useful for better understanding of the ordinary quantum formalism on real space.
## 2 Systems of $`p`$-adic numbers
First we present some facts about $`p`$-adic numbers.
The field of real numbers $`𝐑`$ is constructed as the completion of the field of rational numbers $`𝐐`$ with respect to the metric $`\rho (x,y)`$ $`=`$ $`|xy|`$, where $`||`$ is the usual valuation given by the absolute value. The fields of $`p`$-adic numbers $`𝐐_p`$ are constructed in a corresponding way, but using other valuations. For a prime number $`p`$, the $`p`$-adic valuation $`||_p`$ is defined in the following way. First we define it for natural numbers. Every natural number $`n`$ can be represented as the product of prime numbers, $`n`$ $`=`$ $`2^{r_2}3^{r_3}\mathrm{}p^{r_p}\mathrm{}`$, and we define $`|n|_p`$ $`=`$ $`p^{r_p}`$, writing $`|0|_p`$ $`=0`$ and $`|n|_p`$ $`=`$ $`|n|_p`$. We then extend the definition of the $`p`$-adic valuation $`||_p`$ to all rational numbers by setting $`|n/m|_p`$ $`=`$ $`|n|_p/|m|_p`$ for $`m`$ $``$ $`0`$. The completion of $`𝐐`$ with respect to the metric $`\rho _p(x,y)`$ $`=`$ $`|xy|_p`$ is the locally compact field of $`p`$-adic numbers $`𝐐_p`$. The number fields $`𝐑`$ and $`𝐐_p`$ are unique in a sense, since by Ostrovsky’s theorem (see ) $`||`$ and $`||_p`$ are the only possible valuations on $`𝐐`$, but have quite distinctive properties.
Unlike the absolute value distance $`||`$, the $`p`$-adic valuation satisfies the strong triangle inequality $`|x+y|_p\mathrm{max}[|x|_p,|y|_p],x,y𝐐_p`$
Write $`U_r(a)`$ $`=`$ $`\{x𝐐_p:|xa|_pr\}`$ and $`U_r^{}(a)`$ $`=`$ $`\{x𝐐_p:|xa|_p<r\},`$ where $`r`$ $`=`$ $`p^n`$ and $`n`$ $`=`$ $`0`$, $`\pm 1`$, $`\pm 2`$, $`\mathrm{}`$. These are the “closed” and “open” balls in $`𝐐_p`$ while the sets $`S_r(a)`$ $`=`$ $`\{xK:|xa|_p=r\}`$ are the spheres in $`𝐐_p`$ of such radii $`r`$. These sets (balls and spheres) have a somewhat strange topological structure from the viewpoint of our usual Euclidean intuition: they are both open and closed at the same time, and as such are called clopen sets. Another interesting property of $`p`$-adic balls is that two balls have nonempty intersection if and only if one of them is contained in the other. Also, we note that any point of a $`p`$-adic ball can be chosen as its center, so such a ball is thus not uniquely characterized by its center and radius. Finally, any $`p`$-adic ball $`U_r(0)`$ is an additive subgroup of $`𝐐_p`$, while the ball $`U_1(0)`$ is also a ring, which is called the ring of $`p`$-adic integers and is denoted by $`𝐙_p`$.
Any $`x`$ $``$ $`𝐐_p`$ has a unique canonical expansion (which converges in the $`||_p`$–norm) of the form $`x=a_n/p^n+\mathrm{}a_0+\mathrm{}+a_kp^k+\mathrm{}`$ where the $`a_j`$ $``$ $`\{0,1,\mathrm{},p1\}`$ are the “digits” of the $`p`$-adic expansion. The elements $`x`$ $``$ $`𝐙_p`$ have the expansion $`x=a_0+\mathrm{}+a_kp^k+\mathrm{}`$ and can thus be identified with the sequences of digits $`x=a_0\mathrm{}a_k\mathrm{}.`$
The $`p`$-adic exponential function $`e^x=_{n=0}^{\mathrm{}}\frac{x^n}{n!}.`$ The series converges in $`𝐐_p`$ if
$$|x|_pr_p,\text{where}r_p=1/p,p2\text{and}r_2=1/4.$$
(3)
$`p`$-adic trigonometric functions $`\mathrm{sin}x`$ and $`\mathrm{cos}x`$ are defined by the standard power series. These series have the same radius of convergence $`r_p`$ as the exponential series.
If, instead of a prime number $`p`$, we start with an arbitrary natural number $`m>1`$ we construct the system of so called $`m`$-adic numbers $`𝐐_m`$ by completing $`𝐐`$ with respect to the $`m`$-adic metric $`\rho _m(x,y)`$ $`=`$ $`|xy|_m`$ which is defined in a similar way to above. However, this system is in general not a field as there may exist divisors of zero.
## 3 Dynamics on information spaces
The rings of $`p`$-adic integers $`𝐙_p`$ can be used as mathematical models for information spaces. Each element $`x=_{j=0}^{\mathrm{}}\alpha _jp^j`$ can be identified with a sequence
$$x=\alpha _0\alpha _1\mathrm{}\alpha _N\mathrm{},\alpha _j=0,1,\mathrm{},p1.$$
(4)
Such sequences are interpreted as coding sequences (in the alphabet $`A_p=\{0,1,\mathrm{},p1\}`$ with $`p`$ letters) for some amounts of information. The $`p`$-adic metric $`\rho _p(x,y)=|xy|_p`$ on $`𝐙_p`$ corresponds to the nearness ($`AS`$) for information sequences. We choose the space $`X=𝐙_p`$ (or multidimensional spaces $`X=𝐙_p^N`$) for the description of information. The $`X`$ is said to be information space.
Everywhere below we shall use the abbreviation $`\mathrm{"}I\mathrm{"}`$ for the word information (for example, information space = $`I`$-space).
Remark 3.1. Different information phenomena can be described by different mathematical models for $`I`$-spaces. The $`p`$-adic model for $`I`$-spaces is the simplest from the mathematical point of view.
Objects which ”live” in $`I`$-spaces are said to be transformers of information ($`I`$-transformers). $`I`$-transformers are not characterized by localization in information $`p`$-adic space (or real space). They are characterized by the ability to receive an external information and transform it in a new information.
Each $`I`$-transformer $`\tau `$ has internal clocks. A state of the clocks is described by an $`I`$-vector $`tT=𝐙_p`$ which is called information time. The $`I`$-time can have different interpretations in different $`I`$-models. If $`\tau `$ is a conscious system then $`t`$ is (self-recognized) time of the evolution of this system. We can say about psychological time of an individual or about (collective) social time of a group of individuals. In fact, we have not to image $`t`$ as an ordered sequence of time counts. This is only information with describes evolution of $`\tau .`$ In principle, there is no direct relation between $`I`$-time and ”physical” time that is used in the model over the reals.
At each instant $`tT`$ of $`I`$-time there is defined a total information state ($`I`$-state) $`q(t)X`$ of $`\tau .`$ It describes the position of $`\tau `$ in the $`I`$-space $`X`$. The ”life”-trajectory of $`\tau `$ can be identified with the trajectory $`q(t)`$ in $`X`$.
An $`I`$-transformer can be imagine as a kind of Turing machine. Let us consider a free $`I`$-transformer $`\tau _{\mathrm{fr}}`$ (i.e., an $`I`$-transformer which does not interact with other $`I`$-transformers and $`I`$-fields). At the instant of $`I`$-time $`t`$ the $`\tau _{\mathrm{fr}}`$ has the $`I`$-state $`q(t)=(\alpha _0(t),\alpha _1(t),\mathrm{},\alpha _k(t),\mathrm{})`$ (an infinite ribbon with symbols belonging to the alphabet $`A_p=\{0,1,\mathrm{},p1\}).`$ During an $`I`$-time interval $`\mathrm{\Delta }t`$ this state is transformed in a new state $`q(t+\mathrm{\Delta }t)=(\alpha _0(t+\mathrm{\Delta }t),\alpha _1(t+\mathrm{\Delta }t),\mathrm{},\alpha _k(t+\mathrm{\Delta }t),\mathrm{})`$ (a new infinite ribbon with symbols belonging to the alphabet $`A_p).`$ The law of transformation depends on internal $`I`$-parameters $`s`$ which determine the internal structure of $`\tau _{\mathrm{fr}}.`$ In the general case an $`I`$-transformer $`\tau `$ interact with other $`I`$-transformers $`\tau _j,j=1,\mathrm{},N`$ and $`I`$-fields $`\varphi _i(x),i=1,\mathrm{},M.`$ These interactions change continuously internal $`I`$-parameters $`s=s(t,q_{\tau _j}(t)),\varphi _i(q_\tau (t))).`$
For example, a cognitive system $`\tau `$ which is isolated from external $`I`$-flows can be considered as a free $`I`$-transformer $`\tau _{\mathrm{fr}}.`$ Here $`q(t)`$ gives the evolution of $`\tau _{\mathrm{fr}}`$ in ‘space of ideas’; $`I`$-parameters $`s`$ are determined by the neural structure of $`\tau _{\mathrm{fr}}.`$ In general case the cognitive system $`\tau _{\mathrm{fr}}`$ interact with other cognitive systems and material objects (the latter interactions are also considered by $`\tau _{\mathrm{fr}}`$ as $`I`$-interactions) and $`I`$-fields. These interactions change continuously (with respect to $`I`$-time of $`\tau _{\mathrm{fr}})`$ the transformation law, $`q(t)q(t+\mathrm{\Delta }t).`$
We consider now the motion of a material particle $`\tau `$ from the $`I`$-viewpoint. At the moment we restrict our consideration to classical one dimensional motions. We identify the total $`I`$-state $`q`$ of a particle $`\tau `$ with the spatial coordinate of this particle. $`q𝐙_p`$ has the form $`q=\alpha _0+\alpha _1p+\mathrm{}+\alpha _mp^m+\mathrm{}.`$ This representation can be considered as the expansion of the distance $`q`$ in the $`p`$-scale. The main difference from the real model of the motion of $`\tau `$ is discreetness of space. There is the minimal length element $`l=1.`$ The particle $`\tau `$ could not be observed on distances which are less than $`l=1.`$ Other difference is that $`q`$ can yield infinitely large values (these are $`q`$ for which $`\alpha _j0`$ for an infinite number of $`j).`$ Thus the realization of $`I`$-space as spatial space does not reproduce the ordinary model of motion in continuous real space. It gives a model of motion in discrete space. The ordinary physical interactions can realized in this space (see , , -). In this way they can be interpreted as $`I`$-interactions.
We develop an analogue of the Hamiltonian dynamics on the $`I`$-spaces <sup>5</sup><sup>5</sup>5In fact, this is an application to the $`I`$-theory of the Hamiltonian $`p`$-adic formalism developed in (and generalized in ). As usual, we introduce the quantity $`p(t)=\dot{q}(t)(=\frac{d}{dt}q(t))`$ which is the information analogue of the momentum. However, here we prefer to use a physiological terminology. The quantity $`p(t)`$ is said to be a motivation (for changing of the $`I`$-state $`q(t)`$).
The space $`𝐙_p\times 𝐙_p`$ of points $`z=(q,p)`$ where $`q`$ is the $`I`$-state and $`p`$ is the motivation is said to be a phase $`I`$-space. As in the ordinary Hamiltonian formalism, we assume that there exists a function $`H(q,p)`$ ($`I`$-Hamiltonian) on the phase $`I`$-space which determines the motion of $`\tau `$ in the phase $`I`$-space:
$$\dot{q}(t)=\frac{H}{p}(q(t),p(t)),q(t_0)=q_0,$$
(5)
$$\dot{p}(t)=\frac{H}{q}(q(t),p(t)),p(t_0)=p_0.$$
(6)
The $`I`$-Hamiltonian $`H(p,q)`$ has the meaning of an $`I`$-energy. In principle, $`I`$-energy is not related to the usual physical energy.
If $`\tau `$ is a (material) particle, then (5), (6) gives the Hamiltonian dynamics for the particle; here $`q(t)`$ is the spatial coordinate of the particle in discrete space and $`p(t)`$ is the momentum of the particle (which is also discrete). If $`\tau `$ is a cognitive system, then (5), (6) gives the Hamiltonian dynamics for the cognitive system in the ‘space of ideas’.
The simplest $`I`$-Hamiltonian $`H_{\mathrm{fr}}(p)=\alpha p^2,\alpha Z_p`$ describes the motion of a free $`I`$-transformation $`\tau `$, i.e., an $`I`$-transformer which uses only self-motivations for changing of its $`I`$-state $`q(t)`$. Here by solving the system of the Hamiltonian equations we obtain: $`p(t)=p_0,q(t)=q_0+2\alpha p_0(tt_0)`$ <sup>6</sup><sup>6</sup>6In fact, this simplest $`I`$-system is not trivial from the mathematical viewpoint. There exist other solutions which are nonanalytic (but smooth), see , . These solutions may also have an interesting $`I`$-interpretation. We shall discuss this problem later.. The motivation $`p`$ is the constant of this motion. Thus the free $`I`$-transformer ”does not like” to change its motivation $`p_0`$ in the process of the motion in the $`I`$-space. If, we change coordinates, $`q^{}=(qq_0)/k,k=2\alpha p_0`$, then we see that the dynamics of the free $`I`$-transformer coincides with the dynamics of its $`I`$-time.
If $`\tau `$ is a (material) particle, then $`p_0`$ is its momentum and $`\alpha =1/2m,`$ where $`m(=m_0+m_1p+\mathrm{}+m_lp^l)`$ is the mass of $`\tau `$ (which determined with a finite precision). If $`\tau `$ is a cognitive system, then $`p_0`$ is (internal) motivation of $`\tau `$ and $`\alpha =1/2m,`$ where $`m`$ is so called $`I`$-mass (see section 4).
In general case the $`I`$-energy is the sum of the $`I`$-energy of motivations $`H_f=\alpha p^2`$ (which is an analogue of the kinetic energy) and potential $`I`$-energy $`V(q)`$:
$$H(q,p)=\alpha p^2+V(q).$$
The potential $`V(q)`$ is determined by fields of information.
In the Hamiltonian framework we can consider interactions between $`I`$-transformers $`\tau _1,\mathrm{},\tau _N`$. These $`I`$-transformers have the $`I`$-times $`t_1,\mathrm{},t_N`$ and $`I`$-states $`q_1(t_1),\mathrm{},q_N(t_N)`$. By our model we can describe interactions between these $`I`$-transformers only in the case in that there is a possibility to choose the same $`I`$-time $`t`$ for all of them. In this case we can consider the evolution of the system of the $`I`$-transformers $`\tau _1,\mathrm{},\tau _N`$ as a trajectory in the $`I`$-space $`𝐙_p^N=𝐙_p\times \mathrm{}\times 𝐙_p,q(t)=(q_1(t),\mathrm{},q_N(t)).`$
We think that this conditions of consistency for $`I`$-times of interacting $`I`$-transformers plays the crucial role in many psychological experiments. We can not obtain sensible observations for interactions between arbitrary individuals. There must be a process of learning for the group $`\tau _1,\mathrm{}\tau _N`$ which reduces $`I`$-times $`t_1,\mathrm{},t_N`$ to the unique $`I`$-time $`t`$.
Thus, let us consider a group $`\tau _1,\mathrm{},\tau _N`$ of $`I`$-transformers with the internal time $`t`$. The dynamics of $`I`$-states and motivations is determined by the $`I`$-energy; $`H(q,p),q𝐙_p^N,p𝐙_p^N`$. It is natural to assume that
$$H(q,p)=\underset{j=1}{\overset{N}{}}\alpha _jp_j^2+V(q_1,\mathrm{},q_N),\alpha _j𝐙_p.$$
Here $`H_f(p)=_{j=1}^N\alpha _jp_j^2`$ is the total energy of motivations for the group $`\tau _1,\mathrm{},\tau _N`$ and $`V(q)`$ is the potential energy. It is natural to choose $`V(q)=_{ij}\mathrm{\Phi }(q_iq_j)`$, where $`\mathrm{\Phi }(s),s𝐙_p`$, is the potential of the interaction between $`I`$-transformers.
As usual, to find a trajectory in the phase $`I`$-space $`𝐙_p^N\times 𝐙_p^N`$, we need to solve the system of Hamiltonian equations:
$$q_j=\frac{H}{p_j},p_j=\frac{H}{q_j},q_j(t_0)=q_0,p_j(t_0)=p_0.$$
(7)
(see for such equations).
Consequences for cognitive and social sciences and psychology:
1. Energy and information. In our model a transmission of information is determined by the $`I`$-energy which is the sum of $`I`$-energy of motivations and potential $`I`$-energy. In principle, this process need no physical energy. Therefore, there might be transmissions of information which could not be reduced to transmissions of physical energy. In this case we cannot measure physical interactions (i.e., interactions in real space-time) between two $`I`$-transformers, $`\tau _1`$ and $`\tau _2`$ (but we could measure an information interaction). In particular, $`\tau _1`$ and $`\tau _2`$ can be individuals participating in psychological or social experiments (or even experiments which exhibit anomalous behaviour).
2. Distance and information. $`I`$-processes may evaluate in an $`I`$-space which differs from the real space (absolute Newton space or a space of general relativity). Therefore the real (”physical”) distance between $`I`$-transformers does not play the crucial role in processes of $`I`$-interactions.
3. Time and information. Dynamics of information is dynamics with respect to $`I`$-time $`t.`$ There may be a correspondence $`t_{\mathrm{phys}}=g(t)`$ between real time $`t_{\mathrm{phys}}𝐑`$ and $`I`$-time $`t𝐙_p.`$ This correspondence may not preserve distances.
Let $`\tau `$ be an $`I`$-transformer having a continuous trajectory $`q(t).`$ Small variations of $`t,t^{}=t+\delta t,`$ imply small variations of $`q:`$
$$a^{}=q(t^{})=a+p\delta t,a=q(t).$$
(8)
If (in some way) we find the internal time scale of $`\tau ,`$ then it would be possible to find (via (8) ) its $`I`$-state at the instant of time $`t_{\mathrm{phys}}^{}=g(t^{}).`$ If $`t_{\mathrm{phys}}^{}>t_{\mathrm{phys}}`$ then such an $`I`$-measurement can be considered as a prediction of future events; if $`t_{\mathrm{phys}}^{}<t_{\mathrm{phys}}`$ then we have recalling. The relation (8) gives only unsharp information. Thus such acts of recalling and predictions may give a lot of unfruitful information.
4. Motivation. A motion in the $`I`$-space depends, not only on the initial $`I`$-state $`q_0`$, but also on the initial motivation $`p_0`$. Moreover, the Hamiltonian structure of the equations of motion implies that the motivation $`p(t)`$ plays the important role in the process of the evolution. Thus $`I`$-dynamics is, in fact, dynamics in phase $`I`$-space.
5. Consistency for times. An $`I`$-interaction between $`I`$-transformers is possible only if these $`I`$-transformers have consistent $`I`$-times. Therefore every psychological or social experiment has to contain an element of ”learning” for $`I`$-transformers participating in the experiment. A physical interaction need not be involved in such learning. This can be any exchange of information between individuals (or a study of information about some individual).
6. Future and past. The consistency condition for $`I`$-times does not imply such a condition for real times, because different $`I`$-transformers can have different correspondence laws for $`I`$-time and real time. For example, let us consider two $`I`$-transformers, $`\tau _1`$ and $`\tau _2`$ satisfying the consistency condition for $`I`$-times, i.e., $`t_1=t_2=t`$. We assume that it is possible to transform $`I`$-times of $`\tau _1`$ and $`\tau _2`$ to real times $`t_{1,\mathrm{phys}}=g_1(t_1)`$ and $`t_{2,\mathrm{phys}}=g_2(t_2)`$. Let us also assume that $`\tau _1`$ and $`\tau _2`$ interact by the $`I`$-potential $`V(q_1q_2)`$, i.e., at the instant $`t`$ of $`I`$-time the potential $`I`$-energy of this interaction equals $`V(q_1(t)q_2(t))`$. If $`t_{1,\mathrm{phys}}=g_1(t)t_{2,\mathrm{phys}}=g_2(t)`$ then such an interaction is nothing than an interaction with the future or the past.
7. Social phenomena. By our model any social group $`G`$ can be described by a system $`\tau _1,\mathrm{},\tau _N`$ of coupled $`I`$-transformers. There exists an $`I`$-potential $`V(q_1,\mathrm{},q_N)`$ which determines an $`I`$-interaction between members of $`G`$. For example, democratic societies are characterized by uniform $`I`$-potentials $`V=\mathrm{\Phi }(q_iq_j)`$. Here a contribution into the potential $`I`$-energy does not depend on an individual. On the other hand, hierarchic societies are characterized by $`I`$-potentials of the form:
$$V=A_0\underset{j0}{}\mathrm{\Phi }(q_0,q_j)+A_1\underset{j0,1}{}\mathrm{\Phi }(q_1,q_j)+\mathrm{}$$
$$+A_k\underset{j0,\mathrm{},k}{}\mathrm{\Phi }(q_k,q_j)+B\underset{i,j0,\mathrm{},k}{}\mathrm{\Phi }(q_i,q_j),$$
where $`|A_0|_p>>|A_1|_p>>\mathrm{}>>|A_k|_p>>|B|_p`$. These potentials describe the hierarchy $`\tau _0\tau _1\mathrm{}\tau _k(\tau _{k+1},\mathrm{},\tau _N)`$. The $`I`$-transformer $`\tau _0`$ can be a political, national or state leader or a God.
Remark 4.1. (Transformers of information and classical real fields). If $`p\mathrm{}`$ then the coding alphabet $`\{0,1,\mathrm{},p1\}`$ could be thought as being continuous, i.e., it can be identified with the field of real numbers $`𝐑.`$ Therefore information space $`X=𝐙_p,p\mathrm{},`$ can be identified with the infinite product of real fields, $`X=𝐑^{\mathrm{}}.`$ Thus the $`I`$-state of an $`I`$-transformer $`\tau `$ can be identified with a classical field $`\varphi (x),x𝐑`$ (for example via Fourier coefficients). Therefore we can consider $`I`$-transformers as sources of classical fields (in the limit $`p\mathrm{}`$). Of course, this is just a speculation, because we have no mathematical realization of this limiting procedure.
## 4 Information velocity, acceleration, mass and force, Newton’s law.
We have considered dynamics of $`I`$-transformers of the unit mass. There the coefficient $`v`$ of a proportion between the variation $`\delta q`$ of the $`I`$-state and the variation $`\delta t`$ of $`I`$-time $`t`$: $`\delta q=v\delta t`$, was considered as a motivation. In the general case the motivation $`p`$ may not coincide with $`v`$. Let us assume that the motivation $`p`$ is proportional to $`v`$, $`p=mv,m𝐙_p.`$ This coefficient $`m`$ of proportion is called an $`I`$-mass. We also call $`v`$ an $`I`$-velocity. Thus $`\delta q=\frac{p}{m}\delta t`$.
Let $`\tau _1`$ and $`\tau _2`$ be two $`I`$-transformers with the $`I`$-masses $`m_1`$ and $`m_2`$ and let $`|m_1|_p>|m_2|_p`$. Let $`\tau _1`$ and $`\tau _2`$ have the variations $`\delta t_1`$, $`\delta t_2`$ of $`I`$-time of the same $`p`$-adic magnitude, $`|\delta t_1|_p=|\delta t_2|_p`$, and let these variations generate the variations $`\delta q_1`$ and $`\delta q_2`$ of their $`I`$-states of the same $`p`$-adic magnitude, $`|\delta q_1|_p=|\delta q_2|_p`$. To make such a change of the $`I`$-state, $`\tau _1`$ need a larger motivation: $`|p_1|_p=|\frac{\delta q}{\delta t}|_p|m_1|_p>|p_2|_p=|\frac{\delta q}{\delta t}|_p|m_2|_p`$. Thus the $`I`$-mass is a measure of an inertia of information. We define a kinetic $`I`$-energy by $`T=\frac{1}{2m}p^2`$.
A variation $`\delta t`$ of $`I`$-time $`t`$ implies also a variation $`\delta p`$ of the motivation $`p`$: $`\delta p=f\delta t.`$ The coefficient $`f`$ of proportionality is called an $`I`$-force. Thus any change of the motivation is due to the action of an $`I`$-force $`f`$. If $`f=0`$ then $`\delta p=0`$ for any variation $`\delta t`$ of $`t`$. Thus an $`I`$-transformer cannot change its motivation in the absence of $`I`$-forces.
By analogue with the usual physics we call the coefficient $`a`$ of a proportion between the variation $`\delta v`$ of the $`I`$-velocity $`v`$ and the variation $`\delta t`$ of the $`I`$-time $`t`$, $`\delta v=a\delta t`$, an $`I`$-acceleration. Thus $`\delta p=am\delta t`$. This relation can be rewritten in the form of an information analogue of the second Newton law:
$$ma=f$$
(9)
or
$$\dot{p}=f.$$
(10)
An $`I`$-force $`f`$ is said to be a potential force if there exists a function $`V(q)`$ such that $`f=\frac{V}{q}`$ where $`V`$ is called the potential, or potential energy. The total $`I`$-energy $`H`$ is defined as the sum of the kinetic and the potential $`I`$-energies, $`H(q,p)=\frac{1}{2m}p^2+V(q)`$. The Hamiltonian equation $`\dot{p}=\frac{H}{q}`$ coincides with the Newton equation $`\dot{p}=f`$.
Example 4.1. (Hooke’s $`I`$-system). Let the $`I`$-force $`f`$ be proportional to the $`I`$-state $`q`$, $`f=m\beta ^2q`$, where $`m`$ is the $`I`$-mass and $`\beta 𝐙_p`$ is a coefficient of the interaction. Here (9) gives the equation $`\ddot{q}=\beta ^2q`$. As $`f=\frac{V}{q}`$, $`V(q)=\frac{m\beta ^2}{2}q^2`$ and $`H(q,p)=\frac{p^2}{2m}\frac{m\beta ^2q^2}{2}`$; the Hamiltonian equations are $`\dot{q}=p/m`$ and $`\dot{p}=m\beta ^2q`$. Their solutions have the form $`g(t)=ae^{\beta t}+be^{\beta t}`$. By the condition (3) the $`I`$-state $`q(t)`$ and motivation $`p(t)`$ are defined only for instants of $`I`$-time which satisfy the inequality
$$|\beta t|_pr_p.$$
(11)
This condition can be considered as a restriction for the magnitude of the $`I`$-force. If the coefficient of the interaction $`|\beta |_pr_p,`$ then dynamics $`q(t)`$ of the $`I`$-state is well defined for all $`t𝐙_p.`$ Larger forces imply the restriction condition for $`I`$-time. Let $`|\beta |_p=1`$. If $`p2`$ then (11) has the form $`tU_{1/p}(0)`$, i.e., $`t=\alpha _1p+\alpha _2p^2+\mathrm{}`$. Thus the $`I`$-state $`q(t)`$ of the $`I`$-transformer $`\tau `$ can be defined (observed) only for the instants of time $`t_0=0,t_1=p,\mathrm{},t_{p1}=(p1)p,\mathrm{}.`$ If $`p=2`$ then (11) has the form $`tU_{1/4}(0)`$, i.e., and $`t=\alpha _22^2+\alpha _32^3+\mathrm{}`$. Thus the $`I`$-state $`q(t)`$ of $`\tau `$ can be defined (observed) only for the instants of time $`t_0=0,t_1=4,t_2=8,\mathrm{}.`$
Let $`f=m\beta ^2q`$, i.e., $`V(q)=\frac{m\beta ^2q^2}{2}`$ and $`\ddot{q}=\beta ^2q`$. Here $`q(t)`$ and $`p(t)`$ have the form $`g(t)=a\mathrm{cos}\beta t+b\mathrm{sin}\beta t`$. Here we also have the restriction relation (11). As opposite to the real case the $`p`$-adic trigonometric functions are not periodical. There is no analogue of oscillations for the $`I`$-process described by an analogue of Hooke’s law.
Let us consider the solution of the Hamiltonian equations with the initial conditions $`q(0)=0`$ and $`p(0)=m\beta `$: $`q(t)=\mathrm{sin}\beta t`$, $`p(t)=m\beta \mathrm{cos}\beta t`$. We have $`qp=(m\beta /2)\mathrm{sin}2\beta t`$. By using the $`p`$-adic equality $`|\mathrm{sin}a|_p=|a|_p`$ we get $`|qp|_p=|m\beta |_p|\beta t|_p`$. The relation (11) implies
$$|q|_p|p|_p|m\beta |_pr_p.$$
(12)
This is a restriction relation for the trajectory $`(q(t),p(t))`$ in the phase $`I`$-space (compare with ). Let $`\beta =1/m`$. Then (12) gives $`|q|_p|p|_pr_p`$. If the motivation $`p`$ is strong $`|p|_p=1`$, then $`q`$ can be only of the form $`q=\alpha _1p+\alpha _2p^2+\mathrm{}`$, $`p2`$ and $`q=\alpha _22^2+\alpha _32^3+\mathrm{}`$, $`p=2`$. If the motivation $`p`$ is rather weak then the $`I`$-state $`q`$ of an $`I`$-transformer can be arbitrary.
The restriction relation (12) is natural if we apply our information model to describe psychological (social) behaviour of individuals. Strong psychological (social) motivations imply some restrictions for possible psychological (social) states $`q`$. On the other hand, if motivations are rather weak an individual can, in principle, arrive to any psychological (social) state.
We discuss the role of the $`I`$-mass in the restriction relation (12). There the decrease of the $`I`$-mass implies more rigid restrictions for the possible $`I`$-states (for the fixed magnitude of the motivation). If we return to the psychological (social) applications we get that the individual (or a group of individuals) with a small magnitude of $`I`$-mass and the strong motivations will have quite restricted set of $`I`$-states.
The restriction relation (12) is an analogue of the Heisenberg uncertainty relations in the ordinary quantum mechanics. However, we consider a classical (i.e., not quantized) $`I`$-system. Therefore a classical $`I`$-system can have behaviour that is similar to quantum behaviour.
## 5 Mathematical ”pathologies” in the formalism of the information mechanics and their interpretations
In $`p`$-adic analysis the condition $`f0`$ does not imply that a differentiable function $`f`$ is a constant, see , . Therefore, there exist very complicated continuous motions $`(q(t),p(t))`$ in the $`I`$-phase space for $`I`$-transformers with zero $`I`$-energy ($`\dot{q}0`$ or $`\dot{p}0`$).
In psychological models these motions can be interpreted as motions without any motivation. Such motions need no information force. On the other hand, we can consider an $`I`$-potential $`V(q)`$ such that $`\frac{V}{q}=0`$. Here the potential $`I`$-energy $`V(q)`$ can have very complicated behaviour on the $`I`$-space $`X=𝐙_p`$. At the same time the $`I`$-force $`f=0`$. Thus there may exist $`I`$-fields which do not induce any $`I`$-force.
All mathematical pathologies can be eliminated by the consideration of analytical functions. If $`f0`$ and $`f`$ is analytic then $`f=\text{constant}`$.
In psychological models we can interpret analytical trajectories in the phase $`I`$-space as a ”normal behaviour”, i.e., an individual need a motivation for the change of a psychological state. Here we can observe some psychological (information) force which induces this change. There is a psychological (information) field which generates this force. The model puts trajectories (non-analytical) with zero motivation in relation with abnormal psychological behaviour, mental diseases and anomalous phenomena. Here an individual changes his psychological state without any motivation in the absence of any information force. Here, in fact, a $`p`$-adic generalization of the Hamiltonian formalism does not work. We need to propose a new physical formalism to describe such phenomena.
Not all unusual properties of $`p`$-adic quantities are connected with non-analyticity. For example, in $`p`$-adic analysis we can construct polynomials of the form $`V(x)=\alpha _0+\alpha _1x+\mathrm{}+\alpha _Nx^N,`$ where the coefficients $`\alpha _j`$ are natural numbers, $`|\alpha _j|_p=1`$, such that $`ϵ=sup_{x𝐙_p}|V(x)|_p`$ can be arbitrary small (see ). Therefore the result of the simultaneous action of quite strong $`I`$-potentials $`V_j(x)=\alpha _jx^j,j=0,1,\mathrm{},N,`$ can have arbitrary small magnitude.
## 6 Information work, conservation laws
To eliminate from our consideration all ”pathological” motions in the $`I`$-space, we shall consider only $`I`$-quantities described by analytical functions. Of course, we do not claim that only analytical functions describe real information processes. We like only to simplify mathematical considerations.
Let $`f(x)=_{n=0}^{\mathrm{}}a_nx^n`$, $`a_n𝐐_p`$, and let the series converge for $`|x|_p\delta `$, $`\delta =p^{\pm n},n=0,1,\mathrm{}`$. We define an integral of $`f`$ by the formula (see ):
$$_a^bf(x)𝑑x=\underset{n=0}{\overset{\mathrm{}}{}}\frac{a_n}{n+1}[a^{n+1}b^{n+1}].$$
The series on the right-hand side converges for all $`|a|_p,|b|_p\frac{\delta }{p}`$. In particular, we can find an antiderivative $`F`$ of $`f`$ by the formula $`F(x)=_0^xf(x)𝑑x`$.
Let $`f`$ be an $`I`$-force which is described by the function $`f(x)`$ which is analytic for $`|x|_pp`$. Then this force is potential with the $`I`$-potential $`V(x)=_0^xf(x)𝑑x`$.
Let $`\gamma =\{q(t),|t|_p\lambda \}`$ be an analytic curve in $`𝐙_p`$. We define its length element by $`ds=vdt`$, where $`v=\dot{q}`$ is the $`I`$-velocity. By definition
$$W_{ab}=_{\gamma (a,b)}f𝑑s=_{t_0}^{t_1}f(q(t))v(t)𝑑t$$
where $`q(t_0)=a`$ and $`q(t_1)=b`$. The quantity $`W_{ab}`$ is said to be the work done by the external $`I`$-force $`f`$ upon the $`I`$-transformer in going from the point $`a`$ to the point $`b`$. By (10) we have
$$W_{ab}=_{t_0}^{t_1}m\dot{v}v𝑑t=\frac{m}{2}_{t_0}^{t_1}\frac{d}{dt}v^2𝑑t=\frac{1}{2m}(p^2(b)p^2(a)).$$
Thus the work done is equal to the change in the kinetic energy: $`W_{ab}=T_bT_a`$. As the $`I`$-force $`f`$ is potential then the work $`W`$ done around a closed orbit is zero: $`W=f𝑑s=0`$. Thus the work $`W_{ab}`$ does not depend on an analytic trajectory $`\gamma (a,b)`$.
We also have:
$$W_{ab}=_{\gamma (a,b)}\frac{V}{q}ds=_{\gamma (a,b)}\frac{d}{dt}V(q(t))dt=V(a)V(b).$$
Thus $`T_bT_a=V(a)V(b)`$. We have obtained the energy conservation law for an $`I`$-transformer: if the $`I`$-forces acting on an $`I`$-transformer are described by analytical functions (in particular, they are potential), then the total energy of the $`I`$-transformer, $`H=T+V`$, is conserved.
At the moment the situation with nonanalytic potential $`I`$-forces is not clear. It may be that the energy conservation law is violated in the general case.
## 7 Mechanics of a system of information transformers, constraints on information spaces
Let $`\tau _1,\mathrm{},\tau _N`$ be a system of $`I`$-transformers with $`I`$-masses, $`m_1,\mathrm{},m_N𝐙_p`$. As in ordinary mechanics we must distinguish between the external $`I`$-forces $`F_i^{(e)}`$ acting on $`I`$-transformers due to sources outside the system and internal forces $`F_{ji}`$. As we have already discussed, $`I`$-times $`t_1,\mathrm{},t_N`$ of $`\tau _1,\mathrm{},\tau _N`$ must satisfy the consistency condition:
$$t_1=t_2=\mathrm{}=t_N=t.$$
(13)
Thus the equation of motion for the $`i`$th particle is to be written:
$$\dot{p}_i=F_i^{(e)}+\underset{j}{}F_{ji}.$$
(14)
For some $`I`$-systems we may obey an information analogue of Newton’s third law (a law of information action and reaction): $`F_{ij}=F_{ji}`$.
Set $`x=_im_ix_i/M`$, where $`M=m_i`$. This point in the $`I`$-space is said to be the center of information of the system. If the system satisfies Newton’s third law for $`I`$-forces then we get the equation of motion: $`M\ddot{x}=_iF_i^{(e)}=F^{(e)}`$. The center of information moves as if the total external $`I`$-force was acting on the $`I`$-mass $`M`$ of the system concentrated at the center of information. We introduce the motivation $`P=M\dot{x}`$ of the $`I`$-system. There is the following conservation theorem for motions described by analytic functions $`(q_j(t))_{j=1}^N,t𝐙_p`$: if the total external $`I`$-force is zero, the total motivation of the $`I`$-system is conserved.
Example 7.1. (Social systems). We apply our $`I`$-model for describing a society $`S`$ which consists of individuals (or groups of individuals) $`\tau _1,\mathrm{},\tau _N.`$ There exist the center of information of $`S`$, $`x_S𝐙_p`$ which can be considered as a coding sequence for this society. If $`S`$ satisfies Newton’s law of action-reaction for $`I`$-forces then its evolution is determined by the external $`I`$-forces. If this evolution is not ”pathological” then the motivation of $`S`$ is conserved. Of course, there might be numerous ”pathological” evolutions (for example, evolutions with zero motivation, $`P_S=0`$).
For analytic motions the $`I`$-work done by all $`I`$-forces in moving the system from an initial configuration $`A=\{a_i=q_i(t_0)\}`$ to a final configuration $`B=\{b_i=q_i(t_1)\}`$ is well defined:
$$W_{ab}=\underset{i}{}_{\gamma (a_i,b_i)}F_i𝑑s_i+\underset{ij}{}_{\gamma (a,b)}F_{ji}𝑑s_i$$
and $`W_{ab}=T_BT_A`$, where $`T=\frac{1}{2}_im_iv_i^2`$ is the total kinetic $`I`$-energy of the $`I`$-system. As usual $`T=\frac{1}{2}Mv^2+\frac{1}{2}_im_iv_i^2`$ where $`v`$ is the velocity of the center of information and $`v_i^{}`$ is the velocity of $`\tau _i`$ with respect to the center of information.
In our model of ‘social motion’ (Example 7.1) we can say that the total kinetic energy of the society $`S`$ is the sum of the kinetic energy of the center of information of $`S`$ and the kinetic energy of motions of individuals $`\tau _j`$ about the center of information.
We now consider the case when all $`I`$-forces are (analytical) potential: $`F_i^{(e)}=\frac{V_i}{x_i}`$ and $`F_{ji}=\frac{V_{ij}}{x_i}`$. To satisfy the law of action and reaction we can choose $`V_{ij}=\mathrm{\Phi }_{ij}(x_ix_j)`$ where $`\mathrm{\Phi }_{ij}:𝐙_p𝐙_p,\mathrm{\Phi }_{ij}=\mathrm{\Phi }_{ji}`$ are analytical functions. Then by repeating the considerations of the standard mechanics over the reals we obtain that $`W_{AB}=V(B)+V(A)`$, where $`V=_iV_i+\frac{1}{2}_{i,j}V_{ij}`$ is the total potential energy of the system of $`I`$-transformers. Therefore the total $`I`$-energy $`H=T+V`$ is conserved for every $`I`$-system with (analytical) potential $`I`$-forces (such that $`F_{ij}`$ satisfy the law of information action-reaction).
The consideration of $`I`$-systems induces dynamics in multidimensional $`I`$-spaces; $`X_N=𝐙_p^N`$. Such spaces can be useful for the description, not only systems of $`I`$-transformers, but also individual $`I`$-transformers which have multidimensional information spaces.
For example, let $`\tau `$ be a cognitive system and let $`x=(x_1,\mathrm{},x_N),x_j𝐙_p`$, be a set of ideas with which operates $`\tau `$ (i.e., there are $`N`$ parallel thinking processes $`\pi _1,\mathrm{},\pi _N`$ in $`\tau `$, see - for the details). Then the $`I`$-dynamics for $`\tau `$ is described by the trajectory $`(q(t),p(t))𝐙_p^{2N}.`$
As in standard mechanics, constraints play the important role in $`I`$-mechanics. The simplest constraints (”holonomic”) can be expressed as equations connecting $`I`$-states of $`I`$-transformers $`\tau _1,\mathrm{},\tau _N`$ (or equations coupling different ideas in the cognitive system):
$$f(q_1,\mathrm{},q_N,t)=0.$$
Here $`f`$ may be a function from $`𝐙_p^{N+1}`$ into $`𝐙_p`$ or a function from $`𝐙_p^{N+1}`$ into $`𝐑`$. The simplest constraints of the ”real type” are:
($`C1`$) $`|q_1a|_p=r,\mathrm{},|q_Na|_p=r,r>0,a𝐙_p`$, i.e., all $`I`$-transformers have to move over the surface of the sphere $`S_r(a)`$;
($`C2`$) $`|q_2q_1|=r,\mathrm{},|q_Nq_1|=r`$, i.e., there is the fixed $`I`$-transformer $`\tau _1`$ such that all other $`I`$-transformers must move on the distance $`r`$ from $`\tau _1;`$
($`C3`$) We can also consider an ”information rigid body”, i.e., a system of $`I`$-transformers connected by constraints: $`|q_iq_j|_p=r_{ij}.`$
Example 7.2. (Restricted mentality). In cognitive sciences constraint ($`C1`$) can be used for the description of a ”restricted mentality”. All ideas $`q_1(t),\mathrm{},q_N(t)`$ of a cognitive system $`\tau `$ (generated by the parallel processes $`\pi _1,\mathrm{},\pi _N`$) belong to the restricted domain of ideas $`X=S_r(a).`$
Example 7.3. (Ideology, religion). Let us consider the $`I`$-model of a society $`S`$ with an ideology (or religion) $`a𝐙_p.`$ Then constraint ($`C1`$) can be interpreted as describing a social layer $`=(\tau _1,\mathrm{},\tau _N)`$ of $`S.`$ These are all individuals who accept the ideology (or religion) $`a`$ with an ”information precision” $`r.`$ Let this precision $`r=1/p^k`$ and let $`q_j(t)=(q_{j\alpha }(t))_{\alpha =0}^{\mathrm{}},a=(a_\alpha )_{\alpha =0}^{\mathrm{}}.`$ The constraint ($`C1`$) implies that
$$q_{j0}(t)=a_0,\mathrm{},q_{jk1}(t)=a_{k1},\text{but}q_{jk}(t)a_k.$$
The members of $``$ accept dogmas $`a_0,\mathrm{},a_{k1}`$ of the ideology (or religion), but they deny the dogma $`a_k.`$ In our hierarchical model all other dogmas do not play any role. If the dogma $`a_k`$ is violated then the violation of $`a_{k+j}`$ would not change a status of $`\tau _j.`$
Example 7.4. (Evolution of an idea-fix). Let us consider a cognitive system $`\tau `$ with $`N`$ parallel thinking processes $`\pi _1,\mathrm{},\pi _N.`$ The constraint ($`C2`$) means that there is a thinking process in the cognitive system (in our case this is $`\pi _1`$) which has a strong influence on all other thinking processes $`\pi _j,j1.`$ They could not go far away from $`\pi _1.`$ In psychology $`\pi _1`$ may be interpreted as a process of evolution of an idea-fix. The constraint ($`C2`$) in the $`I`$-space of the cognitive system $`\tau `$ implies that all thinking activity of $`\tau `$ is connected with this idea-fix.
Example 7.5. (Kingdoms, families and lovers). The constraint ($`C2`$) can be interpreted as describing a social layer $`=(\tau _2,\mathrm{},\tau _N)`$ in a kingdom $`K`$ with the king $`\tau _1.`$ The evolution $`q_1(t)`$ of the $`I`$-state of the king induces the information restrictions ($`C2`$) for evolutions of $`I`$-states of members of the layer $`.`$ The same constraint may be used for an information model of evolution of a family $`F.`$ Here $`\tau _1`$ may be the father or mother. In the case $`N=2`$ we obtain the symmetric model which may be used for the description of a pair of lovers. Similar constraints in the $`I`$-space might explain some anomalous information connections between individuals.
Example 7.6. (Scandinavian society). The constraint ($`C3`$) may be used in social sciences for the description of ”Scandinavian societies”. There are nonzero distances $`r_{ij}>0,i,j=1,\mathrm{},N,`$ between individuals in the $`I`$-space. These distances are stable in the process of time evolution.
In the case of holonomic constraints described by the system of analytical functions: $`f_j:𝐙_p^{N+1}𝐙_p,j=1,\mathrm{},K`$, i.e., $`f_j(q_1,\mathrm{},q_N,t)=0`$, we can use the technique of the standard mechanics <sup>7</sup><sup>7</sup>7These methods may not be applied to constraints determined by real valued functions. However, in the latter case we need not eliminate these constraints. These constraints describe open subsets of the configuration $`I`$-space $`𝐙_p^N`$. We can choose such subsets as new configuration $`I`$-spaces.. If the equations are independent then we can introduce generalized $`I`$-coordinates $`\xi _1,\mathrm{},\xi _{NK}`$, and $`q_l=q_l(\xi _1,\mathrm{},\xi _{NK},t)`$ $`l=1,\mathrm{},N`$, and $`q_l(\xi ,t)`$ are analytical functions of $`\xi `$ and $`t`$ (see , for the mathematical details).
Example 7.7. (Hidden basic ideas). If $`q(t)=(q_l(t))_{l=1}^N`$ describes ideas in the cognitive system at the instant $`t`$ of $`I`$-time, then by resolving constraints on these ideas we can find ”independent ideas” $`\xi (t)=(\xi _j(t))_{j=1}^{NK}`$ which, in fact, determine the $`I`$-state of the cognitive system.
Example 7.8. (Hidden leaders). If $`q(t)=(q_l(t))_{l=1}^N`$ describes the system $`S=(\tau _1,\mathrm{},\tau _N)`$ of $`I`$-transformers then the existence of generalized $`I`$-coordinates $`\xi `$ can be interpreted as a possibility to reduce $`I`$-behaviour of $`S`$ to $`I`$-behaviour of the other system $`G=(g_1,\mathrm{},g_{NK})`$ of $`I`$-transformers.
As in the standard mechanics, we introduce general $`I`$-forces:
$$Q_j=\underset{i}{}F_i\frac{q_i}{\xi _j},$$
(15)
where $`F_i`$ is the total $`I`$-force acting to $`i`$th $`I`$-transformer (i.e., $`F_i=F_i^{(a)}+f_i`$ is the sum of applied $`I`$-force $`F_i^{(a)}`$ and the $`I`$-force $`f_i`$ of constraints <sup>8</sup><sup>8</sup>8We can interpret $`I`$-constraints as unknown $`I`$-forces.).
In our theory generalized $`I`$-forces have the natural interpretation (compare with the situation with generalized forces in the usual mechanics). As we have noted, the existence of generalized $`I`$-coordinates which are obtained from equations for constraints means that the initial system $`S=(\tau _1,\mathrm{},\tau _N)`$ of $`I`$-transformers is ”controlled” by the other system $`G=(g_1,\mathrm{},g_{NK})`$ of $`I`$-transformers. The $`I`$-forces (15) are, in fact, reaction $`I`$-forces, i.e., the control of $`G`$ over $`S`$ generates $`I`$-forces applied to elements of $`G`$. By repeating of the usual computations we get the equations of motion:
$$\frac{d}{dt}(\frac{T}{\dot{\xi }_j})\frac{T}{\xi _j}=Q_j,j=1,\mathrm{},NK.$$
(16)
If the $`I`$-forces $`F_i`$ are potential with the analytical potential $`V`$, i.e., $`F_i=\frac{V}{q_i}`$, then generalized $`I`$-forces are also potential: $`Q_j=\frac{V}{\xi _j}`$. In this case the above equation can be written in the form:
$$\frac{d}{dt}(\frac{L}{\dot{\xi }_j})\frac{L}{\xi _j}=0,$$
(17)
where $`L=TV`$ is the $`I`$-Lagrangian.
It is important that the equations (17) can be used to describe $`I`$-motions in the presence of an $`I`$-potential $`V(q,\dot{q})`$ which depends on (generalized) $`I`$-velocities $`v_i=\dot{q}_i`$. In this case
$$Q_j=\frac{V}{q_j}+\frac{d}{dt}(\frac{V}{\dot{q}_j})$$
(18)
and $`L=TV`$.
These velocity-dependent potentials may play an important role in $`I`$-processes. In particular, there might have applications in such an exotic field as anomalous phenomena. It is claimed (see, for example, - ) that a psychokinesis effect can be observed for some random physical processes and it cannot be observed for deterministic processes. It might be tempting to explain this phenomenon on the basis of the assumption that an $`I`$-field generated in experiments on the psychokinesis corresponds to a potential which depends on the $`I`$-velocity (thus the corresponding $`I`$-force is defined by (18) ). The $`I`$-velocity is higher for random processes. Therefore the interaction is stronger for these processes.
## 8 Quantum mechanics on information spaces
It is quite natural to quantize classical mechanics on information spaces over $`𝐙_p`$. We give the following reasons for such quantization. Observations over $`I`$-quantities are statistical observations. We have to study statistical ensembles of $`I`$-transformers (instead studying of an individual $`I`$-transformer). Such statistical ensembles are described by quantum states $`\varphi `$. As usual in quantum formalism, we can assume that a value $`\lambda `$ of an $`I`$-quantity $`A`$ can be measured in the state $`\varphi `$ with some probability $`𝐏_\varphi (A=\lambda )`$. This ideology is nothing than the application of the statistical (ensemble) interpretation of quantum mechanics (see, for example, or ) to the information theory. By this interpretation any measurement process has two steps: (1) a preparation procedure $``$; (2) a measurement of a quantity $`B`$ in the states $`\varphi `$ which were prepared with the aid of $``$.
Let us consider these steps in the information framework. By $``$ we have to select a statistical ensemble $`\varphi `$ of $`I`$-transformers on the basis of some $`I`$-characteristics. Typically in quantum physics a preparation procedure $``$ is realized as a filter based on some physical quantity $`A`$, i.e., we select elements which satisfy the condition $`A=\mu `$ where $`\mu `$ is one of the values of $`A`$. We can do the same in quantum $`I`$-theory. An $`I`$-quantity $`A`$ is chosen as a filter, i.e., $`I`$-transformers for the statistical ensemble $`\varphi `$ are selected by the condition $`A=\mu `$ where $`\mu 𝐙_p`$ is some information. For example, we can choose $`A=p`$, the motivation, and select a statistical ensemble $`\varphi =\varphi (p=\mu )`$ of $`I`$-transformers which have the same motivation $`\mu 𝐙_p`$. Then we realize the second step of a measurement process and measure some information quantity $`B`$ in the state $`\varphi _{(p=\mu )}`$. For example, we can measure the $`I`$-state $`q`$ of $`I`$-transformers belonging to the statistical ensemble described by $`\varphi _{(p=\mu )}`$. We shall obtain a probability distribution $`𝐏(q=\lambda |p=\mu ),\lambda ,\mu 𝐙_p`$ (a probability that $`I`$-transformer has the $`I`$-state $`q=\lambda `$ under the condition that it has the motivation $`p=\mu `$). It is also possible to measure the $`I`$-energy $`E`$ of $`I`$-transformers. We shall obtain a probability distribution $`𝐏(E=\lambda |p=\mu ),\lambda ,\mu 𝐙_p`$<sup>9</sup><sup>9</sup>9 We now try to provide theoretical foundations for quantum $`I`$-theory. We do not discuss concrete measurement procedures for $`I`$-quantities. In particular, at the moment it is not clear how the $`I`$-energy can be measured. It seems natural to use an analogue with usual quantum theory here. The $`I`$-energy can be measured in the process of interactions between $`I`$-transformers or interactions of $`I`$-transformers and $`I`$-fields.. On the other hand, we can prepare a statistical ensemble $`\varphi _{(q=\mu )}`$ by fixing some information $`\mu 𝐙_p`$ and selecting all $`I`$-transformers which have the $`I`$-state $`q=\mu `$. Then we can measure motivations of these $`I`$-transformers and we shall obtain a probability distribution $`𝐏(p=\lambda |q=\mu )`$.
Other possibility is to use a generalization of the individual interpretation of quantum mechanics. By this interpretation a wave function $`\psi (x),x𝐑^n,`$ describes the state of an individual quantum particle. In the same way we may assume that a wave function $`\psi (x),x𝐙_p^n,`$ on the $`I`$-space describes the state of an individual $`I`$-transformer $`\tau .`$
Example 8.1. (Referendum). In some social models we can consider individuals as quantum $`I`$-transformers. A referendum is one of the possible measurement devices. Here the act of a measurement is a procedure of giving answers to questions of the referendum. By the individual interpretation individuals have no definite answers to these questions before the referendum. These answers (information communications) are created in the process of the referendum. In fact, this $`I`$-measurement changes $`I`$-states of individuals.
Example 8.2. (Conscious measurement of quantum subconsciousness) We might describe brain’s functioning by the following quantum $`I`$-model. There is a quantum system, subconsciousness<sup>10</sup><sup>10</sup>10By our model the subconsciousness is a kind of processor in that work a large number of dynamical systems of the form $`x_n=f(x_{n1}),`$ where $`f:𝐙_p𝐙_p`$ is a continuous function. Attractors of these dynamical systems are solutions of problems., which state is described by the wave function $`\psi (x),x𝐙_p.`$ There is a measurement device, consciousness, which measures the $`I`$-state $`q`$ of the subconsciousness. The concrete value (idea) of $`q`$ is not determined before the act of the conscious measurement. It is created only at the instant of a measurement. Of course, this act of a measurement (as in the ordinary quantum mechanics) changes the state of the subconsciousness. The main difference from the standard quantum mechanical scheme is that we consider repeatable measurements over the same quantum system. In ordinary scheme of a quantum measurement we consider an ensemble of identical systems. At the moment we can present only some speculations about nature of the consciousness. The consciousness is an information field generated by the brain. This field interacts continuously with the subconsciousness<sup>11</sup><sup>11</sup>11This is a feedback process : the conscious field sends to the subconsciousness a problem $`x_0`$ which is the initial condition for one of dynamical systems located in the subconsciousness (this is the signal to start the work of the dynamical system). On the other hand, an attractor of this dynamical system (a solution of the problem $`x_0`$) interacts with the conscious field (this is the signal to stop the work of the dynamical system)..
Example 8.3. ( Psychoanalysis). On the basis of the model of the previous example we can interpret psychoanalysis as a series of measurements of the $`I`$-state of the subconsciousness. These measurements continuously change the wave function of the subconsciousness. Thus psychoanalysis is a treatment based on the series of quantum $`I`$-measurements. In fact, psychoanalytic tries to provide some functions of the conscious field.<sup>12</sup><sup>12</sup>12Thus Freud’s theory may be interpreted in the following way. If the interaction between the consciousness and subconsciousness is not sufficiently strong, the consciousness ”cannot see” some attractors of dynamical systems located in the subconsciousness. Thus the consciousness cannot send to the subconsciousness the signal to stop iterations of these dynamical systems. These dynamical systems are continuously busy and they cannot be used for other purposes. Other possibility is that the general interaction is strong. However, the consciousness ”cannot recognize” some attractor as a solution of a problem because a strong external information field (a taboo) might hinder to the interaction. Therefore a psychoanalytic has to find the hidden attractor and by this act the work of the corresponding dynamical system will be stoped. Of course, he must be isolated from the corresponding ”taboo-field”..
The problem of interpretations is an important problem of ordinary quantum mechanics on real space. The same problem arises immediately in our quantum $`I`$-theory. We do not like to start our investigation with a hard discussion on the right interpretation. We can be quite pragmatic and use both interpretations by our convenience. However, the reader, who is interested in foundations of quantum mechanics, can find the extended discussion on the problem of the interpretation in Appendix 2.
In fact, a mathematical model for quantum $`I`$-formalism has been already constructed. This is quantum mechanics with $`p`$-adic valued functions, see , , , , . We present briefly this model. The space of quantum states is realized as a $`p`$-adic Hilbert space $`𝒦`$ (see , about the theory of such spaces). This is a $`𝐐_𝐩`$-linear space which is a Banach space (with the norm $``$) and on which is defined a symmetric bilinear form $`(,):𝒦\times 𝒦𝐐_𝐩`$. This form is called an inner product on $`𝒦`$. It is assumed that the norm and the inner product are connected by the Cauchy-Bunaykovski-Schwarz inequality: $`|(x,y)|_pxy,x,y𝒦.`$
Remark 8.1 It is possible to use more general spaces over different extensions of $`𝐐_𝐩`$ (analogues of complex Hilbert spaces).
By definition quantum $`I`$-state $`\varphi `$ is an element of $`𝒦`$ such that $`(\varphi ,\varphi )=1`$; quantum $`I`$-quantity $`A`$ is a symmetric bounded operator $`A:𝒦𝒦`$, i.e., $`(Ax,y)=(x,Ay),x,y𝒦`$<sup>13</sup><sup>13</sup>13In $`p`$-adic models we do not need to consider unbounded operators, because all quantum quantities can be realized by bounded operators, see , , , , .. We discuss a statistical interpretation of quantum states in the case of a discrete spectrum of $`A`$.
Let $`\{\lambda _1,\mathrm{},\lambda _n,\mathrm{}\},\lambda _j𝐙_p`$ be eigenvalues of $`A`$, $`A\varphi _n=\lambda _n\varphi _n,\varphi _n𝒦,(\varphi _n,\varphi _n)=1`$. The eigenstates $`\varphi _n`$ of $`A`$ are considered as pure quantum $`I`$-states for $`A`$, i.e., if the system of $`I`$-transformers is described by the state $`\varphi _n`$ then the $`I`$-quantity $`A`$ has the value $`\lambda _n𝐙_p`$ with probability 1. Let us consider a mixed state
$$\varphi =\underset{n=1}{\overset{\mathrm{}}{}}q_n\varphi _n,q_n𝐐_p,$$
(19)
where $`(\varphi ,\varphi )=_{n=1}^{\mathrm{}}q_n^2=1`$<sup>14</sup><sup>14</sup>14As in the usual theory of Hilbert spaces, eigenvectors corresponding to different eigenvalues of a symmetric operator are orthogonal.. By the statistical interpretation of $`\varphi `$ if we realize a measurement of the $`I`$-quantity $`A`$ for $`I`$-transformers belonging to the statistical ensemble described by $`\varphi `$ then we obtain the value $`\lambda _n`$ with probability $`P(A=\lambda _n|\varphi )=q_n^2`$.
The main problem (or the advantage?) of this quantum model is that these probabilities belong to the field of $`p`$-adic numbers $`𝐐_p`$. The simplest way is to eliminate this problem by considering only finite mixtures (19) for which $`q_n𝐐_p`$ (the field of rational numbers $`𝐐`$ is a subfield of $`𝐐_p`$). In this case the quantities $`𝐏(A=\lambda _n|\varphi )=q_n^2`$ can be interpreted as usual probabilities (for example, in the framework of Kolmogorov’s theory ). Therefore we may assume that there exist (can be prepared) quantum $`I`$-states $`\varphi `$ which have the standard statistical interpretation: when the number $`N`$ of experiments tends to infinity, the frequency $`\nu _N(A=\lambda _n|\varphi )`$ of an observation of the information $`\lambda _n𝐙_p`$ tends to the probability $`q_n^2`$.
However, we can use a more general viewpoint to this problem. In book a (non-Kolmogorov) probability model with $`p`$-adic probabilities has been developed. If we use a $`p`$-adic generalization of a frequency approach to probability (see R. von Mises, ), then $`p`$-adic probabilities are defined as limits of relative frequencies $`\nu _N`$ with respect to the $`p`$-adic topology <sup>15</sup><sup>15</sup>15It is quite surprising that in the $`p`$-adic framework we can obtain negative rational frequency probabilities , . On the other hand, negative ‘probabilities’ appear in quite natural way in many quantum models (see, for example, - ). P.A.M. Dirac was the first to introduce explicitly the concept of negative probability (in close connection with the concept of negative energy), . R. Feynman also discussed the possibility to use negative probabilities in quantum formalism, see . In particular, he remarked: ”The only difference between a probabilistic classical world and the equations of the quantum world is that somehow or other it appears as if the probabilities would have to go negative, and that we do not know, as far as I know, how to simulate”. These probabilities were used to explain violations of Bell’s inequality (see review ). Here the assumption that a distribution of hidden variables may be a signed ‘probabilistic measure’ implies existence of numerous models with hidden variables in that Bell’s inequality is violated. Wiegner’s distribution on the phase space gives other example of signed quantum ‘probabilistic’ distribution. In works - , , the $`p`$-adic probabilities (which are well defined on the mathematical level of rigorousness) were used to justify the use of negative probabilities in quantum theories.. The relative frequencies $`\nu _N𝐐`$ and they can be considered, not only as elements of $`𝐑`$, $`𝐑𝐐`$, but also as elements of $`𝐐_p`$, $`𝐐_p𝐐`$.
By using the $`p`$-adic frequency probability model for the statistical interpretation of quantum $`I`$-states we may assume that there exists $`I`$-states $`\varphi `$ (ensembles of $`I`$-transformers) such that the relative frequencies $`\nu _N(A=\lambda _n|\varphi )`$ have no limit in $`𝐑`$, i.e., we cannot apply the standard law of the large numbers in this situation. Hence if we realize measurements of an $`I`$-quantity $`A`$ for such a quantum $`I`$-state and study the observed data by using the standard statistical methods (based on real analysis), then we shall not obtain the definite result. There will be only random fluctuations of relative frequencies, see , .
Remark 8.2. Such a behaviour can be related to psychological experiments. Here the possibility of the use of $`p`$-adic probability models gives the important consequence for scientists doing experiments with a statistical $`I`$-data: the absence of the statistical stabilization (random fluctuation) does not imply the absence of an $`I`$-phenomenon. This statistical behaviour may have the meaning that this $`I`$-phenomenon cannot be described by the standard Kolmogorov probability model.
We now discuss other interesting implications of $`p`$-adic probability theory. There exists statistical samples , in which the frequencies $`\nu _N0`$, in the standard real topology, but $`\nu _N\alpha 0`$ in $`𝐐_p.`$ In this case the usual (Mises) frequency probability $`𝐏(A=\lambda |\varphi )=0`$. This implies that we have to consider the event $`\{A=\lambda |\varphi \}`$ (an observation of the information $`\lambda `$) as nonphysical event. However, from the point of view of the $`p`$-adic probability theory this is the physical event (of course, in the sense of $`I`$-physics).
The evolution of a $`p`$-adic wave function is described by an $`I`$-analogue of the Schrödinger equation:
$$\frac{h_p}{i}\frac{\psi }{t}(t,x)=\frac{h_p^2}{2m}\frac{^2\psi }{x^2}(t,x)V(t,x)\psi (t,x),$$
(20)
where $`m`$ is the $`I`$-mass of a quantum $`I`$-transformer. Here a constant $`h_p`$ plays the role of the Planck constant. By pure mathematical reasons (related to convergence of $`p`$-adic exponential and trigonometric series) it is convenient to choose $`h_p=\frac{1}{p}.`$
We may also present some physical arguments for such a choice. In ordinary quantum mechanics the Planck constant is related to the measure of discretization. The constant $`h_p=\frac{1}{p}`$ is related to the level of discretization of information.
If we use the statistical interpretation of quantum mechanics then the parameter $`t`$ plays the role of common $`I`$-time for elements of a statistical ensemble of $`I`$-transformers described by the wave function. Therefore, to be able to describe the evolution of a quantum state $`\psi ,`$ we must have consistent $`I`$-times for elements of this statistical ensemble.
We use the factor $`i=\sqrt{1}`$ in (20), because we like to have the total coincidence with formulas of the ordinary quantum mechanics. As we have already noted, in the $`p`$-adic case the functions $`e^{i\alpha x}`$ and $`e^{\alpha x}`$ have the same (non-oscillating) behaviour. Therefore, in principle, we can use the analogue of (20) in that the factor $`i`$ is omitted.
The use of $`i`$ implies the consideration of the extension $`𝐐_p(i)=𝐐_p\times i𝐐_p`$ of $`𝐐_p.`$ Elements of this extension have the form $`z=a+ib,a,b𝐐_p.`$ This extension is well defined for $`p=3,\mathrm{mod}\mathrm{\hspace{0.33em}4}.`$ As usual, we introduce a congugation $`\overline{z}=aib;`$ here we have $`z\overline{z}=a^2+b^2.`$ In what follows we assume that wave functions take values in $`𝐙_p(i)=𝐙_p\times i𝐙_p.`$
Example 8.4. ( A free $`I`$-transformer). Let the potential $`V=0.`$ Then the solution of the Schrödinger equation corresponding to the $`I`$-energy $`E=\frac{𝐩^2}{2m}`$ has the form<sup>16</sup><sup>16</sup>16We note that formal expressions for analytical solutions of $`p`$-adic differential equations coincide with the corresponding expressions in the real case (in fact, we can consider these equations over arbitrary number field, see ). However, behaviours of these solutions are different. :
$$\psi _𝐩(t,x)=e^{i(𝐩xEt)/h_p}.$$
(21)
By the choice $`h_p=1/p`$ this function is well defined for all $`x𝐙_p`$ and $`t𝐙_p`$. As $`\psi \overline{\psi }1,`$ this wave function describes the uniform ($`p`$-adic probability) distribution, see , on the ring of $`p`$-adic integers $`𝐙_p.`$ Thus an $`I`$-transformer $`\tau `$ in the state $`\psi `$ can be observed with equal probability in any state $`x𝐙_p.`$ In this sense behaviour of a free $`I`$-transformer is similar to behaviour of the ordinary free quantum particle. On the other hand, there is no analogue of oscillations: $`\psi _𝐩(t,x)=\mathrm{cos}(𝐩xEt)/h_p+i\mathrm{sin}(𝐩xEt)/h_p,`$ and $`|\mathrm{cos}(𝐩xEt)/h_p|_p=1,|\mathrm{sin}(𝐩xEt)/h_p|_p=|(𝐩xEt)/h_p|_p.`$
Remark 8.4. Is it possible to reproduce oscillations with respect to ordinary real time on the basis of the information model? It could be done by a time scaling. Let $`f:𝐙_p𝐙_p`$ be an arbitrary continuous function. Then $`f(t+kp^n)f(t)`$ for all $`k𝐙`$ for sufficiently large $`n`$ (uniformly for $`t𝐙_p`$). Let $`t_{\mathrm{phys}}=g(t)`$ be a law of the correspondence between $`I`$-time $`t𝐙_p`$ and real time $`t_{\mathrm{phys}}𝐑.`$ If $`2\pi =g(p^n)`$ then the $`p`$-adic continuity will imply the periodicity in real time. Therefore, the ordinary wave behaviour is nothing other than a consequence of continuity of information flows and the appropriative choice of a time scale. Depending on a time scale an $`I`$-process may or may not exhibit wave behaviour in the real picture of reality.
We consider a psychological (and social) consequence of Example 8.4: in the absence of the external potential the same motivation $`𝐩`$ may imply any $`I`$-state $`x𝐙_p.`$
Let us consider mixtures of states of the form (21). We set $`t=0.`$ Let $`\psi (x)=a_1\psi _{𝐩_1}+a_2\psi _{𝐩_2},a_1,a_2𝐙_p.`$
If we compute $`<\psi ,\psi >=_{𝐙_p}\psi (x)\overline{\psi (x)}𝑑x`$ (where $`dx`$ is a uniform $`p`$-adic valued distribution on $`𝐙_p`$) we see a large difference with ordinary quantum mechanics: $`<\psi ,\psi >a_1\overline{a}_1+a_2\overline{a}_2.`$ There is nonzero correlation term. For $`\alpha =(𝐩_1𝐩_2)/h_p,`$ we have :
$$T(\alpha )=<\psi _{𝐩_1},\psi _{𝐩_2}>+<\psi _{𝐩_2},\psi _{𝐩_1}>=\frac{\alpha \mathrm{sin}\alpha }{1\mathrm{cos}\alpha }.$$
Thus there are correlations between the motivations $`𝐩_1`$ and $`𝐩_2`$ in the state $`\psi .`$ By using the individual interpretation of quantum mechanics we say that an $`I`$-transformer $`\tau `$ with the wave function $`\psi `$ is in the superposition of two motivations $`𝐩_1`$ and $`𝐩_2.`$ Moreover, these motivations could not be measured exactly (compare with , ).
Such a situation is natural for psychological and social phenomena. In fact, a psychological or social motivation may be not represented in the brain in the definite form before the act of a measurement (at least for some quantum information states). Moreover, it cannot be measured exactly. Such information measurements may be used as illustrations of the process of a measurement in ordinary quantum mechanics. By analogy we can say that the definite value of a physical observable is created in a long process of the interaction with an equipment. Moreover, it can be never measured exactly (compare with - ).
Example 8.5. (Quantum Hooke’s system) To give an example of a Hamiltonian with discrete spectrum, we consider the formal $`p`$-adic generalization of the Hamiltonian of a harmonic oscillator:
$$\widehat{H}=\frac{h_p^2}{2m}\frac{d^2}{dx^2}\frac{1}{2}m\omega ^2x^2\frac{1}{2},$$
where $`m`$ is the $`I`$-mass. We consider $`\omega `$ simply as the coefficient of interaction (there is no analogue of harmonic oscillations). The operator $`\widehat{H}`$ has eigenvalues $`E_n=h_p\omega n,n=0,1,\mathrm{}`$ (see ). However, in the $`p`$-adic case the difference between continuous and discrete spectra is not so strong (for each $`E_n,`$ we have $`E_n=lim_k\mathrm{}E_{l_k},l_kn`$). On the other hand, discreetness of a spectrum, of course, induces some restrictions on values (information) which can be observed.
## 9 Appendix 1: Models of reality and number systems
Since Newton’s time, we use a model of physical reality based on a description of all physical processes by real numbers. In fact, the use of real numbers is equivalent to the assumption that any physical quantity can be measured (at least in principle) with an infinite precision. We shall discuss this point more carefully.
To realize a measurement of a physical quantity $`x`$, first we have to fix a unit of a measurement $`l=1`$. We assume that there exists such a natural number $`n`$ that
$$\left(n1\right)lx<nl.$$
(22)
This assumption is a mathematical postulate, the Archimedean axiom. Therefore by (22) we restrict our considerations to physical phenomena which can be described on the basis of the Archimedean mathematical model.
We now consider the next step of the measurement process. If $`y_1=(n1)lx`$ then we have to measure the quantity $`x_1=xy_1`$ by using a smaller unit of the measurement. Typically we fix a natural number $`m>1`$ ( the scale of the measurement) and choose the new unit $`l_1=l/m`$. Then we apply the Archimedean axiom (22) to the quantities $`x_1`$ and $`l_1`$ and obtain a natural number $`\beta _1`$ $`(\beta _1=1,\mathrm{},m))`$: $`(\beta _11)l_1x_1<\beta _1l_1`$. This procedure can be continued. If $`y_2=(\beta _11)l_1x_1`$ then we can use the new unit of measurement $`l_2=l_1/m`$ to measure the quantity $`x_2=x_1y_2`$ and so on. We remark that
$$x=\left(n1\right)+x_1=\left(n1\right)+\frac{\alpha _1}{m}+x_2=n1+\frac{\alpha _1}{m}+\mathrm{}+\frac{\alpha _n}{m^n}+x_{n+1},$$
(23)
where $`\alpha _k=\beta _k1=0,1,\mathrm{},m1`$. To obtain the real numbers model for physical reality, we assume that the above process of measurements of every physical quantity $`x`$ can be continued by an infinite number of steps. We call this postulate the postulate of an infinite precision of measurements or the Newton axiom. By this axiom any physical quantity $`x`$ can be identified with the real number:
$$x=\mathrm{}+\frac{\alpha _n}{m^n}+\mathrm{}+\frac{\alpha _1}{m}+\alpha _0+\alpha _1m+\mathrm{}+\alpha _km^k=\alpha _k\mathrm{}\alpha _0,\alpha _1\mathrm{}\alpha _n\mathrm{},$$
(24)
where $`\alpha _{\pm j}=0,1,\mathrm{},m1`$, (here the number $`(n1)`$, see (23), is also expanded with respect to powers of $`m`$).
Both the Archimedean and Newton axioms are natural for the description of an extended class of physical phenomena. The basis of the Archimedean-Newton model of reality is Newton’s space which is continuous, infinitely divisible and infinitely deep. All physical objects are located in this space and their location can be determined (at least in principle) with an infinite precision.
It would be natural to develop other models of physical reality which are not based on the Archimedean and Newton axioms.
The quantum formalism is one of successful attempts to give a new model of physical reality. The Archimedean and Newton axioms cannot be applied to quantum observables. However, quantum theory uses the old mathematical basis, real numbers. Of course, such a situation when non-Archimedean and non-Newtonean phenomena are considered in ”real” reality should induce paradoxes. One of such paradoxes is the EPR paradox which gives the right consequence, the death of reality (see, for example, , for the details).
In the present paper we started to develop a new model of reality, information reality, based on systems of $`m`$-adic numbers. This is a non-Archimedean model. Here we could not ‘compare’ two arbitrary information quantities $`x,y𝐙_m.`$ This is a non-Newtonian model. Information could not be ‘measured with an infinite precision.’ Information spaces are discrete. There always exists a ‘minimal space length’, a bit of information<sup>17</sup><sup>17</sup>17Our model of reality is closely connected with physical theories based on the fundamental length formalism or discrete space-time, see, for example -.
## 10 Appendix 2. An interpretation of quantum information mechanics.
We shall discuss the problem of interpretation on the basis of an example of a cognitive quantum system. Our analysis of measurement processes for quantum cognitive systems implies that we have to use the ensemble interpretation. There are two main viewpoints on the ensemble interpretation. The first is based on realism (see, for example, ). Here every physical quantity $`A`$ pertaining to a quantum state $`\varphi `$ which describes a statistical ensemble $`S`$ has some definite value for each $`sS.`$ The second is based on empiricism. Here it is not assumed the ”objective existence” of definite values of $`A.`$ The $`\varphi `$ gives only probabilities that $`A`$ would take some values (if a measurement of $`A`$ is performed). In fact, in conventional quantum mechanics the ensemble interpretation is based on the stronger form of empiricism: it could not be assumed that $`A`$ has definite ”objective value” $`a`$ for each $`sS.`$ Our analysis implies rather strange consequences. On the one hand, we understood that it is impossible to interpret a quantum state $`\varphi `$ on the basis of pure realism. On the other hand, we could not follow the conventional ensemble interpretation. On the one hand, a preparation procedure $`_b`$ which can be realized as a filter with respect to values $`\left\{b_j\right\}`$ of an $`I`$-quantity $`B`$ produces a statistical ensemble $`S`$ (described by the quantum state $`\varphi `$) in which each individual $`I`$-system has some fixed (objectively existing) value $`b_j.`$ On the other hand, there exist $`I`$-quantities for which we cannot assume that their have definite (”objective”) values for elements $`sS.`$ If $`A`$ is such an $`I`$-quantity, then its values are generated in the process of a measurement $`_a`$ over elements of the $`S.`$ This measurement does not imply a discontinuous collapse of the state $`\varphi `$ to the state $`\varphi _a`$ corresponding to the fixed value $`a`$ of $`A.`$ In the contrary this value $`a`$ of $`A`$ is created in the long process of the interaction between an $`I`$-system and a measurement device<sup>18</sup><sup>18</sup>18Thus the $`\varphi `$ is also related to individual $`I`$-systems.. There is the clear evidence that at least cognitive systems have hidden $`I`$-variables ($`I`$-states of the brain which are represented by configurations of excited neurons, see, for example, , ) which exist objectively and determine with some probabilities results of the $`_a.`$ In fact, these hidden $`I`$-variables (at least some of them) are not longer hidden. The modern experimental neuroscience gives the possibility to observe configurations of excited neurons corresponding to different reactions (results of $`_a),`$ see, for example, , on experiments (based on functional magnetic resonance imaging machine) for memory neurons configurations.
Remark 10.1. (Distribution of cognitive information in real space) Let $`km>1`$ be natural numbers. We introduce a map
$$j_{mk}:𝐙_m[0,1],x=\underset{l=0}{\overset{\mathrm{}}{}}\alpha _jm^jx_𝐑=\underset{l=0}{\overset{\mathrm{}}{}}\frac{\alpha _l}{k^{l+1}}.$$
We can present the speculation that one of maps $`j_{mk}`$ gives the spatial distribution of information in a cognitive system (in the case of “one dimensional brain”). This spatial distribution can have quite exotic structure. For example, the image $`j_{mk}\left(U_r\left(a\right)\right)`$ of a ball $`U_r\left(a\right)`$ can be a fractal (a kind of dusty set) in the real space. This model can have some connection with frequency domain methods in that populations of cortical oscillators self-organize by frequencies; same-frequency sub-populations of oscillators can interact in the sense that a change in phase deviation in one will be felt by the others in the sub-population. Thus here the spatial nearness of neurons (and even the existence of synaptic connections between two neurons) does not guarantee that they interact.
To simplify our considerations, we consider only states which provide the conventional (Kolmogorov) probability interpretation.
In what follows students can be considered as analogues of quantum particles and professors as analogues of measurement devices. Students of a University have to pass a test $`.`$ They have to give n answer to the question $`L.`$ Denote the set of all possible answers $`\left\{a_j\right\}`$ to $`L`$ by the symbol $`𝒜.`$ If we use the coding alphabet $`\{0,1,\mathrm{},p1\},`$ then the elements of $`𝒜`$ can be presented by $`p`$-adic integers. To prepare to this test, students have to read one of books $`=\left\{b_j\right\}`$ (here $`b_j𝐙_p`$ are coding sequences for books); a student has no time for reading of more than one book. All books give a description of the subject; but these descriptions are not identical. The process of reading is considered as a preparation procedure $`_b.`$ It produces a statistical ensemble $`S`$ of students who have read a book $`b_j.`$ The $`_b`$ can be considered as a filter on the set of all students of University. By quantum formalism (with the ensemble interpretation) $`S`$ is described by a quantum state $`\varphi `$ (which is a vector in a $`p`$-adic Hilbert space). This state is presented in the form:
$$\varphi =c_j\varphi _j,$$
(25)
where $`c_j=g_j+if_j,g_j,f_j𝐐`$ and $`\varphi _j`$ is a quantum state which describes the statistical sub-ensemble $`S_j`$ of $`S`$ consisting of students who have read the book $`b_j.`$ The number $`v_j=c_j\overline{c}_j`$ gives the frequency of students $`sS_j`$ in $`S,`$ i.e., proportional probability $`\left|S_j\right|/\left|S\right|`$ (where, for a set $`O,`$ we denote its cardinality by the symbol $`\left|O\right|`$).
The measurement $``$ is the process of an interaction between a quantum $`I`$-transformer (student’s brain) and a measurement equipment (professor’s brain). Brains are $`I`$-systems with very complicated internal structure. A result of interaction between the brain of a student $`sS`$ and the brain of a professor cannot be uniquely determined by the information $`b_j.`$ Moreover, an attempt to verify the condition $`sS_j`$ (by an additional measurement) may change the result of the measurement $`.`$ Therefore the property to give the answer $`a_k`$ as a result of the measurement $``$ is not an objective property of elements of the statistical ensemble $`S`$ described by $`\varphi .`$
Remark 10.2. We may use the notion of potentia: each $`sS`$ is potentially present in all states $`\psi _k=`$(the answer $`a_k`$ to the question $`L`$). The interaction with the equipment induces a transition from possible to actual.
The state $`\varphi `$ can be presented in the form:
$$\varphi =d_k\psi _k,$$
(26)
where $`d_k=m_k+in_k,m_k,n_k𝐐.`$ Probability to obtain the answer $`a_k`$ is given by the standard formula $`u_k=d_k\overline{d}_k.`$ We could not consider probability $`u_k`$ as probability with respect to a statistical ensemble. This is not proportional probability of the form $`\left|S_k^{}\right|/\left|S\right|,`$ where $`S_k^{}`$ is a statistical sub-ensemble of $`S.`$ In fact, the expansion (26) provides a description of some properties of an individual system $`sS`$ (reactions of $`s`$ to the question $`L`$). However, we cannot assume that $`\varphi `$ provides a complete description of the $`I`$-state of a cognitive system $`s.`$ Thinking systems of students can be very different<sup>19</sup><sup>19</sup>19The complete description of the $`I`$-states can be obtained on the basis of hidden variables models. The use of hidden information parameters is very natural in quantum information theory. For example, the brain contains a large number of information parameters which determine results of $`I`$-measurements. These parameters are really hidden in the subconsciousness. The consciousness cannot control them (see for the details)..
As usual, we introduce a diagonal operator $`A`$ in a $`p`$-adic Hilbert space, $`A\psi _i=a_i\psi _i.`$ The spectrum of $`A`$ coincides with the set of answers $`𝒜.`$ This operator provides the quantum description of the measurement $`.`$
As we have already noted, the act on the observation is a part of the measurement process $`.`$ In our example it is important that a student must give an answer to a professor. If we change the measurement procedure and consider a self-observation instead of an answer to the professor, then the states $`\psi _i`$ will be changed (with the corresponding change of probabilities).
Finally we have to remark that the quantum $`I`$-formalism can be used to construct a new model for Bohm’s pilot wave theory. In fact our approach is quite adequate to ideas of D. Bohm and B. Hiley on active information. Moreover, it seems that the ordinary pilot wave theory might be improved by considering the $`\psi `$-function field as a purely $`I`$-field.
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warning/0003/cond-mat0003322.html | ar5iv | text | # Microscopic Scenario for Striped Superconductors
## 1 Introduction and Scenario
Few recent problems in science have generated so many controversial discussions as the problem of high temperature superconductivity since its experimental discovery in 1986. Two fundamental questions are: Is the superconducting state found in high-$`T_c`$ cuprates homogeneous? Is the superfluid density in these materials homogeneous? The standard approach consists in accepting a homogeneous superconducting state of various forms. We assert that this state is in fact “inhomogeneous.” In this work we summarize and expand on a new scenario for striped superconductors where the interplay between inhomogeneous superfluid density and phase fluctuations determines the critical temperature.
In their undoped state, cuprates behave as antiferromagnetic (AF) Mott insulators and it is precisely upon doping with holes that these strongly correlated materials become superconductors. Recent experiments seem to indicate that inhomogeneously textured (intrinsically nanoscale) phases characterize the quantum state of high temperature superconductors. This is, probably, not surprising in retrospect since these are complex materials with competing time and length scales arising from different interactions. A relevant and non-trivial question is, however, whether those textures are essential to drive the phase coherent state, i.e., a Meissner phase.
Neutron scattering experiments have proven to be a very useful tool in investigating magnetic and superconducting properties of high-$`T_c`$ cuprate oxides. With improved sample quality and resolution there is reliable evidence for an incommensurate structure in the spin susceptibility. On the other hand, recent angle-resolved photoemission spectroscopy (ARPES) data suggest a one-dimensional (1D) like electronic structure consistent with clustering of charge carriers into 1D channels . Therefore, although the orientation, width, length and dynamics of the channels remains to be elucidated, both the above classes of experiments appear to confirm a new paradigm of spin and charge ordering in high-$`T_c`$ superconductors: the “stripe” phase.
Motivated by this new paradigm we introduced a class of inhomogeneous microscopic models which capture the magnetic and superconducting properties of these strongly correlated materials. The origin(s) of the mesoscopic skeleton of stripe segments in the CuO<sub>2</sub> planes is presently unclear and several mechanisms could be responsible, such as local spin-orbit coupling, Jahn-Teller distortions, oxygen buckling at the stripe, and/or other magnetoelastic effects. (Local charge-lattice coupling may be an important source of texture formation.) We showed, in particular, that appropriate inhomogeneous interactions that break magnetic symmetries are distinctive in inducing substantial pair binding of holes, as well as explaining the magnetic neutron scattering properties. Moreover, based upon the phenomenology of our microscopic model we developed a mean-field (“Josephson spaghetti”) model which provides a scenario for the macroscopic superconducting state. We also discussed the connection of the resulting inhomogeneity-induced superconductivity to recent experimental evidence for a linear relation between magnetic incommensurability and the superconducting transition temperature, as a function of doping. In a different work we studied the spectral properties of these inhomogeneous models and found, consistent with experiments , a flat band and the correct distribution of quasiparticle weights.
In previous work we have assumed static magnetic inhomogeneities. Certainly, the stripe segments in real materials are likely to have an intrinsic dynamics on a characteristic time scale $`\tau `$. We assume that this time is large enough for attractive forces to produce bound states of two holes. On the other hand, this stripe dynamics will probably restore the $`SU(2)`$ spin rotation invariance on time scales greater than $`\tau `$. $`SU(2)`$ symmetry does not have to be broken statically. In the present manuscript we also discuss an alternative approach where the magnetic inhomogeneities follow the hole, i.e., the stripe phase is dynamically generated by the holes. This model we call the selfconsistent perturbing hole (SPH) model since the hole itself carries the perturbation. The main qualitative difference between both classes of models is the absence of a broken lattice translational symmetry state in the SPH case. Otherwise, the basic phenomenology is qualitatively the same: Holes pair in stripes as a consequence of the existence of an AF background (avoiding a possible global phase separation). The pairing mechanism is kinetic exchange-interaction based and is provided by magnetic inhomogeneities that locally break spin-rotational invariance. In the superfluid phase, it is argued that a phase-locked state is generated as a consequence of a coherent Josephson tunneling of the hole pairs between and along stripes .
## 2 Microscopic Inhomogeneous Models
In this Section we will present two classes of inhomogeneous models. The first model was already introduced in Ref. , where the basic microscopic scenario starts from a homogeneous $`t`$-$`J`$ Hamiltonian as background
$$H_{tJ}=t\underset{𝐫,\overline{𝐫},\sigma }{}c_{𝐫\sigma }^{}c_{\overline{𝐫}\sigma }^{}+J\underset{𝐫,\overline{𝐫}}{}(𝐒_𝐫𝐒_{\overline{𝐫}}\frac{1}{4}\overline{n}_𝐫\overline{n}_{\overline{𝐫}})$$
(1)
but, to mimic the stripe segments, we add inhomogeneous magnetic interactions. These inhomogeneous terms break translational invariance and spin-rotational $`SU(2)`$ symmetry locally:
$$H_{\mathrm{inh}}=\underset{\alpha ,\beta }{}\delta J_zS_\alpha ^zS_\beta ^z+\frac{\delta J_{}}{2}\left(S_\alpha ^+S_\beta ^{}+S_\alpha ^{}S_\beta ^+\right)$$
with $`\delta J_{}\delta J_z`$, representing the magnetic perturbation of a static local Ising anisotropy, locally lowering spin symmetry (named a $`t`$-$`JJ_z`$ model). Only a few links (where the stripes are located) have this lowered spin symmetry. This Ising anisotropy is also sufficient to produce a spin gap.
Pair binding of holes in this class of models is substantial . This substantial binding energy is achieved as the energy for two holes falls much faster with $`t`$ than twice the energy of one hole.
It is interesting to make the following remarks: Suppose that a stripe segment is represented by the 1D $`t`$-$`J_z`$ model
$$H=t\underset{\alpha ,\sigma }{}(c_{\alpha \sigma }^{}c_{\alpha +1\sigma }^{}+\mathrm{H}.\mathrm{c}.)+J_z\underset{\alpha }{}S_\alpha ^zS_{\alpha +1}^z.$$
(2)
One can show that, within the ground state subspace (for a given number of holes), this Hamiltonian maps into the attractive spinless fermion model
$$H=t\underset{\alpha }{}(b_\alpha ^{}b_{\alpha +1}^{}+\mathrm{H}.\mathrm{c}.)\frac{J_z}{4}\underset{\alpha }{}\stackrel{~}{n}_\alpha ,\stackrel{~}{n}_{\alpha +1}$$
(3)
where $`\stackrel{~}{n}_\alpha =b_\alpha ^{}b_\alpha ^{}`$, and which certainly has a superconducting phase (i.e., correlation exponent $`K_\rho >1`$) . In this particular model, it turns out that (e.g., at half-filling for $`|J_z/8t|<1`$) the isolated stripe segment belongs to the Luttinger liquid universality class. Notice, however, that our stripes are embedded in an AF background. This background provides a strong boundary condition that results in an additional attractive potential for the holes in the stripe. As a result, an enhanced superconducting region is expected , avoiding alternative charge density wave phases or phase segregation.
Knowing that perturbating the system by breaking magnetic symmetries is an efficient mechanism to achieve substantial binding of carriers, some natural questions arise: What would happen if the hole itself carries this perturbation? Would this process be sufficient to generate a stripe phase? Would the binding energy be still appreciable in the thermodynamic limit?
Our SPH model corresponds to the Hamiltonian $`H=H_{tJ}+H_{\mathrm{sph}}`$, where
$$H_{\mathrm{sph}}=\underset{𝐫,𝐝}{}\frac{\delta J_{}}{2}(1\overline{n}_𝐫)\left(S_𝐫^+S_{𝐫+𝐝}^{}+S_𝐫^{}S_{𝐫+𝐝}^+\right).$$
(4)
In this Hamiltonian $`1\overline{n}_𝐫=n_𝐫`$ is the ocupation number of holes at site $`𝐫`$ and $`S_𝐫^+`$, $`S_𝐫^{}`$ are the usual spin operators . The presence of a hole at site $`𝐫`$ perturbs the magnetic links in the directions defined by $`𝐝`$ by lowering the spin symmetry and making them more Ising-like.
This model is perhaps more natural on physical grounds than the $`t`$-$`JJ_z`$ one. Magnetoelastic effects caused by the presence of the hole, or buckling of the oxygens close to the carrier may easily produce an Ising-like anisotropy. Upon doping with holes it is not obvious what the extra hole does to the environment. Since these are strongly correlated materials, the extra holes could have a stronger influence on the system than just those effects produced by simple hopping dynamics.
The models considered above are different in some respects. While in the model of Eqs. 1,2 translational symmetry has been explicitly broken (i.e., adding $`H_{\mathrm{inh}}`$) the model defined by $`H_{\mathrm{sph}}`$ is translationally symmetric, and the only symmetry that has been explicitly broken is the spin $`SU(2)`$ symmetry around each hole. This fact has some direct experimental consequences. As already discussed in Ref. , the $`t`$-$`JJ_z`$ model has an inhomogeneous hole density, as they prefer to occupy sites where the magnetic links have been weakened (the stripes). This is not the case for the SPH model since, as the translational symmetry has not been broken, the holes will be found with equal probability on every lattice site. This fundamental difference could be resolved by a Scanning Tunneling Microscope experiment.
In this work we will not address the issue of stripe formation in the SPH model, but rather concentrate on the hole pairing in those textures. Clearly there will be a competition between kinetic and magnetic energies. While the first will try to delocalize the pair, the former will contribute to the pairing. In Fig. 1 we show the binding energy (defined as $`E_b=E_2+E_02E_1`$, where $`E_i`$ is the ground state energy in the subspace with $`i`$ holes) as a function of $`1/N`$ (where $`N`$ is the size of the system). $`t`$-$`J`$ models in 1D (diamonds) and 2D (triangles) have zero or negligible binding in the extrapolated thermodynamic limit. The model labeled as $`t`$-$`JJ_z`$ corresponding to the Hamiltonian of Eqs. 1 ,2 in 1D has substantial binding in the thermodynamic limit.
That model corresponds to placing a perturbation like Eq. 2 every 4 sites. The model labeled as $`t`$-$`J`$ \+ $`H_{\mathrm{sph}}`$ is the one we are mostly interested in this work. It has appreciable binding of holes in the thermodynamic limit, although not as strong as in the $`t`$-$`JJ_z`$ model. In the calculations we have taken $`|𝐝|=2`$ (i.e., the hole perturbs up to second neighbors). For $`|𝐝|=1`$ the model reduces trivially to the $`t`$-$`J`$ model, as the links close to the holes are magnetically inactive. Therefore $`|𝐝|=d`$ must be greater or equal to 2 to show any new behavior. In order to get information about the bound state we have calculated the correlation function $`n_0n_i`$, which gives the probability of finding a hole at site $`i`$ if there is one at site $`0`$. We remark again that the density of holes is constant and equal to $`n_i=N_h/N`$ (with $`N_h`$ the number of holes and $`N`$ the size of the system) for every site $`i`$ in the case of $`t`$-$`J`$ \+ $`H_{\mathrm{sph}}`$. In Fig. 2, upper panel we show this correlation function for the cases $`d=2`$ and $`d=3`$ when there are two holes in a 16 sites chain. As expected, the pair is more strongly bound for $`d=3`$. In a real material there must be a decreasing function of the perturbation carried by the hole with distance. In our case that perturbation is constant. In Fig. 2 lower panel, we show the same correlation function for $`d=2`$ and different $`t`$ values. The pair evolves from being tightly bound to delocalization as the kinetic energy is increased. It is worth noting that even when the holes are not very close the binding energy scales to a finite value.
We are grateful to C.D. Batista for useful discussions. This work has been supported by the US DOE. |
warning/0003/astro-ph0003470.html | ar5iv | text | # WHY ARE ROTATING ELLIPTICAL GALAXIES LESS ELLIPTICAL AT X-RAY FREQUENCIES?
## 1 INTRODUCTION
Cooling flows in massive elliptical galaxies are expected to rotate since most of the interstellar gas within an optical effective radius $`R_e`$ is produced by mass loss from the slowly rotating stellar systems. If the cooling flow proceeds inward conserving angular momentum, the interstellar gas ultimately flows toward a large disk (Rdisk
>Re
>subscript𝑅𝑑𝑖𝑠𝑘subscript𝑅𝑒R_{disk}\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}R_{e}) and spins up to the local equatorial circular velocity, $`400`$ km s<sup>-1</sup> (Kley & Mathews 1995; Brighenti & Mathews 1996, subsequently referred to as Paper I). As interstellar gas approaches the disk, the X-ray images are significantly flattened toward the equatorial plane on scales $`R_e3R_e`$ when viewed perpendicular to the axis of rotation. For bright elliptical galaxies at distance 17 Mpc in the Virgo cluster – NGC 4472, NGC 4636, and NGC 4649 – $`R_e7\mathrm{kpc}1.4`$’. However, Einstein X-ray images show no evidence of rotational flattening either at the resolution of the IPC ($`55^{\prime \prime }4.5`$ kpc; Trinchieri et al. 1986) or at the much higher resolution or HRI ($`6^{\prime \prime }0.49`$ kpc; Fabbiano, Kim & Trinchieri 1992). In a series of papers Buote & Canizares have detected large scale X-ray flattening in several massive elliptical galaxies which they attribute to a flattened dark matter potential, but the highly flattened inner isophotes anticipated from momentum-conserving cooling flows were not evident (see Buote & Canizares 1998 for a review). More recently, Hanlan & Bregman (1999) have found little or no rotationally enhanced X-ray ellipticity in ROSAT HRI and PSPC images of several bright elliptical galaxies. The implications of these observations are not widely appreciated among the community of astronomers interested in early type galaxies.
Massive, X-ray luminous elliptical galaxies have two principal sources of hot interstellar gas: (1) ejection of stellar envelopes from an evolving population of old stars and (2) inflow of distant circumgalactic gas which accumulates in the galactic halo either by tidal exchange or secondary infall (Brighenti & Mathews 1999a). If the circumgalactic gas has a lower specific angular momentum than the stars, the X-ray flattening would be lessened as inflowing circumgalactic gas mixes with gas lost from the rotating stellar system within several $`R_e`$. But it is unlikely that X-ray disks can be eliminated by this means since most of the gas in the inner galaxy is produced by stellar mass loss and has not flowed in from the galactic halo region. Our recent models for the large elliptical galaxy NGC 4472 require that at least 60 - 70 percent of the hot gas in r
<Re
<𝑟subscript𝑅𝑒r\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}R_{e} originates from stellar mass loss. This rate of stellar mass ejection is expected from normal stellar evolution and is required to explain the radial variation of interstellar temperature and metallicity typically observed in massive elliptical galaxies (Brighenti & Mathews 1999a). In addition, it is unlikely that the absence of observed X-ray flattening can be understood simply by the inflow of halo gas into the optical centers of elliptical galaxies since it is plausible that the extended halo gas has an even larger specific angular momentum than the stars. N-body simulations of galaxy formation typically produce elliptical galaxies and galaxy groups with increasing specific angular momentum at large radii (e.g. Barnes & Efstathiou 1987; Quinn & Zureck 1988). But the intrinsic angular momentum in distant halo gas is uncertain; gas acquired by tidal disruption is expected to have significant rotation, but secondary infalling gas that arrives after most mergers have occurred may have less net rotation than the stellar system.
However, there is no doubt that the stellar systems in most large elliptical galaxies and the gas that they expel are rotating. Although the stellar systems in luminous elliptical galaxies (with LB
>3×1010
>subscript𝐿𝐵3superscript1010L_{B}\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}3\times 10^{10} $`L_{B,}`$) are not flattened by rotation, their rotation about the minor axis is not small, typically 50 - 100 km s<sup>-1</sup> (e.g. Binney, Davies, & Illingworth 1990), and the gas they expel must spin up further as it moves inward.
These general considerations indicate that the net angular momentum of gas ejected from stars must be removed to circularize the X-ray images. This can be accomplished either by removing gas from the flow or by transporting angular momentum away from the axis of rotation by viscous interactions. In this paper we explore the importance of these two processes and determine their relative influence on the X-ray images. The first process, localized dropout cooling of interstellar gas in regions of low specific entropy, is also required to limit the masses of central black holes (or dark stellar nuclei) in elliptical galaxies to their observed values. Central black holes would be about ten times more massive if the cooling flow proceeded all the way to the galactic center with constant mass flux (Brighenti & Mathews 1999b). The presence of distributed optical emission lines from cooling gas also argues for cooling dropout (Mathews & Brighenti 1999a). However, as we show below, the strongest argument for the presence of localized mass dropout in the interstellar gas may be the circularization of X-ray images that results.
We also explore the possibility that interstellar viscosity transports angular momentum outward, reducing the rotation and X-ray flattening in the inner galaxy. An effective viscosity can arise from several natural sources: plasma viscosity, interstellar “turbulent” viscosity driven by stellar mass loss or Type Ia supernovae, the viscosity implicit in regions of localized interstellar cooling as they sink in the galactic potential, and turbulent viscosity that may develop from shear instabilities in the flow. On several occasions Nulsen and Fabian have suggested that the turbulent viscosity is so large that the entire interstellar medium becomes approximately spherical (Nulsen, Stewart, & Fabian 1984; Nulsen & Fabian 1995), but no detailed calculations were provided to support this conjecture.
In the following section we describe rotating cooling flows in which angular momentum is lost by cooling dropout. Then we study the outward transport of angular momentum by turbulent viscosity and the associated diffusion of interstellar metallicity.
## 2 COOLING DROPOUT IN ROTATING COOLING FLOWS
In the earliest studies of non-rotating cooling flows, the interstellar gas was assumed to flow completely to the galactic center where, in the region of increasing gas density, it cooled by intense radiative losses. The notion that the hot interstellar gas in cooling flows cools not just at the centers of elliptical galaxies but throughout the galactic volume was first advanced by Fabian, Nulsen & Canizares (1982) and Thomas (1986). This distributed cooling was introduced to avoid the strong central peaking of the X-ray surface brightness $`\mathrm{\Sigma }_x`$ that occurred in centrally-cooling models but not in the galaxies observed. Our recent studies indicate that the effect of volumetric cooling on $`\mathrm{\Sigma }_x`$ is more complicated. Although $`\mathrm{\Sigma }_x`$ is globally reduced if the gas cools before it has flowed through the entire potential of the galaxy, this reduction is generally rather modest. However, in localized regions of intense cooling dropout $`\mathrm{\Sigma }_x`$ actually increases because of the enhanced density and emission associated with cooling regions; if cooling dropout is restricted to a small but finite central region, $`\mathrm{\Sigma }_x`$ can deviate even further from observed profiles.
Notwithstanding these contrary results, the evidence for spatially distributed cooling of interstellar gas is apparent in other observations. The dynamically determined masses of central black holes (Magorrian et al. 1998) are typically much less than the masses of interstellar gas that have cooled over a Hubble time, typically several $`10^{10}`$ $`M_{}`$ in massive elliptical galaxies (Brighenti & Mathews 1999b). The cooled mass is also far larger than the masses of cold HI or H<sub>2</sub> gas observed in luminous elliptical galaxies (e.g., Bregman, Roberts & Giovanelli 1988; Braine, Henkel & Wiklind 1997), indirectly indicating that stars have formed in the cooled gas. However, if the stars produced from cooled gas within $`R_e`$ are optically dark (e.g., brown dwarfs), the thinness of the fundamental plane is difficult to understand (Mathews & Brighenti 1999b). This suggests that the young, low-mass stars that form are optically luminous, which is also consistent with H$`\beta `$ indices observed in the stellar spectra in some luminous elliptical galaxies (Gonzalez 1993; Faber et al. 1995; Trager 1997; Trager et al. 1998; Mathews & Brighenti 1999c). Another indication of distributed interstellar cooling is the extended optical line emission observed in in the central regions of most or all bright elliptical galaxies (e.g., Macchetto et al. 1996). We have argued that this gas with temperature $`10^4`$ K is cooled interstellar gas ionized by stellar UV (Mathews & Brighenti 1999a).
In the absence of spatially distributed mass dropout or angular momentum redistribution, all the cooled gas in rotating cooling flows is ultimately deposited in a large disk (Paper I; Brighenti & Mathews 1997). Because the gas density and radiative emissivity increase as the cooling flow approaches the disk, the X-ray images in these cooling flow models are highly flattened unless the galaxy is viewed along the axis of rotation. In the following discussion we show that spatially distributed mass dropout can sharply reduce, but not eliminate, the rotational flattening of X-ray images. As a result, the approximate circular shapes of observed X-ray images may be the most convincing evidence for the existence of cooling dropout. We begin by briefly discussing the gravitational potential and stellar dynamics in a simple rotating elliptical galaxy; this is followed with a description of a gas-dynamical model including the cooling dropout process.
### 2.1 Model of a Rotating Galaxy
For comparison with our earlier models of rotating elliptical galaxies and for mathematical simplicity, we adopt here the same E2 model galaxy described in detail in Paper I. The ellipsoidal mass distributions for stars and dark matter in this galaxy are assumed to have King and pseudo-isothermal profiles respectively:
$$\rho _{}(R,z)=\rho _o[1+(R/R_c)^2+(z/z_c)^2]^{3/2}$$
$$\rho _h(R,z)=\rho _{oh}[1+(R/R_{ch})^2+(z/z_{ch})^2]^1.$$
Both distributions are truncated at the same large ellipsoidal boundary defined by $`R_t`$ and $`z_t=R_t(1e^2)^{1/2}`$. The eccentricity of an E2 galaxy ($`n=2`$) is $`e=[1(1n/10)^2]^{1/2}=0.6`$. Scale parameters $`R_c`$ and $`R_{ch}`$ for the E2 elliptical can be derived from a standard spherical E0 galaxy having the same central densities $`\rho _o`$ and $`\rho _{oh}`$ by expanding the characteristic core radii $`R_c^{(s)}`$ for the spherical galaxy in the $`R`$-direction, $`R_c=R_c^{(s)}(10.1n)^{1/3}`$. The physical properties of the fiducial non-rotating E0 galaxy are listed in Table 1. Values of $`R_e`$, $`L_B`$ and $`\sigma _{}`$ have been chosen so that the E0 galaxy lies on the fundamental plane. The total mass of dark matter $`M_{ht}`$ is nine times that of the stellar component, $`M_t`$.
As described in Paper I, we solve the two-dimensional Jeans equation to derive the velocity dispersion in the meridional plane, $`\sigma ^2(R,z)=\sigma _R^2=\sigma _z^2`$. Following Satoh (1980), we decompose the mean square azimuthal speed into a random dispersion $`\sigma _\varphi ^2`$ and a systematic rotation $`\overline{v_\varphi }^2`$ using a single parameter $`k`$ assumed to be constant,
$$\sigma _\varphi ^2\overline{v_\varphi ^2}\overline{v_\varphi }^2=k^2\sigma ^2+(1k^2)\overline{v_\varphi ^2}.$$
When $`k=1`$ the stellar velocity dispersion is isotropic and the ellipsoidal distortion is produced entirely by rotation; when $`k=0`$ there is no net rotation and the ellipsoidal flattening is produced entirely by anisotropic velocity dispersion, $`\sigma _\varphi ^2>\sigma ^2`$. For a typical slowly rotating giant elliptical galaxy in which most (but not all) of the optical ellipticity is due to anisotropic orbits, we choose $`k=0.5`$ as in the “E2;0.25” model described in Paper I. From the Jeans equations we can determine a self consistent rotation for the stellar system, $`v_\varphi (R,z)`$, which is required for the gas dynamical equations. The total galactic potential $`\mathrm{\Phi }(R,z)`$, used in the gas dynamical equation of motion, can be derived from the density distributions above (see Paper I). The mean stellar temperature that appears in the gas dynamical equation for thermal energy is given by
$$T_{}(R,z)=\frac{\mu _mm_p}{3k_B}(2\sigma ^2+\sigma _\varphi ^2),$$
where $`\mu _m=0.62`$ is the mean molecular weight and $`m_p`$ is the proton mass.
### 2.2 Cooling Dropout
The physical origin of spatially localized interstellar cooling and mass dropout is a hypothetical ensemble of entropy fluctuations. For example, a small region in the cooling flow with low entropy (in pressure equilibrium with low $`T`$, high $`\rho `$) will radiate more efficiently and cool before the ambient flow reaches the galactic center or cooling disk. The entropy irregularities have complex origins – supernova explosions, stellar mass loss, galactic winds and other disturbances in the early universe, inhomogeneous magnetic fields, etc. – so it is not currently possible to derive them from first principles. However, by exploring the influence of a variety of assumed mass dropout profiles on the X-ray surface brightness and radial mass distributions in luminous elliptical galaxies, Brighenti & Mathews (1999b) confirmed that the following modification of the equation of continuity gives adequate results:
$$\frac{\rho }{t}+\rho 𝐯=\alpha _{}\rho _{}q\frac{\rho }{t_{do}}.$$
The negative term on the right hand side represents the disappearance of gas from the flow by cooling dropout. This sink term is proportional to the local gas density divided by the dropout time $`t_{do}=5m_pkT/2\mu _m\rho \mathrm{\Lambda }`$, the time for gas to radiatively cool locally at constant pressure. This kind of simple mass loss with $`q=1`$, similar to that used by Sarazin & Ashe (1989), was shown by Brighenti & Mathews (1999b) to be consistent with optical and X-ray observations of massive elliptical galaxies.
The first term on the right side of the equation of continuity above describes the source of new interstellar gas, which is proportional to the local stellar density and to
$$\alpha _{}(t)=\alpha _n[t/(t_nt_s)]^{1.3},$$
the specific mass loss from evolving stars. Here $`\alpha _n=4.7\times 10^{20}`$ s<sup>-1</sup>, $`t_n=13`$ Gyrs is the current cosmic time and $`t_s=1`$ Gyr is the time that the galactic stars are assumed to have formed (see Brighenti & Mathews 1999a for details). We are particularly interested in the flattening of X-ray images within $`R_e`$ from the centers of large ellipticals. In this central region most of the interstellar gas is produced by stellar mass loss, so we shall generally ignore additional gas flowing in from the galactic halo region beyond the optical galaxy.
Interstellar gas is heated and enriched by Type Ia supernovae which we assume occur at a rate SNu$`(t)=`$ SNu$`(t_n)(t_n/t)`$, where SNu$`(t_n)=0.03`$ SNIa per 100 yrs per $`10^{10}`$ $`L_B`$ (Cappellaro et al. 1997). The specific rate of mass injection into the cooling flow by Type Ia supernovae is $`\alpha _{sn}=\nu _{sn}m_{sn}/M_t`$ where $`\nu _{sn}=\mathrm{SNu}(t)(L_B/10^{10})/3.15\times 10^9`$ is the number of supernovae per second. The rate that gas is supplied to the interstellar medium by normal stellar mass loss greatly exceeds the mass input from supernovae, i.e., $`\alpha _{}\alpha _{sn}`$.
### 2.3 Rotating Cooling Flows with and without Dropout
Figures 1a and 1b illustrate the X-ray surface brightness $`\mathrm{\Sigma }_x(R,z)`$ in the ROSAT band (0.1 – 2.4 keV) for the rotating elliptical at time $`t_n=13`$ Gyrs when viewed along the equatorial plane and at 60<sup>o</sup> inclination. As in Paper I, our 2D logarithmic computational grid extends beyond 100 kpc in both $`R`$ and $`z`$, but Figure 1 zooms in to show only the central $``$15 kpc region that is most relevant to published X-ray images. Figure 1a shows the extremely flattened X-ray appearance of the E2 galaxy when no distributed cooling dropout is present. This diagram is similar to Figure 12b of Paper I with several small differences due to the larger value of $`t_n`$ used in that paper and the slightly lower Type Ia supernova rate that we assume here. When the same rotating cooling flow is observed at 60<sup>o</sup> inclination, the bright X-ray emission associated with the outer disk is seen as a separate feature at $`5`$ kpc (Figure 1b). The current ROSAT band luminosity for this galaxy within the optical image is $`L_x=2.0\times 10^{40}`$ ergs s<sup>-1</sup>, about 5 times less than the non-rotating version of the same galaxy (Paper I). However, this $`L_x`$ must be regarded as an underestimate because our computational proceedure is inaccurate when most of the radiative cooling occurs in the central zone (Brighenti & Mathews 1999b).
Figures 1c and 1d show the same X-ray views of the rotating galaxy at $`t_n=13`$ Gyrs but now with cooling dropout ($`q=1`$) and with all other parameters identical to those in Figures 1a and 1b. The surface brightness $`\mathrm{\Sigma }_x(R,z)`$ in Figures 1c and 1d includes the additional X-ray emission from the cooling regions, but the isophotal shapes are influenced very little by this contribution. It is immediately clear from Figure 1c that radiative dropout is a major explanation for the nearly circular appearance of observed X-ray images. Evidently, interstellar gas is cooling from the flow before it has progressed very far from its stellar origin toward the axis of rotation. However, the central X-ray bright core of the image is still noticeably flattened in Figure 1c with an X-ray ellipticity that is slightly greater than that of the E2 stellar image. (The optical isophotal contours would intersect the $`z`$-axis at 0.8 of the $`R`$-axis intersection.) This excess ellipticity of the X-ray image is not greatly diminished if the galaxy is viewed at an inclination $`i=60^o`$ (Figure 1d). The ROSAT luminosity for the $`q=1`$ model E2 galaxy at time 13 Gyrs is $`L_x=8.8\times 10^{40}`$ ergs s<sup>-1</sup>.
When there is no cooling dropout, $`q=0`$ as in Figure 1a, a large and massive disk of cold gas forms in the equatorial plane, rotating at the local circular velocity; such disks are described in detail in Paper I and in Brighenti & Mathews (1997). However, in the calculation with $`q=1`$ mass dropout by radiative losses as in Figure 1c, no disk of cold gas forms in the $`z=0`$ plane at least beyond $`150`$ pc, the size of our innermost computational grid zone. This failure to form a cold disk results from the increased radiative dropout rate as the density in the cooling flow increases toward the equatorial plane, $`(\rho /t)_{do}q\rho /t_{do}\rho ^2`$.
The azimuthal velocity distributions $`v_\varphi (R,z)`$ in the interstellar gas shown in Figures 2a and 2b correspond to the cooling flows imaged in Figures 1a and 1c respectively; interstellar rotation is significantly reduced in the presence of cooling dropout. Figure 2a shows the azimuthal velocity distribution at 13 Gyrs without radiative dropout. As described in Paper I, without dropout mass loss there is no loss of angular momentum and the slow inward velocity of the cooling flow leads to a dramatic spinup in angular velocity. In Figure 2a the gas is flowing toward a cool disk that extends to $`5`$ kpc. As interstellar gas converges toward the disk, its azimuthal velocity $`v_\varphi `$ approaches the local circular velocity, $`400`$ km s<sup>-1</sup>, as it cools and merges with the rotationally supported cooled disk.
With $`q=1`$ (Figure 2b) the rotational velocities are reduced by approximately 30 percent in the region shown within 15 kpc of the center. If the entire major axis (-5 to +5 kpc) were observed, the rotational broadening of X-ray emission lines would be $`400`$ km s<sup>-1</sup>, similar to the thermal width at the virial temperature, but resolution of such lines would be difficult or impossible even with the Chandra Observatory. For purposes of comparison with future high-resolution X-ray observations, we show in Figure 2c the emission-weighted line of sight velocity at every position for the $`q=1`$ calculation,
$$v_{los}(R,z)=\frac{n_en_p\mathrm{\Lambda }_{ROSAT}𝐯_\varphi 𝐧_s𝑑s}{n_en_p\mathrm{\Lambda }_{ROSAT}𝑑s}$$
where $`𝐧_s`$ is a unit vector along the line of sight represented with the $`s`$ coordinate. The contours in Figure 2b and 2c are rather similar due to the strong spatial gradient in electron density, but $`v_{los}(R,z)`$ is about 25 percent smaller than $`v_\varphi (R,z)`$ over most of the region plotted.
Since it is likely that the cooled gas forms into optically luminous stars (Mathews & Brighenti 1999b), it is of interest to compare in Figure 3 the local density of cooled (dropout) gas $`\rho _{do}(R,z)`$ for the $`q=1`$ solution with that of the old stars, $`\rho _{}(R,z)`$. Notice that the density distribution of cooled gas is considerably flatter than the E2 distribution of the old stellar population; the isodensity ellipticity of the gas just prior to dropout reflects rotational flattening in the galactic potential. The orbits of the young stellar population are expected to inhabit the same flattened volume as the interstellar dropout from which they formed, with stellar density contours similar to $`\rho _{do}(R,z)`$. For an optically luminous population of young stars formed from cooled gas, we predict from Figure 3 that the H$`\beta `$ equivalent width in the stellar spectrum (produced mostly by young stars) should exhibit a decidedly greater ellipticity than that of the bulk of the stellar light. In addition, since the interstellar gas in the cooling flow spins up before cooling, we expect that the young stellar population formed from mass dropout has a somewhat greater systemic rotation than that of the old stars.
For a larger dropout coefficient, $`q=4`$, the X-ray images are even more circular and the corresponding azimuthal velocity $`v_\varphi (R,z)`$ is lower than that of the $`q=1`$ solution. However, when $`q=4`$ the apparent gas temperature in the cooling flow is (unrealistically) lowered by the larger emission contribution from cooling sites. We have shown that the observed gas temperature profile in NGC 4472 cannot be fit if $`q4`$ (Brighenti & Mathews 1999b).
## 3 VISCOUS DIFFUSION OF ANGULAR MOMENTUM
We now consider an alternative possible explanation of the circularity of X-ray images of elliptical galaxies: transport of angular momentum away from the axis of galactic rotation by turbulent diffusion.
Flow-induced turbulence can occur when the ratio of inertial to viscous terms in the equation of motion is large, i.e. when the Reynolds number $`=\rho uL/\mu `$ is much larger than unity. Here $`\mu `$ is the viscosity and $`\rho uL`$ represent the characteristic density, velocity, and length scale in the cooling flow. In an unmagnetized plasma the “molecular” viscosity $`\mu `$ is proportional to the plasma mean free path $`\lambda T^2/n`$ which can be very large in the hot interstellar gas. However, Faraday depolarization studies in the interstellar gas in elliptical galaxies (Garrington & Conway 1991) indicate microgauss magnetic fields for which the proton Larmor radius $`r_L`$ is much smaller than $`\lambda `$, greatly increasing the effective Reynolds number and promoting the likelihood of turbulence. To estimate the Reynolds number we replace $`\lambda `$ with the thermal gyroradius $`r_L`$ in the expression for the viscosity, $`\mu \rho v_{th}r_L`$, where $`v_{th}=(3k_BT/m_p)^{1/2}`$ is the thermal velocity. The resulting Reynolds number for rotating cooling flows,
$$=\left(\frac{\rho ur}{\rho v_{th}r_L}\right)$$
$$3\times 10^{11}\left(\frac{v_\varphi }{200\mathrm{km}/\mathrm{s}}\right)\left(\frac{r}{\mathrm{kpc}}\right)\left(\frac{T}{10^7\mathrm{K}}\right)^{1/2}\left(\frac{B}{\mu \mathrm{G}}\right),$$
is enormous, so turbulent flow can be expected.
The Rayleigh criterion establishes a necessary condition to avoid shear turbulence in rotating flows: for stability to linear perturbations the specific angular momentum must increase from the rotation axis. By this criterion alone, rotating galactic cooling flows appear to be stable. Initially, as the gas first enters the flow, its specific angular momentum increases outward; this follows from the flat stellar rotation curves typically observed in slowly rotating massive elliptical galaxies (e.g. van der Marel 1991). We have verified that the sense of this gradient in the angular momentum density is preserved in the final cooling flows shown in Figure 1a.
However, rotating flows that satisfy the Rayleigh criterion may still be unstable to non-infinitesimal perturbations that result from turbulent mixing. A variety of non-linear perturbations are present in elliptical galaxies which can generate interstellar turbulence: transport and dissipation of mass ejected from orbiting stars, Type Ia supernovae, the radial sinking of denser cooling sites, buoyant and shearing magnetic fields, etc. Indeed, interstellar mixing (turbulence) due to stellar and supernova sources is likely to generate large magnetic fields in the cooling flow gas by the turbulent dynamo process. These fields, which can serve to redistribute angular momentum, may become dynamically important in the central regions of the flow where the magnetic energy density is further increased by the converging flow (Soker & Sarazin 1990; Moss & Shukurov 1996; Mathews & Brighenti 1997). However, if interstellar turbulence is invoked to account for the circularization of X-ray images, the associated spatial diffusion must not suppress abundance gradients observed in the hot interstellar gas (Matsushita 1997).
In the following discussion we estimate the effective viscosity corresponding to various sources of interstellar mixing or turbulence. Then we describe modifications to the gas dynamical equations that incorporate a spatially dependent turbulent viscosity. This is followed by a description of rotating cooling flows including both mass dropout and turbulent viscosity.
### 3.1 Sources of Interstellar Viscosity
The effective viscosity due to stellar mass loss and Type Ia supernovae can be estimated by constructing the viscosity from its simplest kinematical representation, $`\mu \rho \mathrm{}v`$, where $`\rho `$ is the local gas density, $`\mathrm{}`$ is the mean free path or coherence length of the disturbed region, and $`v`$ is the characteristic velocity appropriate to each process. To estimate the viscosity $`\mu _{}`$ generated by stellar mass loss, we use a length scale $`\mathrm{}_{}2a_{en}=4.5\times 10^9\rho ^{1/3}`$ based on the radius $`a_{en}`$ of a recently ejected stellar envelope in pressure equilibrium (Mathews 1990) and the mean turbulent velocity $`v_t=5.7\times 10^5r_{kpc}^{0.2}`$ cm s<sup>-1</sup> as estimated by Mathews and Brighenti (1997). In evaluating $`v_t`$ we use a stellar dispersion $`\sigma _{}=351`$ km s<sup>-1</sup> (Table 1) and assume that half of the kinetic energy in stellar ejecta supports the turbulence. With these dimensional values and taking $`\rho 1.05\times 10^{25}r_{kpc}^{1.8}`$ gm cm<sup>-3</sup> we find $`\mu _{}0.06r_{kpc}^{1.4}`$ gm cm<sup>-1</sup> s<sup>-1</sup> for the effective viscosity generated by stellar ejecta. Alternatively, the characteristic length $`\mathrm{}_{}`$ could be based on the total volume of interstellar gas disturbed by the motion of an ejected stellar envelope. If the stopping distance of the stellar envelope after ejection from an orbiting star is $`\mathrm{}_{st}1.2\times 10^{22}r_{kpc}^{0.6}`$ cm (Mathews 1990), then $`\mathrm{}_{}\mathrm{}_{vol}(\pi a_{en}^2\mathrm{}_{st})^{1/3}3.3\times 10^{19}r_{kpc}^{0.6}`$ cm. In this case the viscosity is somewhat larger, $`\mu _{}\rho v_t\mathrm{}_{vol}1.0r_{kpc}^{0.6}`$ gm cm<sup>-1</sup> s<sup>-1</sup>, reflecting the uncertainty in estimating $`\mu _{}`$.
For an equally approximate estimate of the viscosity corresponding to supernova remnants, we use an eddy size $`\mathrm{}_{sn}=3.6\times 10^{11}\rho ^{1/3}`$ cm based on the blast wave calculations of Mathews (1990) and a characteristic mixing velocity $`v_{sn}=2.0\times 10^6r_{kpc}^{0.2}`$ cm s<sup>-1</sup> from Mathews & Brighenti (1997), assuming that half of the supernova energy goes into interstellar turbulence. The resulting viscosity for Type Ia supernova explosions is $`\mu _{sn}16.2r_{kpc}^{1.4}`$ gm cm<sup>-1</sup> s<sup>-1</sup>. For both $`\mu _{}`$ and $`\mu _{sn}`$ we ignore a weak time dependence due to variable mass loss from evolving stars and the time dependent Type Ia supernova rate respectively. However, both $`\mu _{}`$ and $`\mu _{sn}`$ decline with galactic radius in a similar fashion due to the decrease in gas and stellar density, so $`\mu _{sn}\mu _{}`$ is expected at all $`r_{kpc}`$ from these formal estimates.
An additional viscosity arises from the cooling dropout process. The increased density in a cooling site relative to the surrounding cooling flow causes it to sink more rapidly in the galactic potential than the background flow. The net effect of the ensemble of differentially sinking sites is an approximately radial exchange of gas in the flow with a corresponding viscosity. When a site of density $`\rho _s>\rho `$ falls at its terminal velocity $`u_t`$, the drag force on the moving site must balance the gravitational force, $`\rho u_t^2\pi r_s^2(4\pi /3)r_s^3\rho _sg`$, where $`g5.2\times 10^7r_{kpc}^1`$ cm s<sup>-2</sup> is an approximation to the galactic gravitational acceleration. The radius $`r_s`$ of a cooling site can be estimated using the approximate model for an isobaric, steady state cooling site described by Mathews & Brighenti (1999a). Within the cooling site structure, the temperature varies nearly linearly with radius, $`T=[K\mathrm{\Lambda }(T)]^{1/3}r`$ where $`K\mathrm{\Lambda }=(8\pi /5)(\mu _m/k)^3P^2m_pf\mathrm{\Lambda }(T)/\dot{m}_s=5.2\times 10^{38}r_{kpc}^{3.6}`$. Here $`\mu _m=0.62`$ is the molecular weight, $`f(\mu _m)=0.58`$, and $`P1.9\times 10^{10}r_{kpc}^{1.8}`$ dyne is the approximate interstellar pressure variation. The flow into each site is $`\dot{m}_s=10^6`$ $`M_{}`$ yr<sup>-1</sup> since line emission from $`10^6`$ sites is required to match the total H$`\beta `$ luminosity from luminous elliptical galaxies (Mathews & Brighenti 1999a) and approximately 1 $`M_{}`$ of interstellar gas cools each year. Taking $`T0`$ at the site center, the radius where $`T=10^6`$ K, at which the site density is $`10\rho `$, is $`r_s2.7\times 10^{18}r_{kpc}^{1.2}`$ cm and the corresponding terminal velocity is $`u_t=[4gr_s(\rho _s/\rho )/3]^{1/2}=43r_{kpc}^{0.1}`$ km s<sup>-1</sup>. The effective viscosity is then $`\mu _sn_sm_su_tr_s=q\rho u_tr_s`$ where the product of the space density of sites $`n_s`$ and the local site accretion rate $`\dot{m}_s`$ is set equal to $`q\rho /t_{do}`$ and we assume $`m_s/\dot{m}_s=t_{do}`$. Combining all these results, the effective viscosity for the ensemble of sinking cooling sites is roughly $`\mu _s1.2r_{kpc}^{0.5}`$ gm cm<sup>-1</sup> s<sup>-1</sup>. This viscosity is small within most of the bright X-ray core of the galaxy. Although $`\mu _s`$ exceeds $`\mu _{sn}`$ for r
>50
>𝑟50r\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}50 kpc, there is little or no optical line emission from these distant regions.
Obviously, each of these estimates for $`\mu _{}`$, $`\mu _{sn}`$ and $`\mu _s`$ is highly uncertain, perhaps good to an order of magnitude. In addition to these possible sources of turbulent viscosity, the interstellar gas could become turbulent due to continuous mergers of gas-rich smaller galaxies having turbulent wakes, but there is little optical evidence for sufficiently numerous merging galaxies. The torquing action of magnetic stresses can also propagate angular momentum outward, masquerading as viscosity. Finally, we note the considerable energy available in differential shear that could be converted to turbulence if the rotating flow is globally unstable. If the characteristic change of azimuthal velocity over length $`\mathrm{}_\varphi `$ is $`\mathrm{\Delta }v_\varphi |v_\varphi /R|\mathrm{}_\varphi `$, then the effective viscosity is $`\mu _\varphi \rho |v_\varphi /R|\mathrm{}_\varphi ^2`$. From Figure 2a, $`|v_\varphi /R|40`$ km s<sup>-1</sup> kpc<sup>-1</sup> so $`\mu _\varphi 1000r_{kpc}^{1.8}\mathrm{}_{\varphi ,kpc}^2`$, which could be large depending on the unknown scale length $`\mathrm{}_\varphi `$.
### 3.2 Viscous Terms
In rotating cooling flows the flow velocity in the meridional plane is very subsonic, typically more than an order of magnitude smaller than the azimuthal velocity (Paper I). Therefore, only the azimuthal velocity component $`v_\varphi `$ is relevant to our investigation of the importance of the viscous transport of angular momentum in circularizing the X-ray images of elliptical galaxies. In our numerical solutions, we implement the contribution of viscosity by solving the conservation equation for angular momentum density $`\rho h\rho Rv_\varphi `$. The rate of change of angular momentum density due to viscous effects alone is given by
$$\left(\frac{(\rho h)}{t}\right)_{vis}=\mu \frac{^2h}{R^2}+\mu \frac{^2h}{z^2}\frac{\mu }{R}\frac{h}{R}$$
$$\frac{2h}{R}\frac{\mu }{R}+\frac{h}{R}\frac{\mu }{R}+\frac{h}{z}\frac{\mu }{z}$$
where $`\mu =\mu (R,z)`$ is a spatially variable viscosity. We have incorporated these viscous terms into the Eulerian ZEUS2D code (Stone & Norman 1992) as a separate computational process. In addition to the usual consideration of “source” and “transport” steps in the Eulerian numerical procedure, we modify $`v_\varphi `$ for viscous transport by solving the equation above using a standard operator-splitting procedure. The viscosity terms in the angular momentum equation above correspond to a modified two-dimensional diffusion. We solve these terms at each timestep using a second order alternating-direction implicit (ADI) difference scheme (Press et al. 1992) on the 2D logarithmic grid. Angular momentum diffuses spatially because of viscous effects. When there is no cooling dropout, our numerical procedure conserves total angular momentum but also allows it to flow out of the computational grid at the boundaries.
The equation for thermal energy conservation must also be corrected for the contribution of viscous dissipation. An additional numerical step is applied to the thermal energy equation to solve for the following dissipative heating term:
$$\left(\rho \frac{d\epsilon }{dt}\right)_{vis}=\mu R^2\left\{\left[\frac{}{R}\left(\frac{h}{R^2}\right)\right]^2+\left[\frac{}{z}\left(\frac{h}{R^2}\right)\right]^2\right\}$$
where $`\epsilon =3kT/2\mu _mm_p`$ is the specific thermal energy. The remaining (non-viscous) terms in the thermal energy equation, and the complete equation of motion, are discussed in Brighenti & Mathews (1999a).
### 3.3 Effect of Viscosity on X-ray Images
In general we find that the circularization of X-ray images of elliptical galaxies is easier to achieve with cooling dropout than with viscous redistribution of angular momentum. For this reason we begin our discussion of viscous rotating cooling flows with the supernova viscosity, the largest viscosity expected from the estimates above. In this case $`\mu =\mu _{sn}(r)`$ where $`r=[R^2+z^2]^{1/2}`$ is the radial distance from the galactic center. Figure 4a shows the X-ray image of our rotating galaxy viewed along the equatorial plane at time $`t_n=13`$ Gyrs but now including viscous terms based on $`\mu _{sn}=16.2r_{kpc}^{1.4}`$ gm cm<sup>-1</sup> s<sup>-1</sup>. The solution shown in Figure 4a also includes a mass dropout coefficient $`q=1`$ and the additional surface brightness $`\mathrm{\Sigma }_x(R,z)`$ contribution from the collective emission from cooling sites. The contours in Figure 4a are seen to be almost identical to those shown in Figure 1c, the same calculation without viscosity. The velocity field (not shown) is also nearly the same as that in Figure 2b and the spatial distribution of dropout gas in this solution $`\rho _{do}(R,z)`$ is almost identical to the same non-viscous $`q=1`$ solution shown in Figure 3. It is clear that the additional small viscosity due to stellar mass loss and moving cooling sites, $`\mu _{}`$ and $`\mu _s`$, will have almost no effect on the solutions. We conclude that the expected level of turbulent viscosity is insufficient to circularize the X-ray images of cooling flows in elliptical galaxies.
However, our estimates of the viscosity are extremely crude, and it is possible that the real value is considerably larger. To explore this idea, we show in Figure 4b surface brightness $`\mathrm{\Sigma }_x(R,z)`$ contours at $`t_n=13`$ Gyrs computed with a greatly enhanced viscosity, $`\mu =100r_{kpc}^{1.4}`$, which still has the decreasing radial dependence expected if the viscosity has a stellar origin. The surface brightness contours in Figure 4b, also computed with $`q=1`$ and dropout emission, are noticeably rounder in the bright central region. In addition, a projection of the X-ray emission from this solution viewed at inclination $`i=60^o`$ (Figure 4c) has contours that are considerably more circular and compares favorably with the observations of Hanlan & Bregman (1999).
Finally, we consider solutions with very large and uniform turbulent viscosity. Our estimates of the sources of viscosity – $`\mu _{sn}`$, $`\mu _{}`$, and $`\mu _s`$ – are independent of galactic shear. However, as we speculated earlier, it is possible that these (or other) non-linear perturbations could unleash the considerable energy density in the shear flow to maintain a high level of turbulence throughout the flow. Strong, approximately uniform distributed turbulence might result from continuing mergers of small gas-rich galaxies penetrating through the cooling flow. But this hypothesis seems doubtful since there is little or no evidence for this type of continuous galactic bombardment in either X-ray or optical images. If many on-going small mergers were required to establish global interstellar turbulence, numerous small, gas-rich galaxies should be apparent in optical images of large elliptical galaxies within the optical image and beyond.
Another obvious motivation for exploring larger turbulent viscosities is the persistence of X-ray flattening in the 10 - 15 kpc region apparent in Figures 1b, 1d, 4b and 4c, which may still exceed X-ray ellipticities in Einstein HRI and ROSAT HRI images. Isophotal flattening (and twists) have been reported by Boute & Canizares (1998) in some large ellipticals, but this effect occurs at much larger radii and is probably related to misalignments of the dark halo and stellar potentials which we do not consider here.
In Figure 5 we show equator-on views of the X-ray surface brightness at 13 Gyrs for solutions computed with large and spatially uniform viscosities, $`\mu =5`$ and 20 gm cm<sup>-1</sup> s<sup>-1</sup> respectively. Both solutions include dropout with $`q=1`$ and emission produced by cooling sites. Clearly, the largest turbulent viscosity ($`\mu =20`$) can circularize images both in the 10 - 15 kpc region and near the bright center. Turbulent viscosity has very little influence on the current ROSAT luminosity within the optical image of the E2 galaxy: $`L_x=9.1`$, 11.0, 10.0, and $`11.5\times 10^{40}`$ ergs s<sup>-1</sup> for $`q=1`$ models with the four viscosities we consider, $`\mu =16.2r_{kpc}^{1.4}`$, $`100r_{kpc}^{1.4}`$, 5, and 10 respectively.
As we mentioned earlier, the fraction of interstellar gas that derives from stellar mass loss decreases with galactic radius. At r
>10
>𝑟10r\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}10 kpc a larger fraction of the gas has flowed in from the halo (Brighenti & Mathews 1999a) and is not created solely by stellar mass loss as we have assumed. The angular momentum density in gas flowing in from the outer halo may be different (and possibly larger) than that of the stars. However, if most of this extended gas was expelled from the large elliptical at early times, driven by a Type II supernova wind, its angular momentum may differ little from that of the galactic stars. But, even in this scenario, if this extended gas has been tidally influenced by other nearby galaxies, as may be common, when it flows in from the outer halo its specific angular momentum could differ from the galactic stars both in magnitude and direction.
To explore the possible influence of inflowing circumgalactic gas on the apparent X-ray shape of galactic cooling flows, we performed several idealized calculations in which additional, non-rotating circumgalactic gas flows in from the galactic halo. In these models (not illustrated here) we used the same parameters for the circumgalactic gas as described in Brighenti & Mathews (1998). As expected, the ellipticity of the X-ray isophotes at 10 - 15 kpc is reduced by the inflowing gas, but the effect is not very large. Evidently the small but significant global spinup resulting from interaction with gas expelled from the rotating stellar system is important even at these radii.
## 4 VISCOUS DIFFUSION OF METALLICITY
If viscosity is invoked to circularize the X-ray images of elliptical galaxies or to amplify magnetic fields in the interstellar gas (Mathews & Brighenti 1997), then the level of turbulence required must be consistent with observed metallicity gradients in the interstellar gas. When X-ray data reduction procedures are applied to spatially resolved elliptical galaxies, the iron abundance in the interstellar gas has been found to decrease with increasing radius approximately as $`z_{Fe}r^{0.5}`$ (e.g., Matsushita 1997). We restrict our discussion here to iron for which current data is most plentiful. In the presence of spatial diffusion the continuity equation for iron mass is
$$\frac{\rho c}{t}+(\rho c𝐮)=(\rho Dc)$$
(1)
$$+\left[\left(\frac{\alpha _{}z_{}}{1.4}\right)+\left(\frac{\alpha _{sn}y_{Fe,Ia}}{m_{sn}}\right)\right]\rho _{}q\frac{\rho }{t_{do}}c,$$
where $`cc_{Fe}=\rho _{Fe}/\rho `$ is the concentration of iron relative to the total mass density. The abundance of iron in the hot gas relative to hydrogen mass can be found from the concentration by correcting for the mass of helium, $`z=1.4c`$. The diffusion coefficient depends on a characteristic velocity and mixing length, $`Dv_d\mathrm{}_d`$ and can be expressed directly in terms of the viscosity, $`D\mu /\rho `$. Iron enters the interstellar gas both from normal stellar mass loss and from Type Ia supernovae. For the stellar iron abundance we assume
$$z_{}(R,z)=z_{Fe,o}[1+(R/R_c)^2+(z/z_c)^2]^{0.3}$$
with central value $`z_{Fe,o}=1.77\times 10^3`$ (e.g., Arimoto et al. 1997; Ishimaru & Arimoto 1997). Each supernova expels $`y_{Fe,Ia}=0.744`$ $`M_{}`$ in iron and a total mass of $`m_{sn}=1.4`$ $`M_{}`$. With a stellar mass to light ratio $`M_t/L_B=9.14`$, we find $`\alpha _{sn}=4.85\times 10^{21}\mathrm{SNu}(t_n)(t/t_n)^1`$ s<sup>-1</sup>.
For most cases of interest viscous spatial diffusion is insufficient to flatten iron abundance gradients since the gradient is re-established on short timescales by Type Ia supernovae. The timescale for spatial diffusion of iron is $`t_{d,Fe}\rho r^2/\mu 2\times 10^9r_{kpc}^{1.6}`$ years using $`\rho 1.05\times 10^{25}r_{kpc}^{1.8}`$ gm cm<sup>-3</sup> and $`\mu =\mu _{sn}=16.2r_{kpc}^{1.4}`$ gm cm<sup>-1</sup> s<sup>-1</sup>. By comparison, the timescale required for Type Ia supernovae to restore an interstellar iron concentration of $`c=z/1.4`$ is $`t_{sn,Fe}\rho cm_{sn}/\alpha _{sn}y_{Fe,Ia}\rho _{}`$. Adopting $`z=0.5`$, $`z_{Fe,}=1.77\times 10^3`$, $`\rho =1.05\times 10^{25}r_{kpc}^{1.8}`$ gm cm<sup>-3</sup>, $`\rho _{}=4.35\times 10^{22}r_{kpc}^3`$ gm cm<sup>-3</sup>, and SNu$`(t_n)=0.03`$, we find $`t_{sn,Fe}6\times 10^7r_{kpc}^{1.2}`$ years. Therefore $`t_dt_{sn,Fe}`$ and supernovae can maintain interstellar iron abundance gradients even in the presence of supernova-driven turbulence.
To test this more quantitatively, we have solved the iron concentration equation in several evolving cooling flows. In cylindrical geometry the divergence term in the iron continuity equation (1) is
$$(\rho Dc)=\frac{1}{R}\frac{}{R}\left(R\rho D\frac{c}{R}\right)+\frac{}{z}\left(\rho D\frac{c}{z}\right).$$
This diffusion equation was solved using a second-order ADI method similar to that used for the viscosity terms in the equation of motion.
Figure 6 compares iron abundance profiles in two cooling flows at $`t=t_n`$ along the $`R`$ and $`z`$ axes for both the gas and stars. For the flow with $`q=1`$, $`\mu =\mu _{sn}=16.2r_{kpc}^{1.4}`$ gm cm<sup>-1</sup> s<sup>-1</sup> and $`D=\mu /\rho `$, the diffusion has almost no effect on the iron distribution in the hot gas, in agreement with our simple estimate above. Even when the viscosity is increased to $`\mu =100`$ throughout the computational grid, also shown in Figure 6, an iron abundance gradient is still present, but is beginning to flatten, particularly at the outer edge of the galactic stellar distribution at $`\mathrm{log}R\mathrm{log}z2`$. This indicates that observed interstellar iron abundance gradients can be maintained only if the average turbulent viscosity is μ
<100
<𝜇100\mu\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}100 gm cm<sup>-1</sup> s<sup>-1</sup>. To investigate further the importance of the source terms on the metallicity gradient, we performed an additional exploratory computation (with $`q=1`$ and $`\mu =100`$) in which the iron source terms in equation (1) were set to zero beginning at time $`t=10`$ Gyrs. By time $`t=11`$ Gyrs in this calculation the iron gradient in the interstellar gas was completely leveled, resulting in uniform $`z_{Fe}(R,z)`$ throughout the computational grid. In this unrealistic test calculation our solution of the concentration diffusion term above behaved as expected.
The diffusive mixing of supernova enrichment in the interstellar gas is relevant to the possibility that locally metal-rich regions may cool by X-radiation and dropout faster than the rest of the cooling flow. We estimate the metallicity diffusion time by equating the diffusion length from each supernova blast wave bubble to the mean distance between bubbles. The number of Type Ia supernova bubbles in a large elliptical galaxy at any time, $`𝒩_{bubb}\nu _{sn}t_{bubb}1.5\times 10^3(t_{bubb}/10^6\mathrm{yrs})`$, is quite small. Here $`t_{bubb}10^6`$ yrs is a typical lifetime of a supernova bubble to destruction by buoyant mixing (Mathews 1990). However, the iron deposited in the cooling flow by Type Ia supernova persists long after the thermal energy in the hot blast wave bubble has dissipated. To estimate the efficiency that supernova iron mixes throughout the interstellar gas by turbulent diffusion, we imagine a time $`t=0`$ at which the supernova explosions begin. Then the mixing time $`t_{mix}`$ for supernova iron to diffuse between supernova sites is found by equating the decreasing distance between supernova bubbles at $`t>0`$ to the distance that iron has diffused by turbulent mixing from any particular site; $`t_{mix}`$ must then be compared to the local cooling flow time. As Type Ia supernovae events increase with time $`t`$, the mean separation between sites is $`\delta _{bubb}(t\nu _{sn}\rho _{}/M_t)^{1/3}36r_{kpc}t_{yr}^{1/3}`$ kpc, estimated with $`\rho _{}=4.35\times 10^{22}r_{kpc}^3`$ gm cm<sup>-3</sup>. This can be compared with the distance $`\delta _{diff}`$ that diffusion mixes metals from an individual supernova bubble in the same time $`t`$: $`\delta _{diff}=(\mu t/\rho )^{1/2}2.3\times 10^5t_{yr}^{0.5}r_{kpc}^{0.9}`$ kpc where we assume $`\mu =\mu _{sn}=16.2r_{kpc}^{1.4}`$ gm cm<sup>-1</sup> s<sup>-1</sup> and $`\rho =1.05\times 10^{25}r_{kpc}^{1.8}`$ gm cm<sup>-3</sup>. Diffusive mixing of supernova enrichment in the interstellar medium occurs in time $`t=t_{mix}`$ found by setting $`\delta _{diff}\delta _{bubb}`$, which gives $`t_{mix}3\times 10^7r_{kpc}^{0.96}`$ yrs. This mixing time is comparable to the radial flow time in the cooling flow, $`t_{flow}r/u5\times 10^7r_{kpc}[u/20(\mathrm{km}/\mathrm{s})]^1`$ yrs. We conclude that supernova enrichment products are well mixed throughout the cooling flow for the Type Ia supernova rate we assume here; for turbulent viscosities μ
>μsn
>𝜇subscript𝜇𝑠𝑛\mu\mathrel{\mathchoice{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.0pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}\mu_{sn}, mixing is guaranteed. In a cooling flow the local radiative cooling time is comparable to the flow time, so iron produced by supernovae disappears from the flow by mass dropout in time $`t_{flow}`$. The complete absence of diffusive mixing, as envisioned by Fujita, Fukumoto & Okoshi (1996; 1997), in which unmixed, metal-rich supernova bubbles and stellar ejecta cool rapidly by enhanced X-ray line emission and differentially drop out from the flow may not be realistic, but we have not explicitly ruled out this possiblilty.
## 5 SUMMARY AND CONCLUSIONS
In this paper we have addressed the curious absence of X-ray disks in rotating, luminous elliptical galaxies. In our previous hydrodynamic studies of rotating cooling flows (Paper I; Brighenti & Mathews 1997) both angular momentum and mass were conserved and the resulting X-ray images were dramatically flattened out to 1-2 effective radii when viewed perpendicular to the axis of rotation. These results appear to be inconsistent with the more circular X-ray images found by observations with Einstein HRI (Fabbiano, Kim & Trinchieri 1992) and ROSAT HRI (Hanlan & Bregman 1999). Nevertheless, the X-ray shapes of elliptical galaxies are difficult to determine from the observations and we may need to wait for the next generation of X-ray telescopes with higher spatial resolution for definitive determinations of X-ray ellipticities. It is likely that interstellar rotational flattening will eventually be observed.
In the preceding discussion we have explored the circularizing influence of mass dropout and turbulent transfer of angular momentum on the X-ray images of rotating galactic cooling flows. Strong theoretical and observational arguments support both of these possibilities. Localized radiative cooling leading to mass dropout throughout the cooling flow is indicated by the relatively low masses of central black holes and adjacent nuclear regions when compared to the mass of interstellar gas that is expected to cool over cosmic time. Further evidence for cooling dropout is provided by optical line emission within $`R_e`$ and the evidence in this same region for a young stellar population having masses that extend up to 1 - 2 $`M_{}`$ but not beyond (Ferland, Fabian, & Johnstone 1994; Mathews & Brighenti 1999a). Interstellar turbulence is also very likely in elliptical galaxies, generated by mass transport associated with stellar mass loss, Type Ia supernovae and mass dropout. Indirect evidence for cooling flow turbulence is provided by interstellar magnetic fields of several $`\mu `$G at $`10R_e`$, which are typically observed in elliptical galaxies having double radio sources. Fields of this magnitude can be understood as originating from small seed fields ejected from mass-losing stars, followed by subsequent field amplification by an interstellar turbulent dynamo mechanism (Moss & Shukurov 1996; Mathews & Brighenti 1997). Due to the field concentration associated with the inward, converging motion of the cooling flow, interstellar magnetic fields are further magnified and are expected to be particularly strong, perhaps exceeding equipartition values, in the central regions of cooling flows. Owen and Eilek (1998) find fields of 10 - 100 $`\mu `$G within $`r=50`$ kpc in the bright cD galaxy NGC 6166.
Elliptical galaxies of lower optical luminosity LB
<3×1010
<subscript𝐿𝐵3superscript1010L_{B}\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}3\times 10^{10} $`L_{B,}`$ differ in many qualitative ways from the more massive elliptical galaxies similar to the one we have studied here (Faber et al. 1997). Low luminosity elliptical galaxies rotate faster and are rotationally flattened. We have shown (Brighenti & Mathews 1997) that rotating, non-dropout cooling flows in these galaxies also form massive, extended disks of cooled gas. However, when $`q=1`$ mass dropout is included in these calculations, disks of cold HI (or H<sub>2</sub>) do not form from the rotating cooling flow. There is observational evidence for extended disks of cold HI gas in some low luminosity elliptical galaxies, but these HI disks often extend far beyond the optical images of the galaxies, indicating that they were created by a (possibly very old) merging event rather than by cooling flow dropout (Oosterloo, Morganti & Sadler 1999a; 1999b). In some cases the cold gaseous disks have central holes in neutral hydrogen replaced with HII emission (Oosterloo, Morganti & Sadler 1999), suggesting a more complete conversion into stars in this part of the disk where gas pressures and densities are highest and star formation (and ionization by PAGB stellar UV) should be most efficient. In general, systematically larger H$`\beta `$ features are observed in the stellar spectra of elliptical galaxies with LB
<3×1010
<subscript𝐿𝐵3superscript1010L_{B}\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\mathsurround 0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}3\times 10^{10} $`L_{B,}`$, indicating considerable recent star formation (de Jong & Davies 1997). Unfortunately, thermal X-ray emission from the hot interstellar gas in these low-luminosity, low-$`\mathrm{\Sigma }_x`$ ellipticals is masked by the collective emission from low mass X-ray binaries. A detection of X-ray rotational flattening in their cooling flows will require spectral separation of the hot gas and stellar X-rays with high spatial resolution, clearly a job for the next generation of X-ray telescopes.
When mass dropout and turbulence are considered in rotating cooling flows inside luminous elliptical galaxies, we find that:
* Mass dropout alone strongly circularizes the X-ray images in rotating cooling flows since gas is removed from the flow before it has moved very far from its point of origin (stellar mass loss) toward the axis of rotation.
* Conversely, the absence of strong rotational flattening in X-ray images of elliptical galaxies is persuasive evidence for distributed mass dropout.
* In the presence of mass dropout by radiative cooling with $`q=1`$, no cold gaseous disks form on the equatorial plane having radii larger than our innermost grid size, $`150`$ pc. This is due to the sensitivity of radiative cooling (and associated mass dropout) to the local density, $`(\rho /t)_{do}q\rho /t_{do}\rho ^2`$.
* The spatial distribution of cooled gas mass is markedly flatter than that of the old stellar population, especially near the galactic center. Any optical signature of a younger stellar population formed from the cooled gas, such as the H$`\beta `$ index, should also exhibit a significantly higher ellipticity than that of the older, background stars.
* The estimated viscosity from known sources of interstellar turbulence is dominated by motions induced by Type Ia supernovae; stellar mass loss and mass dropout are smaller sources of interstellar turbulent viscosity.
* Supernova-induced turbulent viscosity is insufficient to circularize the X-ray appearance of rotating elliptical galaxies; the viscous effects of stellar mass loss and cooling dropout on the X-ray images are even smaller.
* Much larger, spatially uniform turbulent viscosities can circularize the X-ray isophotes throughout the galaxy; such turbulent viscosities could result from rotational shear instabilities, but we have not demonstrated this instability here.
* The spatial diffusion of interstellar iron due to turbulent viscosity does not appreciably reduce the interstellar iron abundance gradients observed in bright elliptical galaxies, provided the viscosity does not exceed $`100`$ gm cm<sup>-1</sup> s<sup>-1</sup>. This follows since the observed interstellar iron gradients are re-established on short time scales by Type Ia supernovae and, to a lesser extent, by stellar mass loss.
Studies of the evolution of hot gas in elliptical galaxies at UC Santa Cruz are supported by NASA grant NAG 5-3060 and NSF grant AST-9802994 for which we are very grateful. FB is supported in part by Grant MURST-Cofin 98. |
warning/0003/math0003180.html | ar5iv | text | # On complete arcs arising from plane curves
## 1. Introduction and statement of results
A $`(k,d)`$-arc $`𝒦`$ in the projective plane $`𝐏^2(𝐅_q)`$, $`𝐅_q`$ being the finite field with $`q`$ elements, is a set of $`k`$ elements such that no line in $`𝐏^2(𝐅_q)`$ meets $`𝒦`$ in more than $`d`$ points. The $`(k,d)`$-arc is called complete if it is not contained in a $`(k+1,d)`$-arc. For basic facts on arcs the reader is refered to \[7, Ch. 12\] (see also the references therein), \[10, Sec. 5\], , and .
A natural example of a $`(k,d)`$-arc is the set $`𝒳(𝐅_q)`$ of $`𝐅_q`$-rational points of a plane curve $`𝒳`$ without linear components and defined over $`𝐅_q`$, where $`k=\mathrm{\#}𝒳(𝐅_q)`$ and $`d`$ is the degree of $`𝒳`$. As a matter of terminology, we shall say that $`𝒳`$ has the arc property whenever $`𝒳(𝐅_q)`$ is a complete $`(k,d)`$-arc with $`k`$ and $`d`$ as above. As a matter of fact, the interplay between the theory of algebraic curves and finite geometries was initiated by Segre around 1955. In (see also \[7, Sec. 10.4\]) he established an upper bound for the second largest size that a complete $`(k,2)`$-arc in $`𝐏^2(𝐅_q)`$ can have. He proved this result by applying the Hasse-Weil upper bound to the non-singular model of the envelope associated to $`(k,2)`$-arcs. For further results on these arcs see Hirschfeld and Korchmáros’ papers and . For other applications of curves to finite geometries see the surveys and .
In this paper we are concerning with the problem of determining plane curves having the arc property. This was asked around 1988 by Hirschfeld and Voloch \[11, Problem III\]. Only few examples of such curves are known. Among them we have the irreducibles conics in odd characteristic \[7, Ch. 8\], certain cubics \[11, Sect. 5\], \[5, Sect. 6\], and Hermitian curves \[7, Lemma 7.20\].
For $`q`$ a square, the Hermitian curve $`𝒳_{\sqrt{q}+1}`$ can be defined by
(1.1)
$$X^{\sqrt{q}+1}+Y^{\sqrt{q}+1}+Z^{\sqrt{q}+1}=0.$$
Then $`𝒳_{\sqrt{q}+1}(𝐅_q)`$ is a complete $`(q\sqrt{q}+1,\sqrt{q}+1)`$-arc (loc. cit.). The completeness property means that for each $`P𝐏^2(𝐅_q)𝒳_{\sqrt{q}+1}(𝐅_q)`$ there exists a $`𝐅_q`$-rational line $`\mathrm{}`$ such that $`\mathrm{\#}\mathrm{}𝒳_{\sqrt{q}+1}(𝐅_q)=\sqrt{q}+1`$. This property can be easily shown by using the special feature of Eq. (1.1) (see the proof of Theorem 1.1(1)). Moreover, this property is also a consequence of the fact that the image of $`P𝒳_{\sqrt{q}+1}`$ by the $`𝐅_q`$-Frobenius morphism lies on the tangent line of $`𝒳_{\sqrt{q}+1}`$ at $`P`$ (see the proof of Theorem 1.2). A plane curve satisfying the above property for general points is called $`𝐅_q`$-Frobenius non-classical , .
The Hermitian curve is a member of the family of Fermat curves $`𝒳_d(a,b)`$ defined by
(1.2)
$$aX^d+bY^d+Z^d=0,$$
where $`d=(q1)/(q^{}1)`$, $`q=p^n`$, $`q^{}=p^e`$, $`p:=\mathrm{char}(𝐅_q)`$ such that $`e<n`$ and $`e|n`$, and where $`a,b𝐅_q^{}^{}:=𝐅_q^{}\{0\}`$.
The first aim of this paper is to extend the arc property of the Hermitian curve to the curves $`𝒳_n(a,b)`$ above as well as to the Fermat curve $`𝒳_{q1}`$ defined by
(1.3)
$$X^{q1}+Y^{q1}=2Z^{q1},$$
provided that $`p3`$.
###### Theorem 1.1.
For $`p3`$, the following statements hold:
1. the set of $`𝐅_q`$-rational points of the Fermat curve defined by Eq. (1.2) is a complete $`(k,d)`$-arc, where $`k=d(qd+2);`$
2. the set of $`𝐅_q`$-rational points of the Fermat curve defined by Eq. (1.3) is a complete $`((q1)^2,q1)`$-arc.
The arc in part (2) of this theorem is maximal among $`(k,q1)`$-arcs for which there exists an external line, see Remark 3.3. On the other hand, it seems that the arcs in part (1) are new.
We notice that the curves in Theorem 1.1 are among the Fermat curves having a large number of $`𝐅_q`$-rational points . We also notice that the hypothesis $`p3`$ is necessary, see Remark 3.2.
The second aim of this paper is to show that certain $`𝐅_q`$-Frobenius non-classical plane curves do satisfy the arc property.
###### Theorem 1.2.
Let $`𝒳`$ be a non-singular $`𝐅_q`$-Frobenius non-classical plane curve of degree $`d`$. Let $`ϵ`$ be the order of contact of $`𝒳`$ with the tangent at a general point. If
(1.4)
$$d(d1)<(q+1)ϵ,$$
then the set of $`𝐅_q`$-rational points of $`𝒳`$ is a complete $`(k,d)`$-arc, where $`k=\mathrm{\#}𝒳(𝐅_q)=d(qd+2)`$.
The Fermat curve in Theorem 1.1(2) is $`𝐅_q`$-classical by \[4, Thm. 2\]. Hence the hypothesis of being $`𝐅_q`$-Frobenius non-classical in Theorem 1.2 is not necessary. We observe that $`d(d1)/ϵ`$ is the degree of the dual curve of $`𝒳`$; see Remark 2.3.
For the Hermitian curve $`𝒳_{\sqrt{q}+1}`$, $`ϵ=\sqrt{q}`$ (see e.g. ) and therefore it satisfies (1.4) in Theorem 1.2. The Fermat curves in Theorem 1.1(1) are $`𝐅_q`$-Frobenius non-classical, see \[4, Thm. 2\]. For these curves, $`ϵ=q^{}`$ \[6, Thm. 2\], and so they satisfy (1.4) if and only if they are Hermitian curves. So far, we could not find examples of non-singular $`𝐅_q`$-Frobenius non-classical curves fulfilling (1.4) which are not $`𝐅_q`$-isomorphic to Hermitian curves; see Remark 4.3.
## 2. Frobenius non-classical planes curves
The study of Frobenius non-classical curves was initiated by Hefez and Voloch based on a fundamental paper by Stöhr and Voloch , where an approach to the Hasse-Weil bound was given.
In this paper we only consider irreducible non-linear plane curves defined over $`𝐅_q`$. Let $`𝒳𝐏^2(\overline{𝐅}_q)`$ be such a curve. For $`i=0,1,2`$, let $`x_i`$ be the coordinates functions of $`𝐏^2(\overline{𝐅}_q)`$ on $`𝒳`$. Let $`t`$ be a separating variable of $`𝐅_q(𝒳)|𝐅_q`$ and denote by $`D^i=D_t^i`$ the $`i`$-th Hasse derivative on $`𝒳`$. The order sequence of $`𝒳`$ (see \[15, p. 5\]) are the numbers $`0,1`$ and $`ϵ=ϵ(𝒳)`$, where $`ϵ>1`$ is the least integer such that
$$\mathrm{det}\left(\begin{array}{ccc}x_0& x_1& x_2\\ D^1x_0& D^1x_1& D^1x_2\\ D^ϵx_0& D^ϵx_1& D^ϵx_2\end{array}\right)0.$$
Geometrically, the numbers 0, 1 and $`ϵ`$ represent all the possible intersection multiplicities of the curve $`𝒳`$ with lines in $`𝐏^2(\overline{𝐅}_q)`$ at general points.
The $`𝐅_q`$-Frobenius order sequence of $`𝒳`$ (see \[15, p. 9\]) are the numbers 0 and $`\nu =\nu (𝒳,q)`$, where $`\nu >0`$ is the least integer such that
(2.1)
$$\mathrm{det}\left(\begin{array}{ccc}x_0^q& x_1^q& x_2^q\\ x_0& x_1& x_2\\ D^\nu x_0& D^\nu x_1& D^\nu x_2\end{array}\right)0.$$
We have that $`\nu \{1,ϵ\}`$ \[15, Prop. 2.1\]. The plane curve $`𝒳`$ is called $`𝐅_q`$-Frobenius non-classical if $`\nu =ϵ`$ (or equivalently if $`\nu >1`$).
###### Remark 2.1.
From Eq. (2.1) follows that $`𝒳`$ is $`𝐅_q`$-Frobenius non-classical if and only if $`𝐅r_𝒳(P)T_P𝒳`$ for all non-singular points $`P𝒳`$, where $`𝐅r_𝒳`$ is the $`𝐅_q`$-Frobenius morphism on $`𝒳`$ and $`T_P𝒳`$ is the tangent line to $`𝒳`$ at $`P`$.
###### Remark 2.2.
Let $`𝒳`$ be a curve as above.
(i) If $`ϵ(𝒳)>2`$ then it is a power of $`p`$ \[3, Prop. 2\].
(ii) If $`𝒳`$ is $`𝐅_q`$-Frobenius non-classical and $`p>2`$, then $`ϵ(𝒳)>2`$ \[6, Prop. 1\].
(iii) If $`𝒳`$ is $`𝐅_q`$-Frobenius non-classical, then $`ϵ(𝒳)q`$ \[6, p. 266\]. If in addition $`𝒳`$ is non-singular, then $`ϵ(𝒳)\sqrt{q}`$ \[6, Prop. 6\].
(iv) Let $`𝒳`$ be non-singular $`𝐅_q`$-Frobenius non-classical plane curve. Let $`d`$ the degree of $`d`$ and suppose that $`ϵ=ϵ(𝒳)>2`$. Then \[6, Props. 5, 6\]
$$\sqrt{q}+1d(q1)/(ϵ1).$$
###### Remark 2.3.
Let $`𝒳`$ be a non-singular plane curve of degree $`d2`$, $`𝒳^{}`$ the dual curve of $`𝒳`$ and $`T:𝒳𝒳^{}`$ the dual map; i.e, $`T(P)=T_P𝒳`$. Then $`d(d1)=\mathrm{deg}(T)d^{}`$, where $`d^{}`$ is the degree of $`𝒳^{}`$; this follows e.g. from \[12, Lemma 4.3\]. Now by a result of Kaji \[13, Cor. 4.5\], $`T`$ is purely inseparable. If in addition, $`𝒳`$ is $`𝐅_q`$-Frobenius non-classical, then $`\mathrm{deg}(T)=ϵ(𝒳)`$ \[6, Props. 3 and 4\]. Therefore, in Theorem 1.2 we look for $`𝐅_q`$-Frobenius non-classical plane curves such that the degree of its dual curve is upper bounded by $`(q+1)`$.
Let $`F=F(X_0,X_1,X_2)=0`$ be the equation of $`𝒳`$ over $`𝐅_q`$. From \[3, Thm. 1\], $`ϵ=ϵ(𝒳)>2`$ if and only if there exist homogeneous polynomials $`H,P_0,P_1,P_2𝐅_q[X_0,X_1,X_2]\{0\}`$ such that
(2.2)
$$FH=X_0P_0^ϵ+X_1P_1^ϵ+X_2P_2^ϵ.$$
We have that $`H=1`$ if $`𝒳`$ is non-singular (loc. cit, p. 462). Now it is easy to see that $`𝒳`$ is $`𝐅_q`$-non-classical if and only if there exists $`H_1𝐅_q[X_0,X_1,X_2]`$ such that
(2.3)
$$FH_1=X_0^{q/ϵ}P_0+X_1^{q/ϵ}P_1+X_2^{q/ϵ}P_2.$$
Finally, we mention a formula for the precise number of $`𝐅_q`$-rational points of non-singular $`𝐅_q`$-Frobenius non-classical curves.
###### Lemma 2.4.
(\[6, Thm. 1\]) Let $`𝒳`$ be a plane non-singular $`𝐅_q`$-Frobenius non-classical curve of degree $`d`$. Then
$$\mathrm{\#}𝒳(𝐅_q)=d(qd+2).$$
## 3. Proof of Theorem 1.1
(1) That $`\mathrm{\#}𝒳_d(a,b)=d(qd+2)`$ is well known, see e.g. \[4, p. 354\]. (This result also follows from \[4, Thm. 2\] and Lemma 2.4).
Next, for $`P𝐏^2(𝐅_q)𝒳_d(𝐅_q)`$, we will show that there exists a $`𝐅_q`$-rational line $`\mathrm{}`$ which passes through $`P`$ and intersects the curve $`𝒳_d(a,b)`$ in $`d`$ distinct $`𝐅_q`$-rational points. We recall the following easy fact.
###### Claim 3.1.
Let $`A,B𝐅_q^{}^{}`$. Then the equation $`AX^d+B=0`$ has $`d`$ distinct solutions in $`𝐅_q`$.
###### Proof.
Since $`p`$ does not divide $`d`$, the equation has $`d`$ solutions in $`\overline{𝐅}_q`$. If $`x`$ is a solution, then $`x^{q1}=1`$, as $`d=(q1)/(q^{}1)`$, and hence $`x𝐅_q`$. ∎
We consider four cases:
Case 1: $`P=(\alpha :\beta :0)`$. Let $`\mathrm{}:Z=0`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_d(a,b)=\{(\lambda :1:0):a\lambda ^d+b=0\}.$$
Case 2: $`P=(\alpha :\beta :1)`$ and $`a\alpha ^d+10`$. Let $`\mathrm{}:X=\alpha Z`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_d(a,b)=\{(\alpha :\lambda :1):b\lambda ^d+a\alpha ^d+1=0\}.$$
Case 3: $`P=(\alpha :\beta :1)`$ and $`b\beta ^d+10`$. Let $`\mathrm{}:Y=\beta Z`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_d(a,b)=\{(\lambda :\beta :1):a\lambda ^d+b\beta ^d+1=0\}.$$
Case 4: $`P=(\alpha :\beta :1)`$ and $`a\alpha ^d+1=b\beta ^d+1=0`$. Let $`\mathrm{}:\beta X=\alpha Y`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_d(a,b)=\{(\alpha ,\beta ,\beta \lambda ):\lambda ^d+2b=0\}.$$
Now by Claim 3.1 all the sets above have cardinality $`d`$ and are contained in $`𝒳_d(a,b)(𝐅_q)`$. This completes the proof of Theorem 1.1(1).
(2) We have that
$$𝒳_{q1}(𝐅_q)=\{(\alpha :\beta :\gamma )𝐏^2(𝐅_q):\alpha \beta \gamma 0\},$$
so that $`\mathrm{\#}𝒳_{q1}(𝐅_q)=q^2+q+13q=(q1)^2`$. Let $`P=(\alpha :\beta :\gamma )𝐏^2(𝐅_q)𝒳_{q1}(𝐅_q)`$. The proof of the arc property for $`𝒳_{q1}`$ follows from the following five computations.
Case 1: $`\beta =\gamma =0`$. Let $`\mathrm{}:Y=Z`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_{q1}=\{(\lambda :1:1):\lambda 𝐅_q^{}\}.$$
Case 2: $`\beta 0`$ and $`\gamma =0`$. Let $`\mathrm{}:Y=bX`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_{q1}=\{(1:b:\lambda ):\lambda 𝐅_q^{}\}.$$
Case 3: $`\alpha =\beta =0`$. Let $`\mathrm{}:X=Y`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_{q1}=\{(1:1:\lambda ):\lambda 𝐅_q^{}\}.$$
Case 4: $`\alpha 0`$ and $`\gamma =1`$. Let $`\mathrm{}:X=aZ`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_{q1}=\{(a:\lambda :1):\lambda 𝐅_q^{}\}.$$
Case 5: $`\beta 0`$ and $`\gamma =1`$. Let $`\mathrm{}:Y=bZ`$. Then $`P\mathrm{}`$ and
$$\mathrm{}𝒳_{q1}=\{(\lambda :b:1):\lambda 𝐅_q^{}\}.$$
###### Remark 3.2.
The hypothesis $`p3`$ in Theorem 1.1(1) is necessary. Indeed, consider the Fermat curve $`𝒳`$ defined by
$$X^{q1}+Y^{q1}+Z^{q1}=0,$$
where $`q`$ is a power of two. Take $`P=(1:1:1)`$. Then it is easy to see that $`\mathrm{\#}\mathrm{}𝒳<q1`$ for any $`𝐅_q`$-rational line $`\mathrm{}`$ passing through $`P`$.
###### Remark 3.3.
For a $`(k,q1)`$-arc $`𝒦`$ having an external line (i.e. a $`𝐅_q`$-rational line $`\mathrm{}`$ such that $`\mathrm{}𝒦=\mathrm{}`$), we have that $`k(q2)q+1=(q1)^2()`$ \[7, Thm. 12.40\]. Since the line $`\mathrm{}:Z=0`$ is external to the $`((q1)^2,q1)`$-arc in Theorem 1.1(2), we have that the upper bound in $`()`$ is attained by this arc.
###### Remark 3.4.
Notice that the $`((q1)^2,q1)`$-arc in Theorem 1.1(2) is the complement of three non-concurrent lines in $`𝐏^2(𝐅_q)`$. Moreover, it can be shown that any such set is a complete $`((q1)^2,q1)`$-arc which is also the set of $`𝐅_q`$-rational points of a plane curve of degree $`q1`$; it turns out that this curve is $`𝐅_q`$-isomorphic to the Fermat curve $`𝒳_{q1}`$.
###### Example 3.5.
Taking $`q=p^3`$ and $`q^{}=p`$ in Theorem 1.1(1), we have that there exists a complete $`(k,p^2+p+1)`$-arc in $`𝐏^2(𝐅_q)`$ with $`k=(p^2+p+1)(p^3p^2p+1)`$. In particular, there exists a complete $`(208,13)`$-arc in $`𝐅_{27}`$.
## 4. Proof of Theorem 1.2
Theorem 1.2 will be a consequence of the following more general result.
###### Theorem 4.1.
Let $`𝒳`$ be a $`𝐅_q`$-Frobenius non-classical (possible singular) plane curve. Let $`d`$ be the degree of $`𝒳`$ and $`k=\mathrm{\#}𝒳(𝐅_q)`$. Then $`𝒳(𝐅_q)`$ is a complete $`(k,d)`$-arc provided that
$$k>(dϵ)(q+1\mathrm{\#}S)+(d1)\mathrm{\#}S,$$
where $`S`$ is the set of singular points of $`𝒳`$ and $`ϵ`$ is as in Theorem 1.2.
###### Proof.
Suppose that $`𝒳(𝐅_q)`$ is not complete. Then there exists $`P𝐏^2(𝐅_q)𝒳(𝐅_q)`$ such that for any $`𝐅_q`$-rational line $`\mathrm{}`$ through $`P`$,
(4.1)
$$\mathrm{\#}\mathrm{}𝒳(𝐅_q)<d.$$
###### Claim 4.2.
If $`\mathrm{}`$ is a line such that (4.1) holds and that $`\mathrm{}S=\mathrm{}`$, then
$$\mathrm{\#}\mathrm{}𝒳(𝐅_q)(dϵ).$$
###### Proof.
Case 1: $`\mathrm{}𝒳𝒳(𝐅_q)`$. Let $`Q\mathrm{}𝒳𝒳(𝐅_q)`$. Then, as $`\mathrm{}`$ is $`𝐅_q`$-rational, $`𝐅r_𝒳(Q)\mathrm{}`$ and thus, as $`\mathrm{}S=\mathrm{}`$, $`\mathrm{}`$ is the tangent line of $`𝒳`$ at $`Q`$ (see Remark 2.1). Therefore
$$\mathrm{\#}\mathrm{}𝒳(𝐅_q)dϵ,$$
since $`I(𝒳,\mathrm{};Q)ϵ`$; see \[15, p.5\].
Case 2: $`\mathrm{}𝒳𝒳(𝐅_q)`$. From (4.1) and Bezout’s theorem there exists $`Q\mathrm{}𝒳(𝐅_q)`$ such that $`j(Q):=I(𝒳,\mathrm{};Q)>1`$. Then, as $`\mathrm{}S=\mathrm{}`$, $`\mathrm{}`$ is the tangent line of $`𝒳`$ at $`Q`$. We have that
$$\mathrm{\#}\mathrm{}𝒳(𝐅_q)dj(Q)+1$$
and the claim follows from the fact that $`j(Q)ϵ+1`$ \[15, Cor. 2.10\]. ∎
Now there are at most $`N\mathrm{\#}S`$ lines $`\mathrm{}^{}`$ such that $`\mathrm{}^{}S\mathrm{}`$. For each of these lines, $`\mathrm{}^{}𝒳(𝐅_q)d1`$; hence from Claim 4.2 we have
$$k(dϵ)(q+1N)+(d1)N.$$
Then $`k(dϵ)(q+1)+N(ϵ1)(dϵ)(q+1)+(ϵ1)\mathrm{\#}S`$, a contradiction. This finishs the proof of Theorem 4.1. ∎
Proof of Theorem 1.2. We have $`S=\mathrm{}`$ and so the hypothesis in Theorem 4.1 is
(4.2)
$$k>(dϵ)(q+1).$$
Since $`k=d(qd+2)`$ (see Lemma 2.4), it turns out that (4.2) is equivalent to $`d(d1)<(q+1)ϵ)`$ and the result follows.
###### Remark 4.3.
Let $`𝒳`$ be a non-singular $`𝐅_q`$-Frobenius non-classical curve of degree $`d`$. Assume $`ϵ=ϵ(𝒳)>2`$. Then from Eqs. (2.2) and (2.3), $`d=\lambda ϵ+1`$ for some $`\lambda 𝐍`$. If (1.4) holds, then
$$\sqrt{q}/ϵ\lambda <\sqrt{q/ϵ},$$
where the first inequality follows from Remark 2.2(iv).
For a concrete example take $`q=p^3`$. Then $`ϵ=p`$ by Remark 2.2(i)(iii) and so $`\sqrt{p}\lambda <p`$. Therefore $`𝒳`$ will satisfy (1.4) if $`\lambda =p1`$, i.e. if $`𝒳`$ has degree $`d=p^2p+1`$. The existence of a such curve is equivalent to the existence of polynomials $`P_0,P_1,P_2𝐅_q[X_0,X_1,X_2]`$ of degree $`p1`$ and $`H_1𝐅_q[X_0,X_1,X_2]`$ such that Eqs. (2.2) and (2.3) hold true. This seems an involved problem. On the other hand, the curve $`𝒳`$ will rise to the existence of a complete $`((p^2p+1)(p^3p^2+p+1),p^2p+1)`$-arc. Unfortunately, the existence of such an arc is not known.
Acknowledgments. The authors wish to thank J.W.P. Hirschfeld and G. Korchmáros for useful comments. This research was carried out with the support of the Italian Ministry for Research and Technology (project 40% “Strutture geometriche, combinatorie e loro applicazioni”). Part of this paper was written while Torres was visiting Perugia in February and December 1999. |
warning/0003/math0003148.html | ar5iv | text | # Untitled Document
Second-order linear differential equations with two irregular singular
points of rank three: the characteristic exponent
Wolfgang Bühring
Physikalisches Institut, Universität Heidelberg, Philosophenweg 12,
D-69120 Heidelberg, Germany
Abstract
For a second-order linear differential equation with two irregular singular points of rank three, multiple Laplace-type contour integral solutions are considered. An explicit formula in terms of the Stokes multipliers is derived for the characteristic exponent of the multiplicative solutions. The Stokes multipliers are represented by converging series with terms for which limit formulas as well as more detailed asymptotic expansions are available. Here certain new, recursively known coefficients enter , which are closely related to but different from the coefficients of the formal solutions at one of the irregular singular points of the differential equation. The coefficients of the formal solutions then appear as finite sums over subsets of the new coefficients. As a by-product, the leading exponential terms of the asymptotic behaviour of the late coefficients of the formal solutions are given, and this is a concrete example of the structural results obtained by Immink in a more general setting. The formulas displayed in this paper are not of merely theoretical interest, but they also are complete in the sense that they could be (and have been) implemented for computing accurate numerical values of the characteristic exponent, although the computational load is not small and increases with the rank of the singular point under consideration.
AMS classification: 34A20; 34A25; 34A30
Keywords: Irregular singular point; Characteristic exponent; Stokes multiplier
1. Introduction
Let us consider the differential equation
$$z^2f^{\prime \prime }+zf^{}[\underset{m=1}{\overset{6}{}}D_mz^m+L^2+\underset{m=1}{\overset{6}{}}B_mz^m]f(z)=0$$
$`(1.1)`$
with the thirteen parameters $`D_1,\mathrm{},D_6`$, $`L`$, $`B_1,\mathrm{},B_6`$, which for simplicity of presentation are assumed to be real. This differential equation has two irregular singular points, each of rank 3 (if $`D_60`$ and $`B_60`$, respectively), at the origin and at infinity. Without loss of generality, one of the parameters except $`L`$ could be set equal to $`1`$ . We assume that $`B_6`$ is positive, again for simplicity of presentation, and that $`L`$ is not negative.
At infinity, there are formal power series solutions
$$f_\mathrm{}1^{\mathrm{asy}}(z)=\mathrm{exp}(P(z))z^{\tau (1)}\underset{n=0}{\overset{\mathrm{}}{}}a_n(1)z^n,$$
$`(1.2)`$
$$f_\mathrm{}2^{\mathrm{asy}}(z)=\mathrm{exp}(P(z))z^{\tau (1)}\underset{n=0}{\overset{\mathrm{}}{}}a_n(1)z^n,$$
$`(1.3)`$
where the various quantities are determined by
$$P(z)=p_3z^3+p_2z^2+p_1z,$$
$`(1.4)`$
$$p_3=\frac{1}{3}\sqrt{B_6},p_2=\frac{1}{12}B_5/p_3,p_1=\frac{1}{6}(B_44p_2^2)/p_3,$$
$`(1.5)`$
$$\tau (\kappa )=\frac{3}{2}\frac{1}{6}\frac{B_3}{\kappa p_3}+\frac{2}{3}\frac{p_1p_2}{\kappa p_3},$$
$`(1.6)`$
$$\kappa \{1,1\}$$
$`(1.7)`$
and where the coefficients
$$a_n(\kappa )=a_n$$
$`(1.8)`$
are known recursively by
$$a_0=1$$
$`(1.9)`$
and
$$\begin{array}{cc}& 6\kappa p_3na_n=[4\kappa p_2(\tau (\kappa )+n2)+p_1^2B_2]a_{n1}\hfill \\ & +[2\kappa p_1(\tau (\kappa )+n\frac{5}{2})B_1]a_{n2}+(\tau (\kappa )+n3L)(\tau (\kappa )+n3+L)a_{n3}\hfill \\ & \underset{m=1}{\overset{6}{}}D_ma_{nm3},\hfill \end{array}$$
$`(1.10)`$
$$n=1,2,\mathrm{},(a_8=a_7=\mathrm{}=a_1=0).$$
The formal solutions are asymptotic expansions as $`z\mathrm{}`$ in appropriate sectors of the complex plane.
At the origin, there are analogical formal solutions, which may be obtained from those at infinity by the simultaneous replacements $`z1/z,B_mD_m,m=1,\mathrm{},6`$. This symmetry is our main reason for choosing just (1.1) as the standard form for a differential equation with two irregular singular points.
In the ring-shaped region $`0<|z|<\mathrm{}`$ we have (convergent) Floquet solutions $`f^{(\omega )}(z)`$ and $`f^{(\omega )}(z)`$, where
$$f^{(\omega )}(z)=z^\omega \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}c_n^{(\omega )}z^n,$$
$`(1.11)`$
which are linearly independent if $`2\omega `$ is not equal to an integer. Here the coefficients $`c_n^{(\omega )}`$ obey the recurrence relation
$$\underset{m=1}{\overset{6}{}}B_mc_{nm}^{(\omega )}+(L+\omega +n)(L+\omega +n)c_n^{(\omega )}\underset{m=1}{\overset{6}{}}D_mc_{n+m}^{(\omega )}=0,$$
$`(1.12)`$
and we want to normalize them by choosing
$$c_0^{(\omega )}=1.$$
$`(1.13)`$
The requirement that the power series converge determines, modulo 1, the possible values of the characteristic exponent (or circuit exponent or Floquet exponent) $`\omega `$. The problem to compute the characteristic exponent appears also and is best known in the context of Hill’s differential equation . There are methods to compute the characteristic exponent numerically, which require, for $`\omega =0`$, the evaluation of the infinite determinant associated with (1.12), or numerical integration of the differential equation along a suitable contour , or numerical solution of an eigen-value problem .
This paper developes an entirely different method for evaluating the characteristic exponent. We obtain an explicit formula in terms of quantities which are essentially the Stokes multipliers, and these are given explicitly as convergent series, the terms of which are represented by asymptotic expansions. Here certain recursively known coefficients enter which are closely related to but different from the coefficients of the formal solutions (1.2)-(1.3).
Although some of the other authors concerned with irregular singular points of rank larger than one \[1, 2, 4-6, 8-10, 12, 13, 15, 18, 19\] consider the general case of arbitrary rank , we here prefer to restrict our attention to rank three. This is already general enough to give an impression of what can be expected in the case of even higher rank. On the other hand, it is still simple enough so that we can, for the relevant quantities, obtain explicit expressions which are not only theoretically interesting but can also be implemented (and have been implemented) for numerical evaluation. This work may be viewed as an attempt to extend , which was useful for rank one, to the much more complicated case of higher rank.
2. Laplace contour integral solutions
We try to apply the classical method of multiple Laplace contour integral solutions and write for a solution of (1.1)
$$f(z)=z^\lambda (2\pi \text{i})^3_{C_{t_2}}_{C_{t_1}}_{C_s}\mathrm{exp}(z^3s+z^2t_1+zt_2)v(s,t_1,t_2)\text{d}s\text{d}t_1\text{d}t_2,$$
$`(2.1)`$
where a power factor with a still arbitrary parameter $`\lambda `$ has been included in view of later benefits. To derive the appropriate weight function $`v(s,t_1,t_2)`$ is a somewhat lengthy but not principally difficult procedure. We therefore give two lemmata stating the results of this procedure and postpone the proofs to a later section.
Lemma 1. The weight function $`v(s,t_1,t_2)`$ has to be a solution of the partial differential equation
$$\begin{array}{cc}& 9(s^2s_0^2)(^4v/s^4)+12(st_1s_0t_{10})(^4v/s^3t_1)+6(st_2s_0t_{20})(^4v/s^3t_2)\hfill \\ & +4(t_1^2t_{10}^2)(^4v/s^2t_1^2)+(t_2^2t_{20}^2)(^4v/s^2t_2^2)+4(t_1t_2t_{10}t_{20})(^4v/s^2t_1t_2)\hfill \\ & +([816\lambda ]s+[B_34t_{10}t_{20}])(^3v/s^3)+([524\lambda ]t_1+[B_2t_{20}^2])(^3v/s^2t_1)\hfill \\ & +([252\lambda ]t_2+B_1)(^3v/s^2t_2)+([\lambda 12]^2L^2)(^2v/s^2)D_1(^2v/st_1)\hfill \\ & D_2(^2v/st_2)+D_3(v/s)+D_4(v/t_1)+D_5(v/t_2)D_6v=0\hfill \end{array}$$
$`(2.2)`$
with two finite singular points at $`(s,t_1,t_2)=(\kappa s_0,\kappa t_{10},\kappa t_{20})`$ where
$$t_{20}=p_1,t_{10}=p_2,s_0=p_3,$$
$`(2.3)`$
and the contours for each variable have to satisfy the condition that a certain lengthy expression, bilinear in $`K`$ and $`v`$ or their partial derivatives, have the same value at the start and the end of the contour.
Lemma 2. For $`\kappa \{1,1\}`$, there are appropriate weight functions $`v=V(\kappa ;s,t_1,t_2)`$ which at the singular point $`(s,t_1,t_2)=(\kappa s_0,\kappa t_{10},\kappa t_{20})`$ have the power series expansion
$$\begin{array}{cc}& V(\kappa ;s,t_1,t_2)\hfill \\ & =\underset{m=0}{\overset{\mathrm{}}{}}\underset{n_1=0}{\overset{\mathrm{}}{}}\underset{n_2=0}{\overset{\mathrm{}}{}}A(\kappa ;m,n_1,n_2)(s\kappa s_0)^{\mu (\kappa )+m}(t_1\kappa t_{10})^{\nu _1n_1}(t_2\kappa t_{20})^{\nu _2n_2},\hfill \end{array}$$
$`(2.4)`$
where the $`A`$-coefficients satisfy a certain recurrence relation and the exponents are related by
$$3\mu (\kappa )2\nu _1\nu _2=\lambda +\tau (\kappa )6.$$
$`(2.5)`$
An appropriate set of coefficients is
$$A(\kappa ;m,n_1,n_2)=\mathrm{\Gamma }(\mu (\kappa )m)\mathrm{\Gamma }(\nu _1+n_1)\mathrm{\Gamma }(\nu _2+n_2)b(\kappa ;m,n_1,n_2),$$
$`(2.6)`$
where the new coefficients $`b(\kappa ;m,n_1,n_2)`$ are given by the recurrence relation
$$\begin{array}{cc}& 6\kappa s_0(3m2n_1n_2)b(\kappa ;m,n_1,n_2)\hfill \\ & =[(3m2n_1n_2+\tau (\kappa )3)^2L^2]b(\kappa ;m1,n_1,n_2)\hfill \\ & +[4\kappa t_{10}(3m2n_1n_2+\tau (\kappa )2)+t_{20}^2B_2]b(\kappa ;m1,n_11,n_2)\hfill \\ & +[2\kappa t_{20}(3m2n_1n_2+\tau (\kappa )\frac{5}{2})B_1]b(\kappa ;m1,n_1,n_21)\hfill \\ & D_1b(\kappa ;m2,n_11,n_2)D_2b(\kappa ;m2,n_1,n_21)D_3b(\kappa ;m2,n_1,n_2)\hfill \\ & D_4b(\kappa ;m3,n_11,n_2)D_5b(\kappa ;m3,n_1,n_21)D_6b(\kappa ;m3,n_1,n_2)\hfill \end{array}$$
$`(2.7)`$
with the initial conditions
$$b(\kappa ;0,0,0)=1,$$
$$b(\kappa ;m,n_1,n_2)=0\text{if}3m2n_1n_2=0\text{for}(m,n_1,n_2)(0,0,0),$$
$`(2.8)`$
$$(b(\kappa ;m,n_1,n_2)=0\text{if}m<0\text{or}n_1<0\text{or}n_2<0)$$
and with
$$b(\kappa ;0,n_1,n_2)=0\text{for}(n_1,n_2)(0,0),$$
$`(2.9)`$
$$b(\kappa ;m,n_1,n_2)=0\text{for}n_1+n_2>m$$
$`(2.10)`$
as a consequence. In addition, there are weight functions $`U(\kappa ;s,t_1,t_2)`$ which are analytic in $`s`$ at the respective singular point, corresponding to $`\mu (\kappa )=0,1,2`$ (without any further relation such as (2.5)).
In order to avoid unnecessary complications, we want to assume that the non-trivial exponent $`\mu `$ according to (2.5) is not equal to an integer. This can always be guarantied by a suitable choice of the still disposable parameter $`\lambda `$.
Since the exponents $`\mu (\kappa )`$, $`\nu _1`$, $`\nu _2`$ are restricted only by (2.5) but otherwise arbitrary, there are other solutions of the partial differential equation (2.2 ) relevant as weight functions in our contour integrals (2.1). We may assume that $`\nu _1`$ and $`\nu _2`$ are positive integers, preferentially
$$\nu _1=\nu _2=1,$$
$`(2.11)`$
but for the time being we want to keep $`\nu _1`$ and $`\nu _2`$ in the formulas. If $`\nu _1`$ and $`\nu _2`$ are increased by any positive integers $`q_1`$ and $`q_2`$, respectively, and $`\mu (\kappa )`$ simultaneously is decreased by $`(2/3)q_1+(1/3)q_2`$, then (2.5) is still satisfied. We therefore have to consider, for $`q_1,q_2=0,1,2,\mathrm{}`$, the set of solutions
$$\begin{array}{cc}& V(\kappa ;q_1,q_2)=V(\kappa ;q_1,q_2;s,t_1,t_2)\hfill \\ & =\underset{m=0}{\overset{\mathrm{}}{}}\underset{n_1=0}{\overset{m}{}}\underset{n_2=0}{\overset{m}{}}\mathrm{\Gamma }(\nu _1+q_1+n_1)\mathrm{\Gamma }(\nu _2+q_2+n_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2m)\hfill \\ & \times b(\kappa ;m,n_1,n_2)(s\kappa s_0)^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2+m}(t_1\kappa t_{10})^{\nu _1q_1n_1}(t_2\kappa t_{20})^{\nu _2q_2n_2}.\hfill \end{array}$$
$`(2.12)`$
As indicated, we will use a short-hand notation suppressing the dependence on the variables $`s,t_1,t_2`$. We now have to choose appropriate contours for the integral representation (2.1). For each of the integrals over $`t_1`$ or $`t_2`$ a closed circle, traversed once in the positive sense, around the relevant singular point of the integrand is appropriate, since the pertinent exponent is an integer. For the $`s`$ -integral we need an infinite contour which starts somewhere at infinity where the exponential factor of the integral vanishes, surrounds one of the singular points in the positive sense, and returns (on a different sheet) to the starting point. Assuming that $`s_0`$ is real and positive, then, if
$$0<\mathrm{arg}(z)<\frac{1}{3}\pi $$
$`(2.13)`$
the starting- and end-point at infinity has the phase $`\pi /2`$. If we agree that
$$\mathrm{arg}(s\kappa s_0)=0$$
$`(2.14)`$
when $`s`$ is positive and sufficiently large, the integral of a single term of the infinite series (2.4) can be evaluated:
$$\begin{array}{cc}& z^\lambda (2\pi \text{i})^3\stackrel{(\kappa t_{20}+)}{}\stackrel{(\kappa t_{10}+)}{}\underset{\kappa s_0+\mathrm{i}\mathrm{}}{\overset{(\kappa s_0+)}{}}\mathrm{exp}(z^3s+z^2t_1+zt_2)\hfill \\ & \times (s\kappa s_0)^{\mu (\kappa )+m}(t_1\kappa t_{10})^{\nu _1n_1}(t_2\kappa t_{20})^{\nu _2n_2}\text{d}s\text{d}t_1\text{d}t_2\hfill \\ & =\mathrm{exp}(\kappa s_0z^3+\kappa t_{10}z^2+\kappa t_{20}z)\frac{z^{\tau (\kappa )3m+2n_1+n_2}}{\mathrm{\Gamma }(\mu (\kappa )m)\mathrm{\Gamma }(\nu _1+n_1)\mathrm{\Gamma }(\nu _2+n_2)},\hfill \end{array}$$
$`(2.15)`$
where, in the power of $`z`$ on the right-hand side, already use has been made of (2.5). As a consequence, the integral of $`V(\kappa ;q_1,q_2)`$ yields, if the series is integrated term by term, one or the other of the formal solutions (1.2)-(1.3):
$$\begin{array}{cc}\hfill f_\mathrm{}j(z):=& z^\lambda (2\pi \text{i})^3\stackrel{(\kappa t_{20}+)}{}\stackrel{(\kappa t_{10}+)}{}\underset{\kappa s_0+\mathrm{i}\mathrm{}}{\overset{(\kappa s_0+)}{}}\mathrm{exp}(z^3s+z^2t_1+zt_2)\hfill \\ & \times V(\kappa ;q_1,q_2;s,t_1,t_2)\text{d}s\text{d}t_1\text{d}t_2f_\mathrm{}j^{\mathrm{asy}}(z)\hfill \end{array}$$
$`(2.16)`$
in the sector $`0<\mathrm{arg}(z)<\frac{1}{3}\pi `$, where $`j=1`$ if $`\kappa =1`$ or $`j=2`$ if $`\kappa =1`$ . Each of the solutions defined by the integral representation (2.16 ) has one of the formal solutions (1.2) or (1.3) as its asymptotic expansion as $`z\mathrm{}`$ in the indicated sector. It follows by rotation of the contour that the asymptotic expansions are theoretically valid in the larger sectors $`\frac{1}{6}\pi <\mathrm{arg}(z)<\frac{5}{6}\pi `$ for $`j=1`$ or $`\frac{1}{2}\pi <\mathrm{arg}(z)<\frac{1}{2}\pi `$ for $`j=2`$, respectively.
Looking at (2.15) we may see that all the terms for which $`3m2n_1n_2`$ is the same yield the same power of $`z`$. For the coefficients of the asymptotic expansions (1.2) or (1.3) we then have the representation
$$a_n(\kappa )=\underset{(m,n_1,n_2)I_n}{}b(\kappa ;m,n_1,n_2),$$
$`(2.17)`$
where
$$I_n=\{(m,n_1,n_2):3m2n_1n_2=n\}\text{N}_0\times \text{N}_0\times \text{N}_0,n=0,1,\mathrm{}.$$
$`(2.18)`$
Because of the properties of the $`b(m,n_1,n_2)`$, the sum in (2.18) is finite for each finite $`n`$, in particular we have
$$a_0(\kappa )=b_0(\kappa ;0,0,0)=1.$$
$`(2.19)`$
3. Analytic continuation of the integrand
Below we have to consider the integral representation (2.1) with $`t_1`$ and $`t_2`$-contours which are simple closed curves surrounding in the positive sense both the finite singular points $`t_{10}`$ and $`t_{10}`$ or $`t_{20}`$ and $`t_{20}`$, respectively, and with an $`s`$-contour which starts at or near $`s_0+\text{i}\mathrm{}`$, surrounds both the finite singular points $`s_0`$ and $`s_0`$ once in the positive sense and ends at $`s_0+\text{i}\mathrm{}`$. We therefore need the analytic continuation of the integrand between the two singular points along this contour. With appropriate power factors $`\mathrm{\Phi }`$ included for later convenience, the continuation formula reads
$$\begin{array}{cc}& \mathrm{\Phi }(\kappa ;r_1,r_2)V(\kappa ;r_1,r_2)\hfill \\ & =\underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Phi }(\kappa ;q_1,q_2)V(\kappa ;q_1,q_2)+U(\kappa ;r_1,r_2),\hfill \end{array}$$
$`(3.1)`$
where
$$\mathrm{\Phi }(\kappa ;r_1,r_2)=(2\kappa s_0)^{\mu (\kappa ){\scriptscriptstyle \frac{2}{3}}r_1{\scriptscriptstyle \frac{1}{3}}r_2}.$$
$`(3.2)`$
The effect of these power factors is that in (3.1) the total powers with non-integer exponents are powers of $`\frac{1}{2}[1s/(\kappa s_0)]`$ or of $`\frac{1}{2}[1+s/(\kappa s_0)]`$, respectively. We may agree that $`\mathrm{arg}(1s/s_0)=\mathrm{arg}(1+s/s_0)=0`$ when $`s`$ is on the real axis between $`s_0`$ and $`s_0`$. Also, as above, $`\mathrm{arg}(s)=0`$ when s is larger than $`s_0`$. Then, by analytic continuation along the contour under consideration, we have
$$\mathrm{\Phi }(1;r_1,r_2)=(2s_0)^{\mu (1){\scriptscriptstyle \frac{2}{3}}r_1{\scriptscriptstyle \frac{1}{3}}r_2}\mathrm{exp}(\text{i}\pi (\mu (1)+\frac{2}{3}r_1+\frac{1}{3}r_2),$$
$`(3.3)`$
$$\mathrm{\Phi }(1;r_1,r_2)=(2s_0)^{\mu (1){\scriptscriptstyle \frac{2}{3}}r_1{\scriptscriptstyle \frac{1}{3}}r_2}.$$
$`(3.4)`$
Let us rewrite (3.1) using an even more condensed notation, writing $`q`$ for $`(q_1,q_2)`$ in the parameter list of the various functions and writing a sum over $`q`$ in place of a double sum over $`q_1`$ and $`q_2`$, etc. The above continuation formula then reads
$$\mathrm{\Phi }(\kappa ;r)V(\kappa ;r)=\underset{q=r}{}E(\kappa ;r,q)\mathrm{\Phi }(\kappa ;q)V(\kappa ;q)+U(\kappa ;r),$$
$`(3.5)`$
and the second continuation formula
$$\begin{array}{cc}& U(\kappa ;r)=\mathrm{\Phi }(\kappa ,r)V(\kappa ,r)\hfill \\ & \underset{q=r}{}\underset{p=q}{}E(\kappa ;r,q)E(\kappa ;q,p)\mathrm{\Phi }(\kappa ;p)V(\kappa ;p)\underset{q=r}{}E(\kappa ;r,q)U(\kappa ;q)\hfill \end{array}$$
$`(3.6)`$
follows by the requirement of consistency of (3.5) for $`\kappa `$ replaced by $`\kappa `$.
4. Asymptotic expansions for the coefficients in the continuation formula
We now want to determine the $`E`$-coefficients in the continuation formulas (3.5) , (3.6) by means of the asymptotic method of Darboux applied to the variable $`s`$ . The left-hand side of the continuation formula (3.5) is
$$\begin{array}{cc}& (\frac{1}{2}\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}r_1+{\scriptscriptstyle \frac{1}{3}}r_2}\underset{m=0}{\overset{\mathrm{}}{}}\underset{n_1=0}{\overset{m}{}}\underset{n_2=0}{\overset{m}{}}\mathrm{\Gamma }(\nu _1+r_1+n_1)\mathrm{\Gamma }(\nu _2+r_2+n_2)\hfill \\ & \times \mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2m)b(\kappa ;m,n_1,n_2)(2\kappa s_0)^m\hfill \\ & \times (\frac{1}{2}\frac{s}{2\kappa s_0})^m(t_1\kappa t_{10})^{\nu _1r_1n_1}(t_2\kappa t_{20})^{\nu _2r_2n_2}.\hfill \end{array}$$
$`(4.1)`$
The leading singular term, when $`(s,t_1,t_2)(\kappa s_0,\kappa t_{10},\kappa t_{20})`$, on the right-hand side is
$$\begin{array}{cc}& \underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \mathrm{\Gamma }(\nu _1+q_1)\mathrm{\Gamma }(\nu _2+q_2)b(\kappa ;0,0,0)\hfill \\ & \times (\frac{1}{2}+\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2}(t_1+\kappa t_{10})^{\nu _1q_1}(t_2+\kappa t_{20})^{\nu _2q_2}.\hfill \end{array}$$
$`(4.2)`$
By means of the binomial theorem in its hypergeometric-series-form
$$(1x)^\alpha =\underset{j=0}{\overset{\mathrm{}}{}}\frac{(\alpha )_j}{j!}x^j,$$
$`(4.3)`$
where
$$(\alpha )_j=\alpha (\alpha +1)\mathrm{}(\alpha +j1)=\mathrm{\Gamma }(\alpha +j)/\mathrm{\Gamma }(\alpha )$$
means the Pochhammer symbol, we may expand, if $`|t|`$ is sufficiently large,
$$(t+\kappa t_0)^{\nu q}=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(\nu +q)_j}{j!}(2\kappa t_0)^j(t\kappa t_0)^{\nu qj}$$
$`(4.4)`$
and, if $`|s/(\kappa s_0)|`$ is sufficiently small,
$$(\frac{1}{2}+\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2}=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)_m}{m!}(\frac{1}{2}\frac{s}{2\kappa s_0})^m.$$
$`(4.5)`$
Then the leading singular term on the right becomes
$$\begin{array}{cc}& \underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \mathrm{\Gamma }(\nu _1+q_1)\underset{j_1=0}{\overset{\mathrm{}}{}}\frac{(\nu _1+q_1)_{j_1}}{j_1!}(2\kappa t_{10})^{j_1}(t_1\kappa t_{10})^{\nu _1q_1j_1}\hfill \\ & \times \mathrm{\Gamma }(\nu _2+q_2)\underset{j_2=0}{\overset{\mathrm{}}{}}\frac{(\nu _2+q_2)_{j_2}}{j_2!}(2\kappa t_{20})^{j_2}(t_2\kappa t_{20})^{\nu _2q_2j_2}\hfill \\ & \times b(\kappa ;0,0,0)\underset{m=0}{\overset{\mathrm{}}{}}\frac{(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)_m}{m!}(\frac{1}{2}\frac{s}{2\kappa s_0})^m.\hfill \end{array}$$
$`(4.6)`$
The coefficients of this series should agree asymptotically, as $`m\mathrm{}`$, with those of the series on the left-hand side, where the power factor in front of the series (4.1), when $`s\kappa s_0`$, tends to $`1`$ and may be omitted. We therefore obtain
$$\begin{array}{cc}& \underset{n_1=0}{\overset{m}{}}\underset{n_2=0}{\overset{m}{}}\mathrm{\Gamma }(\nu _1+r_1+n_1)\mathrm{\Gamma }(\nu _2+r_2+n_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2m)\hfill \\ & \times b(\kappa ;m,n_1,n_2)(2\kappa s_0)^m(t_1\kappa t_{10})^{\nu _1r_1n_1}(t_2\kappa t_{20})^{\nu _2r_2n_2}\hfill \\ & \underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \mathrm{\Gamma }(\nu _1+q_1)\underset{j_1=0}{\overset{\mathrm{}}{}}\frac{(\nu _1+q_1)_{j_1}}{j_1!}(2\kappa t_{10})^{j_1}(t_1\kappa t_{10})^{\nu _1q_1j_1}\hfill \\ & \times \mathrm{\Gamma }(\nu _2+q_2)\underset{j_2=0}{\overset{\mathrm{}}{}}\frac{(\nu _2+q_2)_{j_2}}{j_2!}(2\kappa t_{20})^{j_2}(t_2\kappa t_{20})^{\nu _2q_2j_2}\hfill \\ & \times b(\kappa ;0,0,0)\frac{(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)_m}{m!},\hfill \end{array}$$
$`(4.7)`$
which holds asymptotically as $`m\mathrm{}`$. On both sides of this asymptotic equation a double power series in the same variables appears, so the coefficients of the corresponding terms must be equal. This yields, for each set $`r_1,r_2,n_1,n_2`$,
$$\begin{array}{cc}& \mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2m)(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{q_1=r_1}{\overset{r_1+n_1}{}}\underset{q_2=r_2}{\overset{r_2+n_2}{}}E(\kappa ;r_1,r_2;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{m!}\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=r_1+n_1q_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=r_2+n_2q_2}b(\kappa ;0,0,0)\hfill \end{array}$$
$`(4.8)`$
or, if we introduce new indices of summation and make use of the reflection formula of the gamma function,
$$\begin{array}{cc}& \pi [\mathrm{\Gamma }(1+m)]^1(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}E(\kappa ;r_1,r_2;r_1+p_1,r_2+p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2\frac{2}{3}p_1\frac{1}{3}p_2+m)}{\mathrm{\Gamma }(1+m)}\hfill \\ & \times \frac{\mathrm{\Gamma }(1+\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2+m)}{\mathrm{\Gamma }(1+m)}\mathrm{sin}(\pi [\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2])\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=n_1p_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=n_2p_2}b(\kappa ;0,0,0).\hfill \end{array}$$
$`(4.9)`$
Now the left-hand side is independent of $`r_1`$ and $`r_2`$ and so is the product of the two ratios of gamma functions, when $`m\mathrm{}`$, on the right. Therefore $`E(\kappa ;r_1,r_2;r_1+p_1,r_2+p_2)\mathrm{sin}(\pi [\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2])`$ must be independent of $`r_1`$ and $`r_2`$ too. This proves
Lemma 3.
$$E(\kappa ;r_1,r_2;r_1+p_1,r_2+p_2)=\frac{\mathrm{sin}(\pi \mu (\kappa ))}{\mathrm{sin}(\pi [\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2])}E(\kappa ;0,0;p_1,p_2).$$
$`(4.10)`$
After Lemma 3 is available, we need to determine the $`E`$-coefficients for $`r_1=r_2=0`$ only, and it is advantageous to introduce the closely related coefficients
$$e(\kappa ;p_1,p_2)=\mathrm{sin}(\pi \mu (\kappa ))E(\kappa ;0,0;p_1,p_2),$$
$`(4.11)`$
which are independent of $`\lambda `$. Eq (4.9), with $`r_1=r_2=0`$ and the factor $`b(\kappa ;0,0,0)`$, which is equal to $`1`$, omitted, then becomes
$$\begin{array}{cc}& \frac{\pi }{\mathrm{\Gamma }(1+\mu (\kappa )+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+m)}{\mathrm{\Gamma }(1+m)}\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=n_1p_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=n_2p_2}.\hfill \end{array}$$
$`(4.12)`$
We then may solve (4.12) for the $`e`$-coefficient with $`p_1=n_1,p_2=n_2`$ and obtain the following asymptotic formula in terms of a $`b`$-coefficient and the earlier $`e`$-coefficients,
$$\begin{array}{cc}& e(\kappa ;n_1,n_2)\pi (2\kappa s_0)^m\hfill \\ & \times \frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}b(\kappa ;m,n_1,n_2)\hfill \\ & \underset{(p_1,p_2)(n_1,n_2)}{\underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}}e(\kappa ;p_1,p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+m)}{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=n_1p_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=n_2p_2}.\hfill \end{array}$$
$`(4.13)`$
By repeated application of this formula, all the $`e`$-coefficients on the right-hand side may be eliminated, and this yields the remarkable explicit limit formula
$$\begin{array}{cc}& e(\kappa ;n_1,n_2)=\pi \hfill \\ & \times \underset{m\mathrm{}}{lim}\frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times (2\kappa s_0)^m\underset{j_1=0}{\overset{n_1}{}}\underset{j_2=0}{\overset{n_2}{}}b(\kappa ;m,j_1,j_2)\frac{(2\kappa t_{10})^{n_1j_1}}{(n_1j_1)!}\frac{(2\kappa t_{20})^{n_2j_2}}{(n_2j_2)!}.\hfill \end{array}$$
$`(4.14)`$
The proof of this formula will be given below in Section 7.
Although approximate asymptotic equations or limit formulas such as (4.13) or (4.14) are interesting from a theoretical point of view, we finally need more, namely a detailed asymptotic expansion suitable for accurate numerical evaluation. This can be obtained, on the basis of Schäfke and Schmidt , essentially in the same way as above, apart from the following two refinements: The power factor in front of (4.1) can no longer be omitted, and we have to include a finite number of singular terms rather than the leading one, (4.2), alone. The result of this procedure, which will be derived in more detail below, is
Theorem 1. The $`E`$-coefficients, or $`e`$-coefficients according to (4.11), in the continuation formula (3.1) or (3.5) are
$$\begin{array}{cc}& e(\kappa ;n_1,n_2)=\pi \hfill \\ & \times [\frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{(q_1,q_2)(n_1,n_2)}{\underset{q_1=0}{\overset{n_1}{}}\underset{q_2=0}{\overset{n_2}{}}}e(\kappa ;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times [1+\underset{k=1}{\overset{K}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_k}{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2m)_k}H(\kappa ;k,l_1,l_2;q_1,q_2)+O(m^{K1})]\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1+q_1+l_1=n_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2+q_2+l_2=n_2}]\hfill \\ & \times \left[1+\underset{k=1}{\overset{K}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2)_k}{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2m)_k}H(\kappa ;k,0,0;n_1,n_2)+O(m^{K1})\right]^1,\hfill \end{array}$$
$`(4.15)`$
where
$$H(\kappa ;k,l_1,l_2;q_1,q_2)=\underset{j=l_1+l_2}{\overset{k}{}}\frac{(\mu (\kappa ))_{kj}}{(kj)!(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_j}(2\kappa s_0)^jb(\kappa ;j,l_1,l_2).$$
$`(4.16)`$
With a suitable choice of $`m`$ and $`K`$, Theorem 1 may be used to compute accurate values of the $`e`$-coefficients, beginning with
$$\begin{array}{cc}& e(\kappa ;0,0)=\pi \frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )+m)}(2\kappa s_0)^mb(\kappa ;m,0,0)\hfill \\ & \times \left[1+\underset{k=1}{\overset{K}{}}\frac{(1+\mu (\kappa ))_k}{(1+\mu (\kappa )m)_k}H(\kappa ;k,0,0;0,0)+O(m^{K1})\right]^1.\hfill \end{array}$$
$`(4.17)`$
While the $`e`$-coefficients are independent of $`\lambda `$, the approximate values computed for any finite $`m`$ do depend on it. Thus $`\lambda `$ here plays the role of a computational parameter which could be adjusted for optimal accuracy. The best choice from this point of view , in particular when $`L`$ is not small, is $`\lambda =L`$ or $`\lambda =L`$, according to the discussion in Section 7 below.
5. Multiplicative solutions and the Floquet exponent
We now want to construct a linear combination of the solutions at infinity,
$$f_p(z)=\alpha f_\mathrm{}1(z)+\beta f_\mathrm{}2(z),$$
$`(5.1)`$
which is a multiplicative solution such that, after analytic continuation along a sufficiently large circle around the origin traversed once in the negative sense, this solution remains the same apart from multiplication by a constant factor, that is
$$f_p(\mathrm{e}^{2\pi \mathrm{i}}z)=pf_p(z).$$
$`(5.2)`$
This solution is proportional to one of the Floquet solutions introduced above in the introduction.
We use the integral representation with contours which surround both the finite singular points as introduced above in Section 3. This integral is equal to the sum of the two integrals with contours, also considered above in Section 2, which surround only one of these singular points. If the contour is kept fixed, then the circle in the $`z`$-plane traversed once in the negative sense corresponds to a circle, around the origin and with radius greater than $`s_0`$, in the $`s`$-plane traversed three times in the positive sense, since then $`z^3s`$ in the exponential part of the integrand does not change.
Let us consider the integral representation with $`V(1;0)`$ in the integrand and with the above phase conventions. From (3.5) we have
$$\mathrm{\Phi }(1;0)V(1;0)=\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)V(1;q)+U(1;0).$$
$`(5.3)`$
The integral then yields a function
$$f^\mathrm{I}(z):=\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)f_\mathrm{}2(z),$$
$`(5.4)`$
where the two terms come from the two singular points, which contribute via the left- or right-hand side of (5.3), respectively. Let us now consider, in the $`s`$-plane, a simple closed loop, homotopic to the circle mentioned above, consisting of small circles around the singular points $`s_0`$ and $`s_0`$ and straight line segments along the real axis between these singular points. Let us start with the continuation formula (5.3) in a neighbourhood of the origin, where the left- as well as the right-hand side of (5.3) are valid, and see what happens when we follow the loop in the positive sense. We shall always refer to the above phase conventions and display any additional phases explicitly. Traversing the circular part around $`s_0`$ multiplies the function $`V(1;0)`$ by the phase factor $`\mathrm{exp}(2\pi \text{i}\mu (1))`$, so that we then have
$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}\mu (1))\mathrm{\Phi }(1;0)V(1;0)\hfill \\ & =\mathrm{exp}(2\pi \text{i}\mu (1))\{\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)V(1;q)+U(1;0)\},\hfill \end{array}$$
$`(5.5)`$
where the right-hand side is the analytic continuation of the left by means of (5.3). Next we have to traverse the circular part around $`s_0`$, which multiplies the function $`V(1;q)`$ by the phase factor $`\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])`$, where
$$\phi (q)=\frac{2}{3}q_1+\frac{1}{3}q_2,$$
$`(5.6)`$
so that we have on the straight line segment from $`s_0`$ to $`s_0`$ , after the loop has been traversed once,
$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}[\mu (1)+\mu (1)])\underset{q=0}{}E(1;0,q)\mathrm{exp}(2\pi \text{i}\phi (q))\mathrm{\Phi }(1;q)V(1;q)\hfill \\ & +\mathrm{exp}(2\pi \text{i}\mu (1))U(1;0)\hfill \end{array}$$
$$\begin{array}{cc}& =\mathrm{exp}(2\pi \text{i}\mu (1))\{\mathrm{\Phi }(1;0)V(1;0)+\underset{q=0}{}\{\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])1\}U(1;q)\hfill \\ & +\underset{q=0}{}\underset{p=q}{}E(1;0,q)\{\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])1\}E(1;q,p)\mathrm{\Phi }(1;p)V(1;p)\},\hfill \end{array}$$
$`(5.7)`$
where again the right-hand side is the analytic continuation of the left. It is evident that following the loop further leads to increasingly lengthier and more complicated formulas, not suitable for being fully displayed in this paper. We therefore want to stop here for a moment and consider the integral representation with the integrand obtained after the loop in the $`s`$-plane has been traversed only once (rather than three times, as finally needed). A representative example of the terms in (5.7) then is
$$\underset{q=0}{}\underset{p=q}{}E(1;0,q)\{\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])1\}E(1;q,p)\mathrm{\Phi }(1;p)V(1;p).$$
$`(5.8)`$
The integral of $`V(1;p)`$ yields $`f_\mathrm{}1(z)`$, independent of $`p`$, so that the sum over $`p`$ can now be performed and , because of Lemma 3, the multiple sum reduces to a product of single sums. After integration we therefore obtain for (5.8)
$$\begin{array}{cc}& \{2\text{i}\mathrm{sin}(\pi \mu (1))\mathrm{exp}(\text{i}\pi \mu (1))\underset{q=0}{}E(1;0,q)(2s_0)^{\phi (q)}\mathrm{exp}(2\pi \text{i}\phi (q))\hfill \\ & \times \underset{\stackrel{~}{p}=0}{}E(1;0,\stackrel{~}{p})(2s_0)^{\phi (\stackrel{~}{p})}\mathrm{exp}(\text{i}\pi \phi (\stackrel{~}{p}))\}\mathrm{\Phi }(1;0)f_\mathrm{}1(z).\hfill \end{array}$$
$`(5.9)`$
It is now convenient to introduce the phase factor
$$\eta :=\mathrm{exp}(\frac{1}{3}\pi \text{i})=\frac{1}{2}(1+\text{i}\sqrt{3})$$
$`(5.10)`$
which satisfies
$$\eta ^6=1,1+\eta ^2+\eta ^4=0.$$
$`(5.11)`$
Also, for the sums occurring here and below we may introduce the ”Stokes multipliers”
$$\begin{array}{cc}& \sigma _n(1):=\mathrm{sin}(\pi \mu (1))\underset{q=0}{}E(1;0,q)(2s_0)^{\phi (q)}\mathrm{exp}(2n\pi \text{i}\phi (q)),\hfill \\ & \sigma _n(1):=\mathrm{sin}(\pi \mu (1))\underset{q=0}{}E(1;0,q)(2s_0)^{\phi (q)}\mathrm{exp}((2n+1)\pi \text{i}\phi (q)),\hfill \end{array}$$
$`(5.12)`$
which become
$$\begin{array}{cc}& \sigma _0(1)=S_0(1)+(2s_0)^{1/3}S_1(1)+(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _0(1)=S_0(1)+\eta (2s_0)^{1/3}S_1(1)+\eta ^2(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _1(1)=S_0(1)+\eta ^2(2s_0)^{1/3}S_1(1)+\eta ^4(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _1(1)=S_0(1)+\eta ^3(2s_0)^{1/3}S_1(1)+(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _2(1)=S_0(1)+\eta ^4(2s_0)^{1/3}S_1(1)+\eta ^2(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _2(1)=S_0(1)+\eta ^5(2s_0)^{1/3}S_1(1)+\eta ^4(2s_0)^{2/3}S_2(1)\hfill \end{array}$$
$`(5.13)`$
in terms of the real partial sums
$$\begin{array}{cc}\hfill S_0(\kappa ):=& \underset{l=0}{\overset{\mathrm{}}{}}(2\kappa s_0)^l\underset{(n_1,n_2)J_{3l}}{}e(\kappa ;n_1,n_2)=e(\kappa ;0,0)+\mathrm{},\hfill \\ \hfill S_1(\kappa ):=& \underset{l=0}{\overset{\mathrm{}}{}}(2\kappa s_0)^l\underset{(n_1,n_2)J_{3l+1}}{}e(\kappa ;n_1,n_2)=e(\kappa ;0,1)+\mathrm{},\hfill \\ \hfill S_2(\kappa ):=& \underset{l=0}{\overset{\mathrm{}}{}}(2\kappa s_0)^l\underset{(n_1,n_2)J_{3l+2}}{}e(\kappa ;n_1,n_2)=e(\kappa ;1,0)+e(\kappa ;0,2)+\mathrm{},\hfill \end{array}$$
$`(5.14)`$
where
$$J_l=\{(n_1,n_2):2n_1+n_2=l\}\text{N}_0\times \text{N}_0,l=0,1,\mathrm{}.$$
$`(5.15)`$
The integral of the representative term (5.9) then becomes
$$\frac{2\text{i}\mathrm{exp}(\text{i}\pi \mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _0(1)\mathrm{\Phi }(1;0)f_\mathrm{}1(z).$$
$`(5.16)`$
In total, we have
$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}\lambda /3)f^\mathrm{I}(\mathrm{e}^{2\pi \mathrm{i}/3}z)=\frac{\mathrm{exp}(2\pi \text{i}[\mu (1)+\mu (1)])}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\mathrm{\Phi }(1;0)f_\mathrm{}2(z)\hfill \\ & +\mathrm{exp}(2\pi \text{i}\mu (1)\{1+\frac{2\text{i}\mathrm{exp}(\pi \text{i}\mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _0(1)\}\mathrm{\Phi }(1;0)f_\mathrm{}1(z),\hfill \end{array}$$
$`(5.17)`$
where
$$f^\mathrm{I}(z)=\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+\frac{1}{\mathrm{sin}(\pi \mu (1))}\sigma _0(1)\mathrm{\Phi }(1;0)f_\mathrm{}2(z).$$
$`(5.18)`$
So far we have traversed the $`s`$-loop once and obtained on it the analytic continuation of the integrand, but we have to traverse it three times. This yields
$$f^\mathrm{I}(\mathrm{e}^{2\pi \mathrm{i}}z)=S_{11}\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+S_{12}\mathrm{\Phi }(1;0)f_\mathrm{}2(z)$$
$`(5.19)`$
with lengthy expressions for $`S_{11}`$ and $`S_{12}`$.
In a similar way, the integral representation with $`V(1;0)`$ yields
$$f^{\mathrm{II}}(z):=\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)f_\mathrm{}1(z)+\mathrm{\Phi }(1;0)f_\mathrm{}2(z)$$
$`(5.20)`$
or
$$f^{\mathrm{II}}(z)=\frac{1}{\mathrm{sin}(\pi \mu (1))}\sigma _0(1)\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+\mathrm{\Phi }(1;0)f_\mathrm{}2(z)$$
$`(5.21)`$
and leads, after the loop has been traversed once, to
$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}\lambda /3)f^{\mathrm{II}}(\mathrm{e}^{2\pi \mathrm{i}/3}z)=\{2\text{i}\mathrm{exp}(\pi \text{i}\mu (1))\sigma _0(1)+\frac{\mathrm{exp}(2\pi \text{i}\mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\hfill \\ & \frac{4\mathrm{exp}(\pi \text{i}[\mu (1)+\mu (1)])}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _1(1)\sigma _0(1)\}\mathrm{\Phi }(1;0)f_\mathrm{}1(z)\hfill \\ & +\mathrm{exp}(2\pi \text{i}\mu (1))\{1+\frac{2\text{i}\mathrm{exp}(\pi \text{i}\mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _1(1)\}\mathrm{\Phi }(1;0)f_\mathrm{}2(z)\hfill \end{array}$$
$`(5.22)`$
and, after three times,
$$f^{\mathrm{II}}(\mathrm{e}^{2\pi \mathrm{i}}z)=S_{21}\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+S_{22}\mathrm{\Phi }(1;0)f_\mathrm{}2(z),$$
$`(5.23)`$
where again the expressions for $`S_{21}`$ and $`S_{22}`$ are too lengthy to be displayed here.
What we really want to obtain are the circuit relations for $`f_\mathrm{}1`$ and $`f_\mathrm{}2`$, that is
$$\begin{array}{cc}& f_\mathrm{}1(\mathrm{e}^{2\pi \mathrm{i}}z)=T_{11}f_\mathrm{}1(z)+T_{12}f_\mathrm{}2(z),\hfill \\ & f_\mathrm{}2(\mathrm{e}^{2\pi \mathrm{i}}z)=T_{21}f_\mathrm{}1(z)+T_{22}f_\mathrm{}2(z),\hfill \end{array}$$
$`(5.24)`$
where
$$\begin{array}{cc}& T_{11}=[S_{11}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{21}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}],\hfill \\ & T_{12}=[\mathrm{\Phi }(1;0)/\mathrm{\Phi }(1;0)][S_{12}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{22}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}],\hfill \\ & T_{21}=[\mathrm{\Phi }(1;0)/\mathrm{\Phi }(1;0)][S_{21}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{11}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}],\hfill \\ & T_{22}=[S_{22}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{12}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}].\hfill \end{array}$$
$`(5.25)`$
It turns out that each numerator contains a common factor which compensates the denominator. The result then is
Theorem 2. The coefficients in the circuit relations (5.24) are
$$\begin{array}{cc}& T_{11}=\mathrm{exp}(2\pi \text{i}\tau (1))\hfill \\ & +4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _0(1)\sigma _0(1)+4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _0(1)+4\sigma _2(1)\sigma _0(1)\hfill \\ & +4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)+4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _0(1)+4\sigma _1(1)\sigma _0(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _0(1)+16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)+16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +64\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1),\hfill \end{array}$$
$`(5.26a)`$
$$\begin{array}{cc}& T_{12}=(2s_0)^{(1/3)[\tau (1)\tau (1)]}\{2\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\hfill \\ & +2\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)+2\text{i}\mathrm{exp}(2\pi \text{i}\tau (1))\sigma _0(1)\hfill \\ & +8\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)+8\text{i}\sigma _2(1)\sigma _1(1)\sigma _0(1)\hfill \\ & +8\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)+8\text{i}\sigma _2(1)\sigma _1(1)\sigma _1(1)\hfill \\ & +32\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\},\hfill \end{array}$$
$`(5.26b)`$
$$\begin{array}{cc}& T_{21}=(2s_0)^{(1/3)[\tau (1)\tau (1)]}\{2\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _0(1)\hfill \\ & 2\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)2\text{i}\sigma _1(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _0(1)\hfill \\ & 32\text{i}\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\},\hfill \end{array}$$
$`(5.26c)`$
$$\begin{array}{cc}& T_{22}=\mathrm{exp}(2\pi \text{i}\tau (1))\hfill \\ & +4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)+4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)+4\sigma _2(1)\sigma _1(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1).\hfill \end{array}$$
$`(5.26d)`$
It turns out that
$$T_{11}T_{22}T_{12}T_{21}=1.$$
$`(5.27)`$
We are looking for a multiplicative solution such that (5.2) holds. It then follows from (5.1), 5.2), (5.24) that
$$\begin{array}{cc}& (T_{11}p)\alpha +T_{21}\beta =0,\hfill \\ & T_{12}\alpha +(T_{22}p)\beta =0\hfill \end{array}$$
$`(5.28)`$
and, as a consequence, that
$$(T_{11}p)(T_{22}p)=T_{12}T_{21}$$
$`(5.29)`$
or
$$p^2(T_{11}+T_{22})p+1=0,$$
$`(5.30)`$
where (5.27) has been used. The roots $`p_1`$, $`p_2`$ of this equation satisfy
$$p_1+p_2=T_{11}+T_{22},$$
$`(5.31)`$
and they may be represented in terms of one (not necessarily real) parameter $`\omega `$ as
$$p_1=\mathrm{exp}(2\pi \text{i}\omega ),p_2=\mathrm{exp}(2\pi \text{i}\omega ).$$
$`(5.32)`$
Then we have
$$p_1+p_2=2\mathrm{cos}(2\pi \omega )=T_{11}+T_{22},$$
$`(5.33)`$
and the final result is
Theorem 3. The characteristic exponent $`\omega `$ of the multiplicative solutions is given by
$$\mathrm{cos}(2\pi \omega )=\mathrm{cos}(2\pi \tau (1))+X,$$
$`(5.34)`$
where
$$\begin{array}{cc}& X=2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _0(1)\sigma _0(1)+2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _0(1)+2\sigma _2(1)\sigma _0(1)\hfill \\ & +2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)+2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _0(1)+2\sigma _1(1)\sigma _0(1)\hfill \\ & +2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)+2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)+2\sigma _2(1)\sigma _1(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _0(1)+8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)+8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +32\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1).\hfill \end{array}$$
$`(5.35)`$
For each of the roots $`p_1`$ and $`p_2`$ we may determine the ratio of $`\alpha `$ and $`\beta `$ from the upper or lower equation of (5.28). Choosing in each case a convenient normalization, we get the desired multiplicative solutions
$$\begin{array}{cc}& f_{p1}(z)=(T_{22}p_1)f_\mathrm{}1(z)T_{12}f_\mathrm{}2(z),\hfill \\ & f_{p2}(z)=(T_{22}p_2)f_\mathrm{}1(z)T_{12}f_\mathrm{}2(z).\hfill \end{array}$$
$`(5.36)`$
Here the lower equation of (5.28) has been used in both cases so that $`T_{11}`$, which consists of a considerably longer expression than $`T_{22}`$, does not appear.
Whenever $`D_m=0`$ for all the $`m=1,2,\mathrm{},6`$, the origin is a regular singular point of the differential equation with exponents $`\omega =L,L`$. We are not able, however, to see analytically that then the lengthy expression (5.35) reduces to $`X=\mathrm{cos}(2\pi L)\mathrm{cos}(2\pi \tau (1))`$, but this is confirmed in examples of numerical computations, as expected.
6. Asymptotic behaviour of the late coefficients of the formal power series solutions
According to (4.12), the leading terms of the asymptotic behaviour of the $`b`$-coefficients for large $`m`$ are given by
$$\begin{array}{cc}& b(\kappa ;m,n_1,n_2)\frac{1}{\pi }(2\kappa s_0)^m\frac{\mathrm{\Gamma }(1+\mu (\kappa )+m)}{\mathrm{\Gamma }(1+m)}\hfill \\ & \times \underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+m)\hfill \\ & \times \frac{1}{(n_1p_1)!}(2\kappa t_{10})^{n_1p_1}\frac{1}{(n_2p_2)!}(2\kappa t_{20})^{n_2p_2}.\hfill \end{array}$$
$`(6.1)`$
This result may be used to discuss the asymptotic behaviour of the late $`a`$-coefficients of the formal solutions. Writing the decomposition (2.17) of the $`a`$-coefficients in terms of the $`b`$-coefficients separately for each of three consecutive indices $`n=3N,3N+1,3N+2`$, we have
$$a_{3N}(\kappa )=\underset{l=0}{\overset{2N}{}}\underset{(n_1,n_2)J_{3l}}{}b(\kappa ;N+l,n_1,n_2)=b(\kappa ;N,0,0)+\mathrm{},$$
$`(6.2)`$
$$a_{3N+1}(\kappa )=\underset{l=1}{\overset{2N+1}{}}\underset{(n_1,n_2)J_{3l1}}{}b(\kappa ;N+l,n_1,n_2)=b(\kappa ;N+1,1,0)+b(\kappa ;N+1,0,2)+\mathrm{},$$
$`(6.3)`$
$$a_{3N+2}(\kappa )=\underset{l=1}{\overset{2N+2}{}}\underset{(n_1,n_2)J_{3l2}}{}b(\kappa ;N+l,n_1,n_2)=b(\kappa ;N+1,0,1)+\mathrm{},$$
$`(6.4)`$
where
$$J_l=\{(n_1,n_2):2n_1+n_2=l\}\text{N}_0\times \text{N}_0,l=0,1,\mathrm{}.$$
$`(6.5)`$
Inserting (6.1) and omitting terms of relative order $`N^1`$, we get from (6.2)
$$\begin{array}{cc}\hfill a_{3N}(\kappa )& C_{3N}(\kappa )\underset{l=0}{\overset{2N}{}}\underset{(n_1,n_2)J_{3l}}{}\underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+l+N)}{\mathrm{\Gamma }(\mu (\kappa )+N)}\hfill \\ & \times (2\kappa s_0)^l\frac{1}{(n_1p_1)!}(2\kappa t_{10})^{n_1p_1}\frac{1}{(n_2p_2)!}(2\kappa t_{20})^{n_2p_2},\hfill \end{array}$$
$`(6.6)`$
where
$$\begin{array}{cc}\hfill C_{3N}(\kappa )& =\frac{1}{\pi }\frac{\mathrm{\Gamma }(1+\mu (\kappa )+N)\mathrm{\Gamma }(\mu (\kappa )+N)}{\mathrm{\Gamma }(1+N)}(2\kappa s_0)^N\hfill \\ & \frac{1}{\pi }\mathrm{\Gamma }(\mu (\kappa )\mu (\kappa )+N)(2\kappa s_0)^N.\hfill \end{array}$$
$`(6.7)`$
Let us now discuss (6.6) in detail: Here only such values of $`n_1`$ and $`n_2`$ occur for which
$$\frac{2}{3}n_1+\frac{1}{3}n_2=l.$$
If again terms of relative order $`N^1`$ are omitted, the ratio of the gamma functions is equal to
$$N^{(2/3)p_1(1/3)p_2+l}=N^{(2/3)(n_1p_1)+(1/3)(n_2p_2)}.$$
The asymptotic behaviour of (6.6) then becomes
$$\begin{array}{cc}& a_{3N}(\kappa )C_{3N}(\kappa )\underset{l=0}{\overset{2N}{}}\underset{(n_1,n_2)J_{3l}}{}\underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\hfill \\ & \times \frac{1}{(n_1p_1)!}[N^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})]^{n_1p_1}\frac{1}{(n_2p_2)!}[N^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})]^{n_2p_2}.\hfill \end{array}$$
$`(6.8)`$
Here the integer power of $`2\kappa s_0`$ has been splitted in two factors with fractional powers, the meaning of which is given in terms of the phase factor $`\eta `$ of (5.10) by
$$(2\kappa s_0)^l=\{\begin{array}{c}(2s_0)^{(2/3)n_1}(2s_0)^{(1/3)n_2}\mathrm{if}\kappa =1\\ (2s_0)^{(2/3)n_1}\eta ^{2n_1}(2s_0)^{(1/3)n_2}\eta ^{n_2}\mathrm{if}\kappa =1\end{array}.$$
$`(6.9)`$
With $`n/3`$ in place of $`N`$, the terms in (6.8) look like the terms of the expansion of the exponential function
$$\begin{array}{cc}\hfill EX_0(n):& =\mathrm{exp}((\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20}))\hfill \\ & =1+(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{},\hfill \end{array}$$
$`(6.10)`$
but because of the three possible values of a third root there are two other such functions,
$$\begin{array}{cc}\hfill EX_1(n):& =\mathrm{exp}(\eta ^4(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^2(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20}))\hfill \\ & =1+\eta ^4(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^2(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{},\hfill \end{array}$$
$`(6.11)`$
$$\begin{array}{cc}\hfill EX_2(n):& =\mathrm{exp}(\eta ^2(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^4(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20}))\hfill \\ & =1+\eta ^2(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^4(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{}.\hfill \end{array}$$
$`(6.12)`$
Therefore (6.8) asymptotically shows exponential behaviour given by a certain linear combination of these three exponential functions. They can be identified by the constant term and the two linear terms of their expansion shown above. It is convenient, and easy because of the properties (5.11) of $`\eta `$, to introduce three functions which are linear combinations of the exponential functions such that only one of the identifying terms is different from zero,
$$\begin{array}{cc}& L_0(n):=\frac{1}{3}[EX_0(n)+EX_1(n)+EX_2(n)]=1+\mathrm{},\hfill \\ & L_1(n):=\frac{1}{3}[EX_0(n)+\eta ^2EX_1(n)+\eta ^4EX_2(n)]=(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\mathrm{},\hfill \\ & L_2(n):=\frac{1}{3}[EX_0(n)+\eta ^4EX_1(n)+\eta ^2EX_2(n)]=(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{}.\hfill \end{array}$$
$`(6.13)`$
We can now determine the asymptotic behaviour of (6.8) by looking at the contributions from $`(p_1,p_2)=(n_1,n_2)`$ or $`(p_1,p_2)=(n_11,n_2)`$ or $`(p_1,p_2)=(n_1,n_21)`$, respectively. This yields
$$a_{3N}(\kappa )C_{3N}(\kappa )\left\{S_0(\kappa )L_0(3N)+(2\kappa s_0)^{1/3}S_1(\kappa )L_1(3N)+(2\kappa s_0)^{2/3}S_2(\kappa )L_2(3N)\right\}.$$
$`(6.14)`$
In a similar way we may obtain from (6.3) and (6.4), respectively,
$$\begin{array}{cc}& a_{3N+1}(\kappa )C_{3N+1}(\kappa )\hfill \\ & \times \left\{(2\kappa s_0)^{2/3}S_2(\kappa )L_0(3N+1)+S_0(\kappa )L_1(3N+1)+(2\kappa s_0)^{1/3}S_1(\kappa )L_2(3N+1)\right\},\hfill \end{array}$$
$`(6.15)`$
$$\begin{array}{cc}& a_{3N+2}(\kappa )C_{3N+2}(\kappa )\hfill \\ & \times \left\{(2\kappa s_0)^{1/3}S_1(\kappa )L_0(3N+2)+(2\kappa s_0)^{2/3}S_2(\kappa )L_1(3N+2)+S_0(\kappa )L_2(3N+2)\right\}.\hfill \end{array}$$
$`(6.16)`$
If we now switch back to a representation in terms of the exponential functions (6.10)-(6.12), Stokes multipliers, according to (5.13) above, appear as their factors,
$$\begin{array}{cc}& a_{3N}(\kappa )\frac{1}{3}C_{3N}(\kappa )\hfill \\ & \times \left\{\sigma _0(\kappa )EX_0(3N)+\sigma _1(\kappa )EX_1(3N)+\sigma _2(\kappa )EX_2(3N)\right\},\hfill \\ & a_{3N+1}(\kappa )\frac{1}{3}C_{3N+1}(\kappa )\hfill \\ & \times \left\{\sigma _0(\kappa )EX_0(3N+1)+\eta ^2\sigma _1(\kappa )EX_1(3N+1)+\eta ^4\sigma _2(\kappa )EX_2(3N+1)\right\},\hfill \\ & a_{3N+2}(\kappa )\frac{1}{3}C_{3N+2}(\kappa )\hfill \\ & \times \left\{\sigma _0(\kappa )EX_0(3N+2)+\eta ^4\sigma _1(\kappa )EX_1(3N+2)+\eta ^2\sigma _2(\kappa )EX_2(3N+2)\right\}.\hfill \end{array}$$
$`(6.17)`$
These three asymptotic equations can be combined to give
Theorem 4. The leading terms of the asymptotic exponential behaviour, as $`n\mathrm{}`$, of the coefficients of the formal solutions are given by
$$\begin{array}{cc}& a_n(\kappa )\frac{1}{3}\pi ^1(2\kappa s_0)^{n/3}\mathrm{\Gamma }(\frac{1}{3}[\tau (\kappa )\tau (\kappa )+n])\hfill \\ & \times [\sigma _0(\kappa )EX_0(n)+\eta ^{2n}\sigma _1(\kappa )EX_1(n)+\eta ^{4n}\sigma _2(\kappa )EX_2(n)].\hfill \end{array}$$
$`(6.18)`$
This is a concrete example, with all the quantities determined explicitly, of the structural results obtained by Immink , and it is interesting also in the context of related work by other authors .
7. Postponed proofs
7.1 Proof of Lemma 1
Let
$$f(z)=z^\lambda u(z)$$
$`(7.1)`$
and multiply the differential equation for $`u(z)`$ by $`z^6`$ in order to remove all the negative powers of $`z`$. We then are concerned with the differential equation
$$\begin{array}{cc}\hfill L_zu(z):=& z^8u^{\prime \prime }+(2\lambda +1)z^7u^{}\hfill \\ & [\underset{m=1}{\overset{6}{}}D_mz^{6m}+(L^2\lambda ^2)z^6+\underset{m=1}{\overset{6}{}}B_mz^{6+m}]u(z)=0.\hfill \end{array}$$
$`(7.2)`$
We are looking for a solution of this differential equation in the form of an integral representation
$$u(z)=_{C_{t_2}}_{C_{t_1}}_{C_s}K(z;s,t_1,t_2)v(s,t_1,t_2)\text{d}s\text{d}t_1\text{d}t_2$$
$`(7.3)`$
with the kernel
$$K=K(z;s,t_1,t_2)=\mathrm{exp}(z^3s+z^2t_1+zt_2).$$
$`(7.4)`$
If we perform the differentiations with respect to $`z`$ under the integrals, the differential equation becomes
$$L_zu=vL_zK\text{d}s\text{d}t_1\text{d}t_2=0,$$
$`(7.5)`$
where
$$\begin{array}{cc}\hfill \frac{1}{K}L_zK& =(9s^2B_6)z^{12}+(12st_1B_5)z^{11}+(6st_2+4t_1^2B_4)z^{10}\hfill \\ & +([6\lambda +9]s+4t_1t_2B_3)z^9+([4\lambda +4]t_1+t_2^2B_2)z^8\hfill \\ & +([2\lambda +1]t_2B_1)z^7+(\lambda ^2L^2)z^6\underset{m=1}{\overset{6}{}}D_mz^{6m}.\hfill \end{array}$$
$`(7.6)`$
It is advisable to rewrite this in terms of the quantities $`p_1`$, $`p_2`$, $`p_3`$, according to (1.5), which determine the exponential factors of the formal solutions (1.2)-(1.3), or in terms of the related quantities $`s_0`$, $`t_{10}`$, $`t_{20}`$ according to (2.3). Then (7.6) becomes
$$\begin{array}{cc}\hfill \frac{1}{K}L_zK& =9(s^2s_0^2)z^{12}+12(st_1s_0t_{10})z^{11}+(6[st_2s_0t_{20}]+4[t_1^2t_{10}^2])z^{10}\hfill \\ & +([6\lambda +9]s+4[t_1t_2t_{10}t_{20}]+[4t_{10}t_{20}B_3])z^9\hfill \\ & +([4\lambda +4]t_1+[t_2^2t_{20}^2]+[t_{20}^2B_2])z^8\hfill \\ & +([2\lambda +1]t_2B_1)z^7+(\lambda ^2L^2)z^6\underset{m=1}{\overset{6}{}}D_mz^{6m}.\hfill \end{array}$$
$`(7.7)`$
By repeated partial integrations of the exponential function with respect to $`s`$ or $`t_1`$ or $`t_2`$, respectively, it is possible to get rid of all the powers of $`z`$. This lengthy task can more conveniently be performed in a formal way as follows: Let us find a partial differential expression $`M=M_{s,t_1,t_2}`$ with respect to $`s`$, $`t_1`$, $`t_2`$, independent of $`z`$, such that
$$L_zKM_{s,t_1,t_2}K.$$
$`(7.8)`$
The powers of $`z`$ correspond to partial derivatives with respect to $`s`$ or $`t_1`$ or $`t_2`$, and there seems to be some ambiguity as to the choice of the partial derivatives which yield the same total power of $`z`$. This ambiguity is resolved by the following reasonning: In order to get the order of each derivative as small as possible, we would prefer derivatives with respect to $`s`$, which account for $`z^3`$, with highest priority, next $`t_1`$, which accounts for $`z^2`$, last $`t_2`$, which accounts only for $`z`$. Conflicting with this policy, however, are some other requirements which ensure that we get a differential expression appropriate for our purpose. So the singularities should be at the right places, which are related to the coefficients of the exponential factor of the formal solutions we want to represent. In particular, bilinear or quadratic factors, such as $`t_1^2`$ for instance, should occur in the form $`t_1^2t_{10}^2`$. For this reason we have in (7.7) already added and subtracted the term $`4t_{10}^2`$. It is then necessary to treat these two terms differently, so that any power is always multiplied by the corresponding derivative. In the example just mentioned this means that $`z^{10}`$ has to be translated into $`^4/s^3t_2`$ for one term, but into $`^4/s^2t_1^2`$ for the other. The unique result now is
$$\begin{array}{cc}\hfill M_{s,t_1,t_2}& =9(s^2s_0^2)^4/s^4+12(st_1s_0t_{10})^4/s^3t_1\hfill \\ & +6(st_2s_0t_{20})^4/s^3t_2+4(t_1^2t_{10}^2)^4/s^2t_1^2\hfill \\ & +([6\lambda +9]s+[4t_{10}t_{20}B_3])^3/s^3+4[t_1t_2t_{10}t_{20}]^4/s^2t_1t_2\hfill \\ & +([4\lambda +4]t_1+[t_{20}^2B_2])^3/s^2t_1+[t_2^2t_{20}^2]^4/s^2t_2^2\hfill \\ & +([2\lambda +1]t_2B_1)^3/s^2t_2+(\lambda ^2L^2)^2/s^2\hfill \\ & D_1^2/st_1D_2^2/st_2D_3/sD_4/t_1D_5/t_2D_6.\hfill \end{array}$$
$`(7.9)`$
Next let us introduce the adjoint differential expression $`\overline{M}=\overline{M}_{s,t_1,t_2}`$ defined, with any sufficiently differentiable function $`v=v(s,t_1,t_2)`$, by
$$\begin{array}{cc}& \overline{M}_{s,t_1,t_2}v=(^4/s^4)\{9(s^2s_0^2)v\}+(^4/s^3t_1)\{12(st_1s_0t_{10})v\}\hfill \\ & +(^4/s^3t_2)\{6(st_2s_0t_{20})v\}+(^4/s^2t_1^2)\{4(t_1^2t_{10}^2)v\}\hfill \\ & (^3/s^3)\{([6\lambda +9]s+[4t_{10}t_{20}B_3])v\}+(^4/s^2t_1t_2)\{4(t_1t_2t_{10}t_{20})v\}\hfill \\ & (^3/s^2t_1)\{([4\lambda +4]t_1+[t_{20}^2B_2])v\}+(^4/s^2t_2^2)\{(t_2^2t_{20}^2)v\}\hfill \\ & (^3/s^2t_2)\{([2\lambda +1]t_2B_1)v\}+(\lambda ^2L^2)(^2v/s^2)\hfill \\ & D_1(^2v/st_1)D_2(^2v/st_2)+D_3(v/s)+D_4(v/t_1)+D_5(v/t_2)D_6v.\hfill \end{array}$$
$`(7.10)`$
The difference as compared with $`M`$ is that the factors in front of each derivative have here to be differentiated too and that all the terms of odd order change their sign. The usefulness of the adjoint expression lies in the formula, known in the one-variable case as the identity of Lagrange,
$$vMKK\overline{M}v=RHS,$$
$`(7.11)`$
where the right-hand side $`RHS`$ is a lengthy expression, bilinear in $`v`$ and $`K`$ or their partial derivatives, which may be so arranged that it consists of a sum of terms of which each is a total derivative with respect to one of the variables. For the term with the factor $`D_1`$, for instance, this reads
$$v\frac{^2K}{st_1}K\frac{^2v}{st_1}=\frac{1}{2}\frac{}{s}[v\frac{K}{t_1}K\frac{v}{t_1}]+\frac{1}{2}\frac{}{t_1}[v\frac{K}{s}K\frac{v}{s}].$$
$`(7.12)`$
By means of (7.5) and (7.8), the equation to be satisfied now is
$$\begin{array}{cc}\hfill L_zu=& vM_{s,t_1,t_2}K\text{d}s\text{d}t_1\text{d}t_2\hfill \\ \hfill =& K\overline{M}_{s,t_1,t_2}v\text{d}s\text{d}t_1\text{d}t_2+RHS\text{d}s\text{d}t_1\text{d}t_2=0.\hfill \end{array}$$
$`(7.13)`$
The first term in the second line can be made to vanish if $`v`$ is required to be a solution of the partial differential equation
$$\overline{M}_{s,t_1,t_2}v(s,t_1,t_2)=0,$$
$`(7.14)`$
and the second term by a suitable choice of the contours of integration, for it is a sum of semi-integrated terms, each involving the difference of the values of the integrand at the termini of the contour of one variable and only two remaining integrals with respect to the other two variables. The partial differential equation (2.2) is the same as (7.14), after the derivatives of the products in (7.10) have been resolved. This completes the proof of Lemma 1.
7.2 Proof of Lemma 2
In terms of the shifted variables
$$S=s\kappa s_0,T_1=t_1\kappa t_{10},T_2=t_2\kappa t_{20},$$
$`(7.15)`$
the differential equation (2.2) reads
$$\begin{array}{cc}& 9S(S+2\kappa s_0)(^4v/S^4)+12(ST_1+\kappa t_{10}S+\kappa s_0T_1)(^4v/S^3T_1)\hfill \\ & +6(ST_2+\kappa t_{20}S+\kappa s_0T_2)(^4v/S^3T_2)+4T_1(T_1+2\kappa t_{10})(^4v/S^2T_1^2)\hfill \\ & +T_2(T_2+2\kappa t_{20})(^4v/S^2T_2^2)+4(T_1T_2+\kappa t_{20}T_1+\kappa t_{10}T_2)(^4v/S^2T_1T_2)\hfill \\ & +([816\lambda ]\kappa s_0+B_34t_{10}t_{20})(^3v/S^3)+(816\lambda )S(^3v/S^3)\hfill \\ & +([524\lambda ]\kappa t_{10}+B_2t_{20}^2)(^3v/S^2T_1)+(524\lambda )T_1(^3v/S^2T_1)\hfill \\ & +([252\lambda ]\kappa t_{20}+B_1)(^3v/S^2T_2)+(252\lambda )T_2(^3v/S^2T_2)\hfill \\ & +([\lambda 12]^2L^2)(^2v/S^2)D_1(^2v/ST_1)D_2(^2v/ST_2)\hfill \\ & +D_3(v/S)+D_4(v/T_1)+D_5(v/T_2)D_6v=0.\hfill \end{array}$$
$`(7.16)`$
Inserting a power series solution
$$v=\underset{m=0}{}\underset{n_1=0}{}\underset{n_2=0}{}a(m,n_1,n_2)S^{\mu +m}T_1^{\nu _1n_1}T_2^{\nu _2n_2},$$
$`(7.17)`$
we obtain for the coefficients the recurrence relation
$$\begin{array}{cc}& 6\kappa s_0(\mu +m)(\mu +m1)(\mu +m2)[3(\mu +m)2(\nu _1+n_1)(\nu _2+n_2)\lambda \tau (\kappa )+6]\hfill \\ & \times a(m,n_1,n_2)\hfill \\ & +(\mu +m1)(\mu +m2)\{[3(\mu +m)2(\nu _1+n_1)(\nu _2+n_2)\lambda +3]^2L^2\}\hfill \\ & \times a(m1,n_1,n_2)\hfill \\ & +(\mu +m1)(\mu +m2)(\nu _1+n_11)\{4\kappa t_{10}[3(\mu +m)+2(\nu _1+n_1)+(\nu _2+n_2)\hfill \\ & +\lambda 4]+t_{20}^2B_2\}a(m1,n_11,n_2)\hfill \\ & +(\mu +m1)(\mu +m2)(\nu _2+n_21)\{2\kappa t_{20}[3(\mu +m)+2(\nu _1+n_1)+(\nu _2+n_2)\hfill \\ & +\lambda \frac{7}{2}]+B_1\}a(m1,n_1,n_21)\hfill \\ & +D_1(\mu +m2)(\nu _1+n_11)a(m2,n_11,n_2)\hfill \\ & +D_2(\mu +m2)(\nu _2+n_21)a(m2,n_1,n_21)\hfill \\ & +D_3(\mu +m2)a(m2,n_1,n_2)D_4(\nu _1+n_11)a(m3,n_11,n_2)\hfill \\ & D_5(\nu _2+n_21)a(m3,n_1,n_21)D_6a(m3,n_1,n_2)=0,\hfill \end{array}$$
$`(7.18)`$
valid for $`m0`$, $`n_10`$, $`n_20`$ provided that we agree that all the $`a`$-coefficients are equal to zero if any of the indices $`m,n_1,n_2`$ is less than zero. Assuming that $`a(0,0,0)0`$, we get from the equation for $`m=n_1=n_2=0`$ the indicial equation
$$\mu (\mu 1)(\mu 2)(3\mu 2\nu _1\nu _2\lambda \tau (\kappa )+6)=0,$$
$`(7.19)`$
with $`\tau (\kappa )`$ according to (1.6). Possible values of the exponent $`\mu `$ are therefore $`\mu =0,1,2,`$ or
$$\mu =\frac{1}{3}[2\nu _1+\nu _2+\lambda +\tau (\kappa )]2.$$
$`(7.20)`$
In order to avoid complications, we may assume that the last possibility does not yield an integer value. This can always be guarantied be a suitable choice of the still disposable parameter $`\lambda `$. Inspection of the recurrence relation then shows that each of the possible values of $`\mu `$ leads to a solution of the partial differential equation. We here need not further consider the three solutions which are regular with respect to $`S`$ at $`S=0`$, but we continue to discuss the singular one with the exponent (7.20), writing for the associated coefficients $`A(\kappa ;m,n_1,n_2)`$ rather than $`a(m,n_1,n_2)`$. Simplifying by means of (7.20) and introducing the $`b`$-coefficients according to (2.6) , we get (2.7). The coefficients for which $`3m2n_1n_2=0`$ are constants of integration and may be chosen to be zero. This completes the proof of Lemma 2.
7.3 Proof of the limit formula
We have to verify that the limit formula (4.14) satisfies (4.13 ). Substituting it for the $`e`$-coefficients in (4.13) and interchanging the summations, we have to evaluate sums such as
$$\underset{p=j}{\overset{n}{}}\frac{(1)^{np}}{(pj)!(np)!}=\frac{(1)^{nj}}{(nj)!}\underset{q=0}{\overset{nj}{}}(1)^q\left(\genfrac{}{}{0pt}{}{nj}{q}\right)=\{\begin{array}{c}1\mathrm{if}j=n\\ 0\mathrm{if}0jn1,n>0\end{array}.$$
$`(7.21)`$
Therefore only one term survives on the right-hand side, which then becomes equal to the left.
7.4 Proof of Theorem 1
In order to prove Theorem 1, we first multiply the continuation formula (3.5), with $`r_1=r_2=0`$, by
$$(\frac{1}{2}\frac{s}{2\kappa s_0})^{\mu (\kappa )}=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(\mu (\kappa ))_j}{j!}(\frac{1}{2}+\frac{s}{2\kappa s_0})^j,$$
$`(7.22)`$
where the left-hand side is used on the left and the right-hand side on the right of the continuation formula. Then the left-hand side is the same as above (4.1) with the power factor in front of the series removed, but the right becomes, after multiplication of the two power series,
$$\begin{array}{cc}& \underset{q_1=0}{\overset{\mathrm{}}{}}\underset{q_2=0}{\overset{\mathrm{}}{}}E(\kappa ;0,0;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \underset{k=0}{\overset{\mathrm{}}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\mathrm{\Gamma }(\nu _1+q_1+l_1)\mathrm{\Gamma }(\nu _2+q_2+l_2)H(\kappa ;k,l_1,l_2;q_1,q_2)\hfill \\ & \times (\frac{1}{2}+\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2+k}(t_1+\kappa t_{10})^{\nu _1q_1l_1}(t_2+\kappa t_{20})^{\nu _2q_2l_2},\hfill \end{array}$$
$`(7.23)`$
with $`H`$ defined in (4.16). Proceeding as above, and making use of the formula
$$(\alpha k)_m=(\alpha )_m\frac{(1\alpha )_k}{(1\alpha m)_k},$$
$`(7.24)`$
we obtain
$$\begin{array}{cc}& \frac{\pi }{\mathrm{sin}(\pi \mu (\kappa ))}\frac{1}{\mathrm{\Gamma }(1+\mu (\kappa )+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{q_1=0}{\overset{n_1}{}}\underset{q_2=0}{\overset{n_2}{}}E(\kappa ;0,0;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{m!}\hfill \\ & \times \underset{k=0}{\overset{\mathrm{}}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_k}{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2m)_k}H(\kappa ;k,l_1,l_2;q_1,q_2)\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1+q_1+l_1=n_1}\left(\frac{(1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2+q_2+l_2=n_2}.\hfill \end{array}$$
$`(7.25)`$
Keeping the first $`K+1`$ singular terms on the right and solving for the $`E`$-coefficient with $`q_1=n_1`$, $`q_2=n_2`$, we have
$$\begin{array}{cc}& E(\kappa ;0,0;n_1,n_2)\hfill \\ & \times [1+\underset{k=1}{\overset{K}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2)_k}{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2m)_k}H(\kappa ;k,0,0;n_1,n_2)+O(m^{K1})]\hfill \\ & =\frac{\pi }{\mathrm{sin}(\pi \mu (\kappa ))}\frac{m!}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{(q_1,q_2)(n_1,n_2)}{\underset{q_1=0}{\overset{n_1}{}}\underset{q_2=0}{\overset{n_2}{}}}E(\kappa ;0,0;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times [1+\underset{k=1}{\overset{K}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_k}{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2m)_k}H(\kappa ;k,l_1,l_2;q_1,q_2)+O(m^{K1})]\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1+q_1+l_1=n_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2+q_2+l_2=n_2}.\hfill \end{array}$$
$`(7.26)`$
This is essentially (4.15) because of (4.11) and completes the proof of Theorem 1.
7.5 Choice of the computational parameter $`\lambda `$
If $`D_3=D_6=0`$, the recurrence relation (2.7) for $`b(\kappa ;m,0,0)`$ reduces to a two-term relation, and we obtain
$$b(\kappa ;m,0,0)=(2\kappa s_0)^m\frac{(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_m(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_m}{m!}.$$
$`(7.27)`$
Rewriting (4.16) by means of the identity
$$\frac{(x)_{kj}}{(kj)!}=\frac{(x)_k}{k!}\frac{(k)_j}{(1xk)_j},$$
$`(7.28)`$
we then have
$$H(\kappa ;k,0,0;0,0)=\frac{(\mu (\kappa ))_k}{k!}\underset{j=0}{\overset{k}{}}\frac{(k)_j(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_j(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_j}{(1\mu (\kappa )k)_j(1+\mu (\kappa ))_jj!}.$$
$`(7.29)`$
Because of $`3\mu (\kappa )=\lambda +\tau (\kappa )3`$ according to (2.5) and (2.11) and $`\tau (\kappa )=3\tau (\kappa )`$ according to (1.6), the series is a terminating one-balanced hypergeometric series at unit argument, which can be summed by the theorem of Saalschütz , so that
$$H(\kappa ;k,0,0;0,0)=\frac{(\frac{1}{3}L+\frac{1}{3}\lambda )_k(\frac{1}{3}L+\frac{1}{3}\lambda )_k}{(1+\mu (\kappa ))_kk!}.$$
$`(7.30)`$
In total, we get
$$\begin{array}{cc}& e(\kappa ;0,0)=\frac{\pi }{\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))}\frac{\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa )+m)\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa )+m)}{\mathrm{\Gamma }(\frac{1}{3}\lambda +\frac{1}{3}\tau (\kappa )+m)\mathrm{\Gamma }(\frac{1}{3}\lambda +\frac{1}{3}\tau (\kappa )+m)}\hfill \\ & \times \left[1+\underset{k=1}{\overset{K}{}}\frac{(\frac{1}{3}L+\frac{1}{3}\lambda )_k(\frac{1}{3}L+\frac{1}{3}\lambda )_k}{(\frac{1}{3}\lambda +\frac{1}{3}\tau (\kappa )m)_kk!}+O(m^{K1})\right]^1.\hfill \end{array}$$
$`(7.31)`$
For $`\lambda =L`$ or $`\lambda =L`$, the terms with $`k=1,2,\mathrm{}`$ of the asymptotic series all vanish and the factor in front of the series becomes independent of $`m`$. More generally this means that some terms which might become quite large when L is not small can be removed by such a choice of $`\lambda `$ from the asymptotic series and incorporated in the $`m`$-dependence of the function in front of the series.
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Wolfgang Bühring
Physikalisches Institut
Universität Heidelberg
Philosophenweg 12
69120 Heidelberg
GERMANY
buehring@physi.uni-heidelberg.de |
warning/0003/cs0003077.html | ar5iv | text | # DATALOG with constraints — an answer-set programming system
## Introduction
Many important computational problems in combinatorial optimization, constraint satisfaction and artificial intelligence can be cast as search problems. Answer-set programming (ASP) (??) was recently identified as a declarative programming paradigm appropriate for such applications. Logic programming with the stable-model semantics (stable logic programming, for short) was proposed as an embodiment of this paradigm. Disjunctive logic programming with the answer-set semantics is another implementation of ASP currently under development (?). Early experimental results demonstrate the potential of answer-set programming approaches in such areas as planning and constraint satisfaction (???).
In this paper we describe another formalism that implements the ASP approach. We call it DATALOG with constraints and denote by $`\mathrm{DC}`$. Our goal is to design an ASP system with a semantics more readily understandable than the semantics of stable models. We seek a semantics that would be as close as possible to propositional satisfiability yet as expressive and as effective, especially from the point of view of conciseness of representations and time performance, as the stable logic programming. We argue that $`\mathrm{DC}`$ has a potential to become a practical declarative programming tool. We show that it yields intuitive and small-size encodings, we characterize its complexity and expressive power and present computational experiments demonstrating its effectiveness.
Answer-set programming is a paradigm in which programs are built as theories in some formal system $``$ with a well-defined syntax, and with a semantics that assigns to a theory $`P`$ in the system a collection of subsets of some domain. These subsets are referred to as answer sets of $`P`$ and specify the results of computation based on $`P`$. To solve a problem $`\mathrm{\Pi }`$ in an ASP formalism, we find a program $`P`$ so that the solutions to $`\mathrm{\Pi }`$ can be reconstructed, in polynomial (ideally, linear) time, from the answer sets to $`P`$.
The definition of the answer-set programming given above is very general. Essentially any logic formalism can be a basis for an answer-set programming system. For instance, the propositional logic gives rise to an ASP system: programs are collections of propositional clauses, their models are answer sets. To solve, say, a planning problem, we encode the constraints of the problem as propositional clauses in such a way that legal plans are determined by models of the resulting propositional theory. This approach, called satisfiability planning, received significant attention lately and was shown to be quite effective (???).
Recently, several implementations of the ASP approach were developed that are based on nonmonotonic logics such as smodels (?), for stable logic programming, dlv (?), for disjunctive logic programming with answer-set semantics, and deres (?), for default logic with Reiter’s extensions. All these systems have been extensively studied. Promising experimental results concerning their performance were reported (???).
The question arises which formal logics are appropriate as bases of answer-set programming implementations. To discuss such a general question one needs to formulate quality criteria with respect to which ASP systems can be compared. At the very least, these criteria should include:
1. expressive power
2. time performance
3. simplicity of the semantics
4. ease of coding, conciseness of programs.
We will discuss these criteria in detail elsewhere. We will make here only a few brief comments on the matter. From the point of view of the expressive power all the systems that we discussed are quite similar. Propositional logic and stable logic programming are well-attuned to the class NP (?). Disjunctive logic programming and default logic capture the class $`\mathrm{\Sigma }_P^2`$ (??). However, this distinction is not essential as recently pointed out in (?). The issue of time performance can be resolved only through comprehensive experimentation and this work is currently under way.
As concerns inherent complexity of the system and intuitiveness of the semantics, ASP systems based on the propositional logic seem to be clear winners. However, propositional logic is monotone and modeling indefinite information and phenomena such as the frame problem is not quite straightforward. In applications involving the computation of transitive closures, as in the problem of existence of hamilton cycles, it leads to programs that are large and, thus, difficult to process. In this respect, ASP systems based on nonmonotonic logics have an edge. They were designed to handle incomplete and indefinite information. Thus, they often yield more concise programs. However, they require more elaborate formal machinery and their semantics are more complex.
Searching for the middle ground between systems such as logic programming with stable model semantics and propositional logic, we propose here a new ASP formalism, $`\mathrm{DC}`$. Our guiding principle was to design a system which would lead to small-size encodings, believing that small theories will lead to more efficient solutions. We show that $`\mathrm{DC}`$ is nonmonotonic, has the same expressive power as stable logic programming but that its semantics stays closer to that of propositional logic. Thus, it is arguably simpler than the stable-model semantics. We present experimental results that demonstrate that $`\mathrm{DC}`$ is competitive with ASP implementations based on nonmonotonic logics (we use smodels for comparison) and those based on propositional logics (we use csat (?) in our experiments). Our results strongly indicate that formalisms which provide smaller-size encodings are more effective as practical search-problem solvers.
## DATALOG with constraints
A $`\mathrm{DC}`$ theory (or program) consists of constraints and Horn rules (DATALOG program). This fact motivates out choice of terminology — DATALOG with constraints. We start a discussion of $`\mathrm{DC}`$ with the propositional case. Our language is determined by a set of atoms $`At`$. We will assume that $`At`$ is of the form $`At=At_CAt_H`$, where $`At_C`$ and $`At_H`$ are disjoint.
A DC theory (or DC program) is a triple $`T=(T_C,T_H,T_{PC})`$, where
1. $`T_C`$ is a set of propositional clauses $`\neg a_1\mathrm{}\neg a_mb_1\mathrm{}b_n`$ such that all $`a_i`$ and $`b_j`$ are from $`At_C`$,
2. $`T_H`$ is a set of Horn rules $`a_1\mathrm{}a_mb`$ such that $`bAt_H`$ and all $`a_i`$ are from $`At`$,
3. $`T_{PC}`$ is a set of clauses over $`At`$.
By $`At(T)`$, $`At_C(T)`$ and $`At_{PC}(T)`$ we denote the set of atoms from $`At`$, $`At_C`$ and $`At_{PC}`$, respectively, that actually appear in $`T`$.
With a $`\mathrm{DC}`$ theory $`T=(T_C,T_H,T_{PC})`$ we associate a family of subsets of $`At_C(T)`$. We say that a set $`MAt_C(T)`$ satisfies $`T`$ (is an answer set of $`T`$) if
1. $`M`$ satisfies all the clauses in $`T_C`$, and
2. the closure of $`M`$ under the Horn rules in $`T_H`$, $`M^c=LM(T_HM)`$ satisfies all clauses in $`T_{PC}`$ ($`LM(P)`$ denotes the least model of a Horn program $`P`$).
Intuitively, the collection of clauses in $`T_C`$ can be thought of as a representation of the constraints of the problem, Horn rules in $`T_H`$ can be viewed as a mechanism to compute closures of sets of atoms satisfying the constraints in $`T_C`$, and the clauses in $`T_{PC}`$ can be regarded as constraints on closed sets (we refer to them as post-constraints). A set of atoms $`MAt_C(T)`$ is a model if it (propositionally) satisfies the constraints in $`T_C`$ and if its closure (propositionally) satisfies the constraints in $`T_{PC}`$. Thus, the semantics of $`\mathrm{DC}`$ retains much of the simplicity of the semantics of propositional logic.
$`\mathrm{DC}`$ can be used as a computational tool to solve search problems. We define a search problem $`\mathrm{\Pi }`$ to be determined by a set of finite instances, $`D_\mathrm{\Pi }`$, such that for each instance $`ID_\mathrm{\Pi }`$, there is a finite set $`S_\mathrm{\Pi }(I)`$ of all solutions to $`\mathrm{\Pi }`$ for the instance $`I`$. For example, the problem of finding a hamilton cycle in a graph is a search problem: graphs are instances and for each graph, its hamilton cycles (sets of their edges) are solutions. A $`\mathrm{DC}`$ theory $`T=(T_C,T_H,T_{PC})`$ solves a search problem $`\mathrm{\Pi }`$ if solutions to $`\mathrm{\Pi }`$ can be computed (in polynomial time) from answer sets to $`T`$. Propositional logic and stable logic programming are used as problem solving formalisms following the same general paradigm. To illustrate all the concepts introduced here and show how $`\mathrm{DC}`$ programs can be built by modeling problem constraints, we will now present a $`\mathrm{DC}`$ program that solves the hamilton-cycle problem.
Consider a directed graph $`G`$ with the vertex set $`V`$ and the edge set $`E`$. Consider a set of atoms $`\{hc(a,b):(a,b)E\}`$. An intuitive interpretation of an atom $`hc(a,b)`$ is that the edge $`(a,b)`$ is in a hamilton cycle. Include in $`T_C`$ all clauses of the form $`\neg hc(b,a)\neg hc(c,a)`$, where $`a,b,cV`$, $`bc`$ and $`(b,a),(c,a)E`$. In addition, include in $`T_C`$ all clauses of the form $`\neg hc(a,b)\neg hc(a,c)`$, where $`a,b,cV`$, $`bc`$ and $`(a,b),(a,c)E`$. Clearly, the set of propositional variables of the form $`\{hc(a,b):(a,b)F\}`$, where $`FE`$, satisfies all clauses in $`T_C`$ if and only if no two distinct edges in $`F`$ end in the same vertex and no two distinct edges in $`F`$ start in the same vertex. In other words, $`F`$ spans a collection of paths and cycles in $`G`$.
To guarantee that the edges in $`F`$ define a hamilton cycle, we must enforce that all vertices of $`G`$ are reached by means of the edges in $`F`$ if we start in some (arbitrarily chosen) vertex of $`G`$. This can be accomplished by means of a simple Horn program. Let us choose a vertex, say $`s`$, in $`G`$. Include in $`T_H`$ the Horn rules $`hc(s,t)vstd(t)`$, for every edge $`(s,t)`$ in $`G`$. In addition, include in $`T_H`$ Horn rules $`vstd(t),hc(t,u)vstd(u)`$, for every edge $`(t,u)`$ of $`G`$ not starting in $`s`$. Clearly, the least model of $`FT_H`$, where $`F`$ is a subset of $`E`$, contains precisely these variables of the form $`vstd(t)`$ for which $`t`$ is reachable from $`s`$ by a nonempty path spanned by the edges in $`F`$. Thus, $`F`$ is the set of edges of a hamilton cycle of $`G`$ if and only if the least model of $`FT_H`$, contains variable $`vstd(t)`$ for every vertex $`t`$ of $`G`$. Let us define $`T_{PC}=\{vstd(t):tV\}`$ and $`T_{ham}(G)=(T_C,T_H,T_{PC})`$. It follows that hamilton cycles of $`G`$ can be reconstructed (in linear time) from answer sets to the $`\mathrm{DC}`$ theory $`T_{ham}(G)`$. In other words, to find a hamilton cycle in $`G`$, it is enough to find an answer set for $`T_{ham}(G)`$.
This example illustrates the simplicity of the semantics — it is only a slight adaptation of the semantics of propositional logic to the case when in addition to propositional clauses we also have Horn rules in theories. It also illustrates the power of $`\mathrm{DC}`$ to generate concise encodings. All known propositional encodings of the hamilton-cycle problem require that additional variables are introduced to “count” how far from the starting vertex an edge is located. Consequently, propositional encodings are much larger and lead to inefficient computational approaches to the problem. We present experimental evidence to this claim later in the paper.
The question arises which search problems can be represented (and solved) by means of finding answer sets to appropriate $`\mathrm{DC}`$ programs. In general, the question remains open. We have an answer, though, if we restrict our attention to the special case of decision problems. Consider a $`\mathrm{DC}`$ theory $`T=(T_C,T_H,T_{PC})`$, where $`T_H=T_{PC}=\mathrm{}`$. Clearly, $`M`$ is an answer set for $`T`$ if and only if $`M`$ is a model of the collection of clauses $`T_C`$. Thus, the problem of existence of an answer set is at least as hard as the propositional satisfiability problem. On the other hand, for every $`\mathrm{DC}`$ theory $`T`$ and for every set $`MAt_C(T)`$, it can be checked in linear time whether $`M`$ is an answer set for $`T`$. Thus, we obtain the following complexity result.
###### Theorem 1
The problem of existence of an answer set for a finite propositional $`\mathrm{DC}`$ theory $`T`$ is NP-complete.
It follows that every problem in NP can be polynomially reduced to the problem of existence of an answer set for a propositional $`\mathrm{DC}`$ program. Thus, given a problem $`\mathrm{\Pi }`$ in NP, for every instance $`I`$ of $`\mathrm{\Pi }`$, $`\mathrm{\Pi }`$ can be decided by deciding the existence of an answer set for the $`\mathrm{DC}`$ program corresponding to $`\mathrm{\Pi }`$ and $`I`$.
Propositional $`\mathrm{DC}`$ can be extended to the predicate case. It is important as it significantly simplifies the task of developing programs for solving problems with $`\mathrm{DC}`$. In the example discussed above, the theory $`T_{ham}(G)`$ depends heavily on the input. Each time we change the input graph, a different theory has to be used. However, when constructing predicate $`\mathrm{DC}`$-based solutions to a problem $`\mathrm{\Pi }`$, it is often possible to separate the representation of an instance (input) to $`\mathrm{\Pi }`$ from that of the constraints that define $`\mathrm{\Pi }`$. As a result only one (predicate) program describing the constraints of $`\mathrm{\Pi }`$ needs to be written. Specific input for the program, say $`I`$, can be described separately as a collection of facts (according to some uniform schema). Both parts together can be combined to yield a $`\mathrm{DC}`$ program whose answer sets determine solutions to $`\mathrm{\Pi }`$ for the input $`I`$. Such an approach, we will refer to it as uniform, is often used in the context of DATALOG, DATALOG<sup>¬</sup> or logic programming to study complexity of these systems as query languages. The part representing input is referred to as the extensional database. The part representing the query or the problem is called the intensional database or program. Due to the space limitations we do not discuss the details of the predicate case here. They will be given in the full version of the paper. We only state a generalization of Theorem 1.
###### Theorem 2
The expressive power of $`\mathrm{DC}`$ is the same as that of stable logic programming. In particular, a decision problem $`\mathrm{\Pi }`$ can be solved uniformly in $`\mathrm{DC}`$ if and only if $`\mathrm{\Pi }`$ is in the class NP.
## Implementation
Some types of constraints appear frequently in applications. For instance, when defining plans we may want to specify a constraint that says that exactly one action from the set of allowed actions be selected at each step. Such constraints can be modeled by collections of clauses. To make sure $`\mathrm{DC}`$ programs are as easy to write and as concise as possible we have extended the syntax of $`\mathrm{DC}`$ by providing explicit ways to model constraints of the form “select at least (at most, exactly) $`k`$ elements from a set”. Having these constraints results in shorter programs which, as we believe, has a significant positive effect of the performance of our system.
An example of a select constraint with a short explanation is presented here. Let $`PRED`$ be the set of predicates occurring in the IDB. For each variable $`X`$ declared in the IDB the range $`R(X)`$ of $`X`$ is determined by the EDB.
where $`n,m`$ are nonnegative integers such that $`nm,qPRED`$ and $`p_1,\mathrm{},p_i`$ are EDB predicates or logical conditions (logical conditions can be comparisons of arithmetic expressions or string comparisons). The interpretation of this constraint is as follows: for every $`\stackrel{}{x}R(\stackrel{}{X})`$ at least $`n`$ atoms and at most $`m`$ atoms in the set $`\{q(\stackrel{}{x},\stackrel{}{y}):\stackrel{}{y}R(\stackrel{}{Y})\}`$ are true.
We implemented $`\mathrm{DC}`$ in the predicate setting. Thus, our system consists of two main modules. The first of them, referred to as grounder, converts a predicate $`\mathrm{DC}`$ program (consisting of both the extensional and intensional parts) into the corresponding propositional $`\mathrm{DC}`$ program. The second module, $`\mathrm{DC}`$ solver, denoted dcs, finds the answer sets to propositional $`\mathrm{DC}`$ programs. Since we focus on the propositional case here, we only describe the key ideas behind the $`\mathrm{DC}`$ solver, dcs.
The $`\mathrm{DC}`$ solver uses a Davis-Putnam type approach, with backtracking, propagation and lookahead (also called literal testing), to deal with constraints represented as clauses, select constraints and Horn rules, and to search for answer sets. The lookahead in $`\mathrm{DC}`$ is similar to local processing performed in csat (?). However, we use different methods to determine how many literals to consider in the lookahead phase. Other techniques, especially propagation and search heuristics, were designed specifically for the case of $`\mathrm{DC}`$ as they must take into account the presence of Horn rules in programs.
The lookahead procedure selects a number of literals which have not yet been assigned a value. For each such literal, the procedure tries both truth values: true and false. For each assignment, the theory is evaluated using propagation. If in both cases a contradiction is reached, then it is necessary to backtrack. If for only one evalution a conflict is reached, then the literal is assigned the other truth value and we proceed to the next step. If neither evaluation results in a contradiction, we cannot assign a truth value to this literal but we save the data such as the number of forced literals and the number of clauses satisfied, computed during propagation.
Clearly, if all unassigned literals were tested it would prune the most search space. At the same time, the savings might not be large enough to compensate for the increase in the running time caused by extensive lookahead. Thus, we select only a portion of all unassigned literals for lookahead. The number of literals to consider was established empirically (it does not depend on the size of the theory). Since not all literals are selected, it is important to focus on those literals that are likely to result in a contradiction for at least one of the truth values. In our implementation, we select the most constrained literals, as determined by their weights.
Specifically, each constraint is assigned a weight based on its current length and types (recall that in addition to propositional clauses, we also allow other types of constraints, e.g., select constraints). The shorter the constraint the greater its weight. Also, certain types of constraints force more assignments on literals and are given a greater weight than other constraints of the same length. Every time a literal appears in an unsatisfied constraint, the weight of that literal is incremented by the weight of the clause.
After testing a predetermined number of literals without finding a forced truth assignment and without backtracking, the information computed during propagation is used to choose the next literal for which both possible truth assignments have to be tested (branching literal). The choice of the next branching literal is based on an approximation of which literal, once assigned a truth value, will force the truth assignments onto the largest number of other literals and will satisfy the largest number of clauses. Using the data computed during propagation gives more accurate information on which to base such approximations. The methods used for determining which literals to select in the lookahead and which data to collect and save during the propagation phase are two key ways in which the literal testing procedure differs from the local processing of csat.
## Experimentation
We compared the performance of $`\mathrm{DC}`$ solver dcs with smodels, a system for computing stable models of logic programs (?), and csat, a system for testing propositional satisfiability (?). In the case of smodels we used version 2.24 in conjunction with the grounder lparse, version 0.99.41. These versions of lparse and smodels implement the expressive rules described by (?). The expressive rules were used whenever applicable during the testing. The programs were all executed on a Sun SparcStation 20. For each test we report the cpu user times for processing the corresponding propositional program or theory. We tested all three system to compute hamilton cycles and colorings in graphs, to solve the $`N`$-queens problem, to prove that the pigeonhole problem has no solution if the number of pigeons exceeds the number of holes, and to compute Schur numbers.
The Hamilton cycle problem has already been described. We randomly generated one thousand graphs with the edge-to-vertex ratio such that $`50\%`$ of the graphs contained Hamilton cycles (crossover region). The number of vertices ranged from 30 to 80. We used encodings of the problem as a $`\mathrm{DC}`$ program, logic program (in smodels syntax) and as a propositional theory. dcs performed better than smodels and smodels performed significantly better than csat (Fig. 1). We believe that a major factor behind poorer performance of csat is that all known propositional encodings of the hamilton cycle problem are much larger than those possible with $`\mathrm{DC}`$ or logic programs (under the stable model semantics). Propositional encodings, due to their size, rendered csat not practical to execute for graphs with more than 40 vertices.
The $`N`$-queens problem consists of finding a placement of $`N`$ queens on an $`N\times N`$ board such that no queen can remove another. Both csat and dcs execute in much less time than smodels for these problems (Fig. 2). Again the size of the encoding seems to be a major factor. One thing to consider in this case is that the number of rules for smodels is approximately five times that for $`\mathrm{DC}`$ and more than twice that of propositional encodings.
The Schur problem consists of placing $`N`$ numbers $`1,2,\mathrm{},N`$ in $`B`$ bins such that no bin is closed under sums. That is, for all numbers $`x`$, $`y`$, $`z`$, $`1x,y,zN`$, if $`x`$ and $`y`$ are the same bin, then $`z`$ is not ($`x`$ and $`y`$ need not be distinct). The Schur number $`S(B)`$ is the maximum number $`N`$ for which such a placement is still possible. It is known to exist for every $`B1`$. We considered the problem of the existence of the placement for $`B=3`$ and $`N=13`$ and $`14`$, and for $`B=4`$ and $`N=43,44`$ and $`45`$. In each case we used all three systems to process the corresponding encodings. The results are shown in Fig. 3. It follows that $`S(3)=13`$ and $`S(4)=44`$. Again, dcs outperforms both smodels and csat.
Results for graph 3-coloring for graphs with the number of vertices ranging from $`50`$ to $`300`$ are shown in Fig. 4 (for every choice of the number of vertices, 100 graphs from the crossover region were randomly generated). Both dcs and csat performed better than smodels. Again the size of the theory seems to be a factor. The CNF theory for coloring is smaller than a logic program encoding the same problem. The sizes of propositional and $`\mathrm{DC}`$ encodings are similar.
Results for the pigeonhole placement problem show a similar performance of all three algorithms, with csat doing slightly better than the others and dcs outperforming (again only slightly) smodels.
## Conclusions
We described a new system, $`\mathrm{DC}`$, for solving search problems. We designed $`\mathrm{DC}`$ so that its semantics was as close as possible to that of propositional logic. Our goal was to design a system that would result is short problem encodings. Thus, we provided constructs for some frequently occurring types of constraints and we built into $`\mathrm{DC}`$ elements of nonmonotonicity by including Horn rules in the syntax. As a result, $`\mathrm{DC}`$ programs encoding search problems are often much smaller than those possible with propositional theories. Experimental results show that dcs often outperforms systems based on propositional satisfiability as well as systems based on nonmonotonic logics, and that it constitutes a viable approach to solving problems in AI, constraint satisfaction and combinatorial optimization. We believe that our focus on short programs is the key to the success of $`\mathrm{DC}`$ and its reasoning engine dcs. Our results show that when building general purpose solvers of search problems, the size of encodings should be a key design factor. |
warning/0003/quant-ph0003112.html | ar5iv | text | # Transverse Isotropy in Identical Particle Scattering
## Abstract
It is pointed out that the cross section for the scattering of identical charged bosons is isotropic over a broad angular range around 90<sup>o</sup> when the Sommerfeld parameter has a critical value, which depends exclusively on the spin of the particle. A discussion of systems where this phenomenon can be observed is presented.
PACS: 03.65.Nk
The scattering of identical particles is a routine exercise in quantum mechanics and its discussion can be found in most text books on the subject . However, in recent years, the rapid oscillation seen in the angular distribution of the Mott elastic scattering of identical charged particles such as nuclei, was utilized to test models such as QCD and to discuss small deviations from the Coulomb force law owing to QED-related corrections such as vacuum polarization .
In the present work we point out a hitherto unknown feature of the Mott scattering cross-section for bosons, namely an almost isotropic angular distribution over a very wide angular range when the Sommerfeld parameter attains a critical value determined entirely by the spin, $`s`$, of the particles viz, $`\eta _C=(3s+2)^{1/2}`$. For the purpose of completeness we first give a short account of the theory of scattering of identical particles. We then turn to the derivation of $`\eta _C`$ and apply to several boson-boson scattering systems. We also briefly discuss the fermion-fermion case.
If in a scattering process the projectile and target particles are identical, when one of them reaches a detector the experiment cannot tell if this particle is the projectile or the target. On the other hand, momentum conservation guarantees that whenever a particle emerges in one direction, the other emerges in the opposite orientation, in the CM frame of reference. Therefore, the amplitude for scattering at the orientations $`𝐫`$ and $`𝐫`$ will be mixed in some way. In Quantum mechanics, the total wave function for pairs of identical particles with integer spins, i.e. two bosons, must be symmetric with respect to the exchange of these particles, while in the case of particles with half-integer spin, i.e. two fermions, it must be anti-symmetric.
The situation is simple for spinless bosons or in collisions where the projectile and the target are polarized so that their spins are aligned. In such cases, the wave function in the spin space is always symmetric and projectile-target exchange reduces to reflection of the relative vector position $`𝐫`$. For spherically symmetric potentials, there is axial symmetry and space reflection corresponds to the transformation $`\theta \pi \theta `$. The elastic cross section is then given by
$$\sigma _\pm (\theta )=\left|f(\theta )\pm f(\pi \theta )\right|^2,$$
(1)
where $`f(\theta )`$ is the scattering amplitude for discernible particles of the same mass under the same potential $`V(r)`$. The $`+()`$ sign in eq.(1) applies when the particles involved are bosons (fermions).
Eq.(1) may be rewritten in the form
$$\sigma _\pm (\theta )=\sigma _{inc}(\theta )\pm \sigma _{int}(\theta ),$$
(2)
with
$$\sigma _{inc}(\theta )=\left|f(\theta )\right|^2+\left|f(\pi \theta )\right|^2$$
(3)
and
$$\sigma _{int}(\theta )=2\mathrm{R}\mathrm{e}\left\{f^{}(\theta )f(\pi \theta )\right\}.$$
(4)
The first term in eq.(2) is the incoherent sum of the contributions to the cross sections arising from projectile and target, if they were distinguishable. While this term is independent of the particle statistics, the sign of the second term is responsible for the difference in the expressions for the cross section of bosons and fermions. This interference term has no classical analogue.
The situation is more complicated in the case of unpolarized spins. In this case, the cross section mixes different parities as the spins couple to produce symmetric or anti-symmetric states in the spin space. However, taking the proper average over spin orientations one obtains the simple formula
$$\sigma _\pm (\theta )=\sigma _{inc}(\theta )\pm \frac{\sigma _{int}(\theta )}{2s+1},$$
(5)
where $`s`$ is the spin of the particle in units of $`\mathrm{}`$. Eq.(5) indicates that the relevance of the interference term decreases with the spin value, vanishing in the classical limit $`s\mathrm{}`$.
The cross sections of eq.(5) are symmetric with respect to $`\theta =90^o`$ and their particular shape depends on several factors such as the statistics of the colliding particles, their interaction and the bombarding energy. A particularly interesting situation is the Coulomb scattering of bosons, where $`\sigma _+`$ is known as the Mott cross section and denoted $`\sigma _{Mott}`$. In this case, we have the analytical expressions
$$\sigma _{inc}(\theta )=\frac{a^2}{4}\left[\frac{1}{\mathrm{sin}^4\left(\theta /2\right)}+\frac{1}{\mathrm{cos}^4\left(\theta /2\right)}\right]$$
(6)
and
$$\sigma _{int}(\theta )=\frac{a^2}{4}\left[\frac{2}{\mathrm{sin}^2\left(\theta /2\right)\mathrm{cos}^2\left(\theta /2\right)}\mathrm{cos}\left(2\eta \mathrm{ln}(\mathrm{tan}^1(\theta /2))\right)\right],$$
(7)
where $`\eta `$ is the Sommerfeld parameter,
$$\eta =\frac{q^2}{\mathrm{}v}$$
(8)
and $`a`$ is half the distance of closest approach in a head-on collision,
$$a=\frac{q^2}{2E}.$$
(9)
Above, $`q`$ is the charge of each of the two collision partners, $`E`$ is the bombarding energy in the center of mass (CM) reference frame and $`v`$ is the corresponding velocity of the relative motion. In the present case, $`\sigma _{inc}`$ exhibits a minimum at $`\theta =90^o`$, with the value $`\sigma _{inc}\left(\theta =90^o\right)=2a^2`$, with an energy-independent shape. On the other hand, $`\sigma _{int}`$ has always a maximum at this angle, with the same value $`2a^2`$. However, its shape depends on the collision energy through the Sommerfeld parameter $`\eta `$. The behavior of $`\sigma _{Mott}`$ in the vicinity of $`\theta =90^o`$ results from a competition between these two opposing trends. For small $`\eta `$ values, $`\sigma _{int}`$ is a slowly varying function of $`\theta `$. The shape of $`\sigma _{Mott}`$ is then dominated by that of $`\sigma _{inc}`$ and it presents a minimum at $`\theta =90^o`$. For large $`\eta `$, the opposite situation takes place and $`\sigma _{Mott}`$ has a maximum at $`\theta =90^o`$. An interesting situation occurs at the critical value of the Sommerfeld parameter, $`\eta _C`$, where the cross section goes through this transition. The value of $`\eta _C`$ is obtained from the condition
$$\left[\frac{d^2\sigma _{Mott}(\theta )}{d\theta ^2}\right]_{\theta =90^o}=0.$$
(10)
Using eqs.(6) and (7), we obtain
$`\left[{\displaystyle \frac{d^2\sigma _{Mott}(\theta )}{d\theta ^2}}\right]_{\theta =90^o}=16a^2\left[{\displaystyle \frac{12\eta ^2}{2s+1}}+3\right]`$
and according to eq.(10) we get
$$\eta _C=\sqrt{3s+2}.$$
(11)
In figure 1, we show cross sections normalized to the value of the Rutherford cross section at $`90^o,`$ $`\sigma _{Ruth}(90^o)=a^2`$, for collisions of identical bosons with spins $`s=0`$ and $`s=1`$. In each case, the calculations were performed at $`\eta _C`$. I.e., $`\eta =\sqrt{2}`$ for $`s=0`$ and $`\eta =\sqrt{5}`$ for $`s=1`$. Also shown for comparison is the incoherent cross section of eq.(6), normalized in the same way. Clearly, $`\sigma _{inc}(90^o)/a^2=2`$ and $`\sigma _b(90^o)/a^2=4`$ as shown in the figure. The striking feature of the figure is the flatness of $`\sigma _{Mott}`$ over a very wide angular region around $`\theta =90^o`$. It is essentially constant for $`60^o<\theta <120^o`$, for $`s=0`$, and $`70^o<\theta <110^o`$, for $`s=1`$. This ‘transverse isotropy’ (TI) is universal as the only relevant parameter which enters the discussion is the Sommerfeld parameter. In principle, it could be observed in atomic or nuclear systems at the appropriate energy. However, as we shall show below, the most appropriate case to investigate the above ‘transverse isotropy’ is that of low-energy scattering of light identical nuclei, such us d-d or $`\alpha \alpha `$.
Investigating a physical system which shows TI could shed light on several small effects related to QED and possibly to QCD, as well as to atomic effects in nuclear scattering. Also the assumed pure bosonic nature of the multifermionic cluster could be nicely examined by a careful analysis of data taken at $`\eta _C`$.
At this stage, it is important to investigate the optimal conditions for observation of TI. The effective forces between identical nuclei or ionized atoms are composed of the long range Coulomb part plus a shorter range nuclear or Van der Waals force. Since the above discussion was based on a pure Coulomb force, it is important to seek the physical conditions that allows the neglect of the short range forces. Calling $`E_C`$ the collision energy corresponding to $`\eta _C`$, the above condition corresponds to the requirement that $`E_C`$ be sufficiently below the Coulomb barrier $`V_B`$, the outermost maximum in the effective potential for $`l=0`$. In the collision of identical particles of charge $`q`$ and mass $`M`$, the critical collision energy $`E_C`$ is given by
$`E_C={\displaystyle \frac{Mq^4}{4\mathrm{}^2(3s+2)}}.`$
If one approximates $`V_B`$ by the coulomb potential at the barrier radius $`R_B`$, namely
$$V_B=\frac{q^2}{R_B},$$
(12)
the condition $`E_C<V_B`$ yields
$$\frac{Mq^2R_B}{\mathrm{}^2(3s+2)}<1.$$
(13)
Since the barrier radius in atomic collisions is very large, the above condition cannot be satisfied. We then consider nuclear collisions. In the nuclear physics one usually takes $`R_B=1.4`$ $`(M/m_0)^{1/3}`$fm, where $`m_0`$ is the nucleon mass. For light nuclei one can assume equal number of protons and neutrons and set $`M/m_0=2Z`$, where $`Z`$ is the atomic number. Eq.(13) then reduces to
$$Z^{10/3}<25.4\times (2s+1).$$
(14)
It can immediately be checked that the above condition is only satisfied for $`\alpha `$particles, in the case of $`s=0`$, and for $`dd`$ and $`{}_{}{}^{6}Li^6Li`$ collisions, in the case of $`s=1`$ (if one relaxes the condition of equal numbers of protons and neutrons, a couple of additional nuclei can be included in this set). The Mott cross section at $`\theta =90^o`$ at for collisions with the critical bombarding energy, $`E_C,`$ is then given by the simple expression
$$\sigma _{Mott}(\theta =90^o)=\frac{\sigma _0}{Z^6}\left(3s+2\right)^2,$$
(15)
with $`\sigma _0=33.7`$ barn.
In table I, we give the relevant quantities in the three above mentioned cases. It is clear from this table that the transverse isotropy is an observable phenomenon, although experimentally difficult in the case of $`dd`$. An important question to address now is the sensitivity of our result to the uncertainty of the energy of the beam. Since a slight change in $`\eta `$ form $`\eta _C`$ might wash out the TI, it is interesting to assess the range of $`\eta `$values around $`\eta _C`$ which can still tolerate a meaningfull study of the phenomenon. Figure 2 shows the way the shape of the Mott cross section changes as $`\eta `$ is varied from $`\eta _C`$ by $`5\%`$. This uncertainty in the $`\eta `$value corresponds to about 10$`\%`$ uncertainty in the collision energy, which is certainly attainable in existing accelerators. It is clear from figure 2 that the transverse isotropy should be visible within a few per cent energy resolutions.
Although we have restricted our discussion so far to $`E_C<V_B`$, we emphasize that the TI may come up at higher energies. However, in this case the details of the short-range interaction becomes important and the discussion becomes model- and system-dependent. We also point out that no transverse isotropy is expected for Coulomb collisions of identical fermions, since in this case both $`\sigma _{inc}`$ and $`\sigma _{int}`$ have minima at $`\theta =90^o`$. However, the situation may be different when the short range interaction dominates. We have briefly looked into this question by examining the extreme case of hard sphere scattering with no Coulomb interaction. The physical parameter that characterizes the collision process is $`kR`$ with $`R`$ being the sum of the radii of the two colliding particles. Our preliminary results indeed show a TI for identical fermions with the rather large spin values $`s9/2`$. The critical value of $`kR`$ is the order of 2.5, which is of the same order of magnitude as the value we found for bosons (1.5). More details of the present work with extensions to other potentials will be published elsewhere .
This work was supported in part by CNPq and the MCT/FINEP/CNPq(PRONEX) under contract no. 41.96.0886.00.
Tables
* Table I: The relevant quantities associated with the d-d, $`{}_{}{}^{6}Li^6Li`$ and $`\alpha \alpha `$ collisions, at the critical value of the Sommerfeld parameter.
| System | $`s`$ | $`E_C(`$keV) | $`V_B(`$keV) | $`\sigma _{Mott}(90^o)(`$barn) |
| --- | --- | --- | --- | --- |
| d + d | 1 | 5.0 | 400 | 135 |
| <sup>6</sup>Li $`+`$<sup>6</sup>Li | 1 | 1200 | 2500 | 1.17 |
| $`\alpha +\alpha `$ | 0 | 400 | 1260 | 2.3 |
Figure Captions
* Figure 1: The Mott cross sections for collisions of identical bosons at the critical value of the Sommerfeld parameter. Results are shown in the case of bosons with spin 0 (solid line) and spin 1 (dashed line), to which corresponds respectively $`\eta _C=\sqrt{2}`$ and $`\eta _C=\sqrt{5}`$. For comparison, the incoherent part of the cross section is also shown (dot-dashed line).
* Figure 2: Sensitivity of the transverse isotropy as $`\eta `$ deviates from $`\eta _C`$ by 5 %. The results are shown in the case of $`s=0`$. |
warning/0003/hep-th0003159.html | ar5iv | text | # 1 Introduction
## 1 Introduction
By now it has been well-established that the smallest eigenvalues of the QCD Dirac operator are correlated according to a Random Matrix Theory with the global symmetries of the QCD partition function (see for recent reviews and a complete list of references). In particular, this has been confirmed by the analysis of the low-energy effective theory , universality studies , lattice QCD simulations and the study two-sublattice theories with disorder . This means that the dynamical details of QCD are not important on energy scales of the order of the average level spacing. The natural question that can be asked is at what energy scale the dynamics of QCD becomes relevant and how does this manifest itself in the Dirac spectrum.
The answer to this question has been understood within the context of effective theories . The effective theory for the QCD Dirac spectrum is known as the partially quenched effective partition function and was originally introduced to study the quenched approximation in QCD . The central observation is that in the domain where the mass dependence of the effective partition function is given by the contribution of the constant fields (the zero momentum modes) the Dirac eigenvalues are correlated according to chiral Random Matrix Theory. The relevant mass scale can thus be identified as the scale for which the Compton wavelength of the lightest particle becomes equal to the size of the box i.e., $`M1/L_s`$. In QCD, in the phase of spontaneously broken chiral symmetry, the lightest particles are the Goldstone modes with a mass given by $`M^2Km`$ (with $`m`$ the quark mass and $`K=\mathrm{\Sigma }/F^2`$ in terms of the pion decay constant $`F`$ and the chiral condensate $`\mathrm{\Sigma }`$). The critical scale is thus given by
$`E_c={\displaystyle \frac{1}{KL_s^2}}.`$ (1)
In the context of disordered condensed matter systems this energy scale is known as the Thouless energy and also in this article we will adopt this name.
A more intuitive interpretation of the Thouless energy has been given in the theory of mesoscopic systems . The time scale $`\mathrm{}/E_c`$ is the time for which an initially localized wave packet diffuses all over space. For this reason the eigenvalues are correlated according to Random Matrix Theory for energy differences below $`E_c`$ (known as the ergodic regime). At shorter time scales, different wave functions do not necessarily overlap resulting in a weakening of correlations of the corresponding eigenvalues. For energy differences beyond the inverse elastic collision time $`\tau _e`$ the corresponding eigenvalues are completely uncorrelated (the Poisson ensemble). The domain inbetween $`E_c`$ and $`\mathrm{}/\tau _e`$ is known as the diffusive or Altshuler-Shklovskii domain. A third energy scale is the average level spacing $`\mathrm{\Delta }\lambda `$. The ratio $`E_c/\mathrm{\Delta }\lambda `$ is identified in mesoscopic physics as the dimensionless conductance. It is equal to the number of subsequent levels correlated according to Random Matrix Theory. The existence of these domains has been confirmed by numerical simulations of the Anderson model .
In QCD the average level spacing is related to the order parameter of the chiral phase transition, the chiral condensate, by the Banks-Casher formula according to $`\mathrm{\Delta }\lambda =\pi /\mathrm{\Sigma }V`$ (with $`V`$ the volume of Euclidean space-time). The prediction from (1) is that the number of eigenvalues correlated according to chRMT is of the order $`F^2\sqrt{V}`$.
For increasing disorder the number of subsequent eigenvalues described by Random Matrix Theory decreases. For strong disorder we expect that all states become localized with uncorrelated eigenvalues. We thus expect a critical value of the disorder for which the three scales, $`\mathrm{\Delta }\lambda `$, $`E_c`$ and $`\mathrm{}/\tau _e`$ coincide. In particular, the dimensionless conductance becomes volume independent . It has been conjectured that at this point the eigenvalue correlations are described by a new universality class known as critical statistics (see for a review). In this class, only the short range correlations of the eigenvalues are described by the usual random matrix ensembles whereas the number variance, $`\mathrm{\Sigma }^2\left(n\right)`$, shows a linear $`n`$-dependence beyond this domain. What is relevant for QCD is that such behavior has been observed in numerical simulations of the 4-dimensional Anderson model .
The volume dependence of the Thouless energy has been investigated by means of lattice QCD simulations and instanton-liquid simulations . In essence, results from lattice QCD simulations are in complete agreement with theoretical results from partial quenched chiral perturbation theory. However, the results from instanton liquid simulations seem to deviate from the prediction (1) with a scale independent constant $`K`$; in that case the Thouless energy only shows a weak volume dependence. This raises the question whether the Dirac eigenvalues might be described by critical statistics. To address this issue we generalize a random matrix model for critical statistics to include the chiral symmetry of the QCD partition function (section 2). In section 3 we map our model on a partition function of noninteracting fermions. This model is solved in the semi-classical limit in section 4, where we obtain analytical expressions for the microscopic spectral density and the two-point correlation function. Comparisons with instanton simulations are shown in section 5 and concluding remarks are made in section 6.
## 2 Definition of the Model
The random matrix model of Moshe and Neuberger and Shapiro is defined by the partition function
$`Z={\displaystyle 𝑑He^{\mathrm{Tr}HH^{}}𝑑Ue^{b\mathrm{Tr}\left([U,H][U,H]^{}\right)}},`$ (2)
where $`H`$ is a Hermitian and $`U`$ a Unitary $`N\times N`$ matrix. The integration measures $`dH`$ and $`dU`$ are given by the Haar measure. This model can be interpreted as the zero-dimensional limit of the Kazakov-Migdal model . It interpolates between the Gaussian Unitary Ensemble $`(b=0`$) and the Poisson Ensemble $`\left(b\mathrm{}\right)`$. Using the invariance of the measure, the integral over $`U`$ can be replaced by an integral over the eigenvalues of $`U`$. For $`b\mathrm{}`$ this partition function is dominated by matrices $`H`$ that commute with arbitrary diagonal unitary matrices. This set of matrices is the ensemble of diagonal Hermitian matrices which is known as the Poisson ensemble. What is nice about this model is that it preserves the unitary invariance which enables us to take full advantage of the existing random matrix theory methods. In order to obtain a nontrivial $`b`$-dependence in the thermodynamic limit, the parameter $`b`$ has to be scaled as
$`b=h^2N^2.`$ (3)
It has been shown that in this limit the model (2) is equivalent to both a banded random matrix model with a power-like cutoff and to random matrix models with a $`\mathrm{Tr}\mathrm{log}^2H`$ probability potential . The correlation functions of the latter model have been derived by means of $`q`$-orthogonal polynomials and Painlevé equations .
In this paper we are interested in chiral random matrix ensembles defined as ensembles of $`N\times N`$ random matrices with the structure
$`D=\left(\begin{array}{cc}0& C\\ C^{}& 0\end{array}\right),`$ (6)
where $`C`$ is an arbitrary complex matrix $`n\times \left(n+\nu \right)`$ matrix ($`N=2n+\nu `$). Since the matrix $`D`$ has exactly $`\nu `$ zero eigenvalues $`\nu `$ is interpreted as the topological quantum number. The chiral Gaussian Unitary Ensemble (chGUE) with $`N_f`$ flavors is defined as the ensemble of matrices $`D`$ with matrix elements $`C`$ distributed according to the Gaussian probability distribution
$`P\left(C\right)\stackrel{N_f}{det}\left(D+m\right)e^{N\mathrm{Tr}CC^{}}.`$ (7)
Here, for simplicity we have taken all quark masses equal to $`m`$. The probability distribution $`P\left(C\right)`$ has the unitary invariance
$`CUCV^1,`$ (8)
where $`U`$ and $`V`$ are unitary matrices. Since an arbitrary complex matrix can always be brought to diagonal form by this transformation this invariance allows us to factorize the probability distribution in a product over the eigenvalues of $`C`$ and the unitary matrices that diagonalize $`C`$.
The generalization of the model of Moshe, Neuberger and Shapiro to the chiral ensembles is immediate. The interpolating model is defined by the partition function
$`Z_\nu ={\displaystyle 𝑑C\stackrel{N_f}{det}\left(D+m\right)e^{\frac{\mathrm{\Sigma }^2}{4h}\mathrm{Tr}D^{}D}𝑑𝒰e^{\frac{\mathrm{\Sigma }^2hN^2}{4}\mathrm{Tr}[D,𝒰][D,𝒰]^{}}},`$ (9)
where $`𝒰`$ has the chiral block structure
$`𝒰=\left(\begin{array}{cc}U& 0\\ 0& V\end{array}\right),`$ (12)
with $`U`$ an $`n\times n`$ unitary matrix and $`V`$ an $`\left(n+\nu \right)\times \left(n+\nu \right)`$ unitary matrix. If $`\mathrm{\Sigma }`$ and $`h`$ are constants as in (9) we will refer to this model as the critical chiral unitary ensemble. As we will see below, in order to make contact with the Thouless energy in the partially quenched effective partition function we have to scale $`h`$ with an additional factor $`1/\sqrt{N}`$. The unitary invariance of this partition function follows from the invariance of the Haar measure. In comparison to , an additional factor $`\mathrm{\Sigma }^2/h`$ has been included in the probability distribution of the matrix elements. As we will see below, this will guarantee that in the thermodynamic limit the spectral density is $`h`$-independent to leading order in $`h`$. We will also find that the partition function is normalized such that the $`\mathrm{\Sigma }`$ represents the chiral condensate by means of the Banks-Casher relation $`\mathrm{\Sigma }=\pi \rho \left(0\right)/N`$ (with $`\rho \left(0\right)`$ the spectral density around $`\lambda =0`$).
Decomposing into the blocks of $`D`$ and $`𝒰`$ the partition function can be written as
$`Z_\nu =m^\nu {\displaystyle 𝑑C\stackrel{N_f}{det}\left[C^{}C+m^2\right]e^{\left(1+2h^2N^2\right)\frac{\mathrm{\Sigma }^2}{2h}\mathrm{Tr}CC^{}}𝑑U𝑑Ve^{\mathrm{\Sigma }^2hN^2\mathrm{ReTr}UCV^1C^{}}}.`$ (13)
The arbitrary complex matrix $`C`$ can be decomposed according to
$`C=U_1\mathrm{\Lambda }U_2,`$ (14)
and the integral over $`C`$ can be expressed as an integral over the eigenvalues $`\mathrm{\Lambda }_k`$ and the unitary matrices $`U_1`$ and $`U_2`$. Up to a irrelevant constant, the Jacobian of this transformation is given by
$`J\left(\mathrm{\Lambda }\right)=\mathrm{\Delta }^2\left(\left\{\lambda _i^2\right\}\right){\displaystyle \underset{k}{}}\lambda _k^{2\nu +1}.`$ (15)
where
$`\mathrm{\Delta }\left(\left\{\lambda _i^2\right\}\right)={\displaystyle \underset{k<l}{}}\left(\lambda _k^2\lambda _l^2\right)`$ (16)
is the Vandermonde determinant. Because of the unitary invariance, the $`U_1`$ and $`U_2`$ dependence in the second exponent can be absorbed in a redefinition of $`U`$and $`V`$, and the integrations over $`U_1`$ and $`U_2`$ just result in an overall constant. Remarkably, the integral over $`U`$ and $`V`$ in (13) is an Itzykson-Zuber type integral which is known analytically
$`{\displaystyle 𝑑U𝑑Ve^{z\mathrm{ReTr}U\mathrm{\Lambda }V^1\mathrm{\Lambda }}}=C{\displaystyle \frac{det_{k,l}\left|I_\nu \left(z\lambda _k\lambda _l\right)\right|}{\mathrm{\Delta }^2\left(\left\{\lambda _i^2\right\}\right)_{k=1}^n\lambda _k^{2\nu }}}.`$ (17)
The Vandermonde determinants cancel resulting in the partition function
$`Z_\nu =m^\nu {\displaystyle \underset{k=1}{\overset{n}{}}d\lambda _k\underset{k=1}{\overset{n}{}}\left(\lambda _k^2+m^2\right)^{N_f}\underset{k,l}{det}\left|\left(\lambda _k\lambda _l\right)^{\frac{1}{2}}e^{\left(\frac{1}{2}+h^2N^2\right)\frac{\mathrm{\Sigma }^2}{2h}\left(\lambda _k^2+\lambda _l^2\right)}I_\nu \left(h\mathrm{\Sigma }^2N^2\lambda _k\lambda _l\right)\right|}.`$
(18)
The instanton-liquid Dirac spectra that will be described by this model were obtained in the quenched approximation and for zero total topological charge. Therefore we will not attempt to solve this model for arbitrary $`N_f`$ but instead focus on the technically simpler case of $`N_f=0`$. Since the topological charge does not give rise to additional complications we will consider the case of arbitrary $`\nu `$. Below we will thus analyze the joint probability distribution
$`\rho _\nu (\lambda _1,\mathrm{},\lambda _n)\underset{k,l}{det}\left|\left(\lambda _k\lambda _l\right)^{\frac{1}{2}}e^{\left(\frac{1}{2}+h^2N^2\right)\frac{\mathrm{\Sigma }^2}{2h}\left(\lambda _k^2+\lambda _l^2\right)}I_\nu \left(h\mathrm{\Sigma }^2N^2\lambda _k\lambda _l\right)\right|.`$
(19)
## 3 Fermi-Gas Representation of the Interpolating chiral Random Matrix Model
In this section we rewrite the joint probability distribution (19) in terms of the fermionic $`n`$-particle matrix element
$`\rho _\nu (x_1,\mathrm{},x_n)Cx_1\mathrm{}x_n\left|e^{\beta H}\right|x_1\mathrm{}x_n,`$ (20)
where $`H`$ is the separable Hamiltonian
$`H={\displaystyle \underset{k}{}}\left(_k^2+{\displaystyle \frac{4\nu ^21}{4x_k^2}}+\omega ^2x_k^2\right),`$ (21)
and $`C`$ is an irrelevant constant. The eigenfunctions of the single particle Hamiltonian are known in terms of Laguerre polynomials. Specifically,
$`\left(_x^2+{\displaystyle \frac{4\nu ^21}{4x^2}}+\omega ^2x^2\right)\varphi _n=\omega _n\varphi _n,`$ (22)
has the solutions
$`\omega _n`$ $`=`$ $`\left(4n+2\nu +2\right)\omega ,`$ (23)
$`\varphi _n`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\omega ^{1/4}}{\sqrt{h}_n}}e^{\omega x^2/2}\left(x\sqrt{\omega }\right)^{\nu +\frac{1}{2}}L_n^\nu \left(x^2\omega \right).`$ (24)
The normalization factor $`h_n=\left(n+\nu \right)!/n!`$ ensures that the $`\varphi _n`$ are normalized to unity
$`{\displaystyle _0^{\mathrm{}}}\varphi _n\left(x\right)\varphi _n\left(x\right)𝑑x=1.`$ (25)
Using completeness, the many-particle matrix element can be written as
$`x_1\mathrm{}x_n\left|e^{\beta H}\right|x_1\mathrm{}x_n=\underset{ij}{det}{\displaystyle \underset{n}{}}\varphi _n\left(x_i\right)e^{\beta \omega _n}\varphi _n\left(x_j\right)`$ (26)
The sum over $`n`$ can be performed analytically using the identity
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{L_n^\nu \left(x\right)L_n^\nu \left(y\right)z^n}{h_n}}={\displaystyle \frac{\left(xyz\right)^{\nu /2}}{1z}}\mathrm{exp}\left(z{\displaystyle \frac{x+y}{1z}}\right)I_\nu \left({\displaystyle \frac{2\sqrt{xyz}}{1z}}\right).`$ (27)
With $`z=\mathrm{exp}\left(4\beta \omega \right)`$ this results in
$`x_1\mathrm{}x_n\left|e^{\beta H}\right|x_1\mathrm{}x_n=\underset{ij}{det}\left|{\displaystyle \frac{\omega \sqrt{x_ix_j}}{\mathrm{sinh}2\beta \omega }}\mathrm{exp}\left({\displaystyle \frac{\omega \mathrm{cosh}2\beta \omega }{2\mathrm{sinh}2\beta \omega }}\left(x_i^2+x_j^2\right)\right)I_\nu \left({\displaystyle \frac{\omega }{\mathrm{sinh}2\beta \omega }}x_ix_j\right)\right|.`$
Comparing this expression with (19) we find that
$`\omega `$ $`=`$ $`\sqrt{4h^2N^2+1}{\displaystyle \frac{\mathrm{\Sigma }^2}{2h}},`$ (29)
$`\mathrm{cosh}2\beta \omega `$ $`=`$ $`1+{\displaystyle \frac{1}{2h^2N^2}}.`$ (30)
The joint probability distribution of the eigenvalues is thus given by an $`n`$particle diagonal matrix element of the density operator. The average spectral density is equal to the average particle density. It is obtained by integrating over the positions of all particles except one. The integral can be performed by rewriting the matrix elements (26) in terms of a sum over permutations $`\pi `$ and $`\pi ^{}`$,
$`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{\omega _1<\mathrm{}<\omega _n}{}}{\displaystyle \underset{\pi \pi ^{}}{}}\sigma \left(\pi \right)\sigma \left(\pi ^{}\right)\varphi _{\omega _{\pi \left(1\right)}}\left(x_1\right)\mathrm{}\varphi _{\omega _{\pi \left(n\right)}}\left(x_n\right)e^{\beta _k\omega _k}\varphi _{\omega _{\pi ^{}\left(1\right)}}\left(x_1\right)\mathrm{}\varphi _{\omega _{\pi ^{}\left(n\right)}}\left(x_n\right),`$ (31)
where $`\sigma ()`$ is the sign of the permutation. Performing the integrations over $`x_2,\mathrm{},x_n`$, by orthogonality we find that the only nonzero contribution is for $`\pi =\pi ^{}`$. Then the remaining sum over $`\pi \left(\omega _1\right)`$ is just a sum over $`\omega _1,\mathrm{},\omega _n`$. For the canonical ensemble we thus find the one-particle density,
$`\rho \left(x\right){\displaystyle \underset{k}{}}\delta \left(xx_k\right)={\displaystyle \frac{1}{Z_n}}{\displaystyle \underset{\omega _1<\mathrm{}<\omega _N}{}}{\displaystyle \underset{i}{}}\left|\varphi _{\omega _i}\left(x_1\right)\right|^2e^{\beta _k\omega _k}.`$ (32)
In an occupation number representation this can be rewritten as
$`\rho \left(x\right)={\displaystyle \frac{1}{Z_n}}{\displaystyle \underset{_in_i=n}{}}{\displaystyle \underset{i}{}}n_i\left|\varphi _{\omega _i}\left(x\right)\right|^2e^{\beta _kn_i\omega _k},`$ (33)
where the sum is over all $`n_i\{0,\mathrm{\hspace{0.17em}1}\}`$ subject to the condition given below the summation sign and the canonical partition function is given by
$`Z_n={\displaystyle \underset{_in_i=n}{}}e^{\beta _in_i\omega _i}.`$ (34)
For completeness we give the following exact expressions for the canonical partition function
$`Z_n={\displaystyle \frac{e^{2\beta \omega n^22\beta \nu \omega n}}{\left(1e^{4\beta \omega }\right)\left(1e^{8\beta \omega }\right)\mathrm{}\left(1e^{4n\beta \omega }\right)}},`$ (35)
and the one-particle density
$`\rho \left(x\right)={\displaystyle \underset{k}{}}\left|\varphi _k\left(x\right)\right|^2{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left(1\right)^re^{4\beta \omega \left(krnr+\frac{1}{2}r\left(r+1\right)\right)}{\displaystyle \underset{p=Nr}{\overset{N}{}}}\left(1e^{4\beta \omega p}\right).`$ (36)
They have been obtained by writing the constraint $`_in_i=n`$ as
$`\delta \left(N{\displaystyle \underset{i}{}}n_i\right)={\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑\theta e^{i\theta \left(N_in_i\right)}},`$ (37)
and summing the geometric series after using the explicit expression (24) for the $`\omega _n`$.
Instead of working with the canonical ensemble we eliminate the constraint on the sum over the $`n_i`$ by working with the grand canonical ensemble. The grand canonical partition function is defined by
$`Z={\displaystyle \underset{n}{}}z^nZ_n={\displaystyle \underset{k}{}}\left(1+e^{\beta \omega _k+\beta \mu }\right),`$ (38)
where $`z=e^{\beta \mu }`$ is the fugacity. The one particle density in the grand canonical ensemble is given by
$`\rho \left(x\right)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{n}{}}z^nZ_n\rho \left(x\right)`$ (39)
$`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{n_i\{0,1\}}{}}{\displaystyle \underset{i}{}}n_i\left|\varphi _i\left(x\right)\right|^2e^{\beta _kn_k\left(\omega _k\mu \right)},`$
which can be evaluated to be
$`\rho \left(x\right)={\displaystyle \underset{i}{}}\left|\varphi _i\left(x\right)\right|^2{\displaystyle \frac{1}{1+e^{\beta \left(\omega _i\mu \right)}}}.`$ (40)
The fugacity $`z=e^{\beta \mu }`$ is determined by the condition that the total number of particles is equal to $`n`$, i.e.
$`{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{1+e^{\beta \left(\omega _i\mu \right)}}}=n.`$ (41)
The connected two-particle correlation function is defined by
$`R_2(x,y)={\displaystyle \underset{kl}{}}\delta \left(xx_k\right)\delta \left(yy_l\right)\rho \left(x\right)\rho \left(y\right)`$ (42)
It can be obtained from the $`n`$-particle matrix element (20) by integration over all coordinates (eigenvalues) except $`x_1`$ and $`x_2`$. The remaining sum over $`\omega _{\pi \left(1\right)}`$ and $`\omega _{\pi \left(2\right)}`$ can be rewritten as
$`{\displaystyle \underset{\omega _1<\mathrm{}<\omega _N}{}}{\displaystyle \underset{ij}{}}\varphi _{\omega _i}\left(x_1\right)\varphi _{\omega _j}\left(x_2\right)\left[\varphi _{\omega _i}\left(x_1\right)\varphi _{\omega _j}\left(x_2\right)\varphi _{\omega _j}\left(x_1\right)\varphi _{\omega _i}\left(x_2\right)\right]e^{\beta _k\omega _k}.`$ (43)
The diagonal term, $`i=j`$, is zero allowing us to include it in the summation. The first term can then be identified as the square of the average particle density. It is the disconnected contribution to the two-particle correlation function. The connected part of the two-particle distribution function can thus be written as
$`R_2(x_1,x_2)={\displaystyle \frac{1}{Z_n}}{\displaystyle \underset{\omega _1<\mathrm{}<\omega _n}{}}{\displaystyle \underset{i,j}{}}\varphi _{\omega _i}\left(x_1\right)\varphi _{\omega _j}\left(x_2\right)\varphi _{\omega _j}\left(x_1\right)\varphi _{\omega _i}\left(x_2\right)e^{\beta _i\omega _i}.`$ (44)
This correlation function does not include the contributions from the self-correlations of the eigenvalues. After all, our starting point was the joint distribution of different eigenvalues. In an occupation number representation $`R_2(x_1,x_2)`$ simplifies to
$`R_2(x_1,x_2)={\displaystyle \frac{1}{Z_n}}{\displaystyle \underset{n_1+n_2+\mathrm{}=n}{}}\left[{\displaystyle \underset{i}{}}n_i\varphi _{\omega _i}\left(x_1\right)\varphi _{\omega _i}\left(x_2\right)\right]^2e^{\beta _kn_k\omega _l}.`$ (45)
Using this representation, the two-particle density in the grand canonical ensemble is found to be
$`R_2(x,y)=\left|{\displaystyle \underset{i}{}}\varphi _i\left(x\right)\varphi _i\left(y\right){\displaystyle \frac{1}{1+e^{\beta \left(\omega _i\mu \right)}}}\right|^2.`$ (46)
Using similar manipulations as for the canonical partition function and the corresponding one-particle density it is possible to simplify the exact analytical expression for the connected two-particle correlation function in the canonical ensemble,
$`R_2(x,y)={\displaystyle \underset{k,l}{}}\left|\varphi _k\left(x\right)\varphi _l\left(y\right)\right|^2{\displaystyle \underset{r,s=0}{\overset{\mathrm{}}{}}}\left(1\right)^re^{4\beta \omega \left(kr+lsn\left(r+s\right)+\frac{1}{2}\left(r+s\right)\left(r+s+1\right)\right)}{\displaystyle \underset{p=Nrs}{\overset{N}{}}}\left(1q^p\right).`$
This result can be used to compare the two ensembles but we will not address this question in this article.
## 4 Semiclassical Calculation
In this section we calculate the microscopic spectral density and the two-point correlation function using semi-classical methods starting from the expressions for the grand canonical ensemble derived in previous section. We are thus interested in the region around $`x=0`$. Because of the hard edge at $`x=0`$ we cannot simply do a WKB approximation by replacing the wave functions by plane waves but instead have to use Bessel functions. This follows immediately from the wave equation for $`x0`$,
$`\left[_x^2+{\displaystyle \frac{4\nu ^21}{4x^2}}\right]\sqrt{x}J_\nu \left(kx\right)=k^2\sqrt{x}J_\nu \left(kx\right).`$ (48)
Alternatively, one can exploit the asymptotic relation between Laguerre polynomials and Bessel functions.
### 4.1 The Microscopic Spectral Density
Taking into account the normalization of the Bessel functions, $`_0^{\mathrm{}}k𝑑k\sqrt{xx^{}}J_\nu \left(kx\right)J_\nu \left(kx^{}\right)=\delta \left(xx^{}\right)`$, to fix the constants in the integration measure, we arrive at the following expression for the single particle density
$`\rho \left(x\right)={\displaystyle _0^{\mathrm{}}}kx𝑑k{\displaystyle \frac{J_\nu ^2\left(kx\right)}{1+e^{\beta k^2\beta \mu +\beta \omega ^2x^2}}}.`$ (49)
The chemical potential is determined by the condition $`𝑑x\rho \left(x\right)=n`$. Since this integral is over all $`x`$, the use of the semi-classical expressions for the wave functions is not justified but instead we have to rely on the exact wave functions. In the limit, $`\beta \mu \omega `$, the sum over $`i`$ in (41) can be replaced be an integral which can be performed analytically resulting in
$`e^{\left(2\nu +2\right)\beta \omega \beta \mu }={\displaystyle \frac{1}{e^{4n\beta \omega }1}}.`$ (50)
In the limit $`n\beta \omega 1`$ the semi-classical expression for the spectral density is thus given by
$`\rho \left(x\right)={\displaystyle _0^{\mathrm{}}}kx𝑑k{\displaystyle \frac{J_\nu ^2\left(kx\right)}{1+e^{\beta k^2+\beta \omega ^2x^2\left(4n2\nu 2\right)\beta \omega }}}.`$ (51)
In this limit the semiclassical expression (49) leads to the correct value of the chemical potential. Below we will show for $`n\beta \omega 1`$ and finite $`x`$ (in units of the average level spacing) the term $`\beta \omega ^2x^2`$ can be neglected relative to $`4n\beta \omega `$.
An estimate for the average spectral density near zero but many level spacings away from $`x=0`$ is obtained by using the leading order asymptotic expansion of the Bessel functions and calculating the integral in (51) in the limit of a degenerate Fermi-gas. This results in
$`\overline{\rho }={\displaystyle \frac{\overline{k}}{\pi }},`$ (52)
where $`\overline{k}\sqrt{4n\omega }`$ is the “radius”of the Fermi-sphere. At fixed $`h`$ in the limit $`n1`$ we have
$`\omega \mathrm{\Sigma }^2n\mathrm{and}\beta \omega {\displaystyle \frac{1}{2hn}}.`$ (53)
Using these results and invoking the Banks-Casher formula the parameter $`\mathrm{\Sigma }`$ can be identified as the chiral condensate,
$`\underset{n\mathrm{}}{lim}{\displaystyle \frac{\pi \overline{\rho }}{2n}}=\mathrm{\Sigma }.`$ (54)
We will now show that our approximations are self-consistent. The condition that we are close to the degenerate Fermi gas can be written as $`\delta k/\overline{k}1`$. We thus have to impose the requirement that
$`{\displaystyle \frac{\delta k}{\overline{k}}}{\displaystyle \frac{1}{2\beta \overline{k}^2}}={\displaystyle \frac{h}{4}}1.`$ (55)
The conditions $`kx1`$ and $`\beta \omega ^2x^24n\beta \omega `$ can be combined into
$`{\displaystyle \frac{1}{4n}}\omega x^24n,`$ (56)
or, in units of the average level spacing, $`x=u/\overline{\rho }`$, the range of validity of the above asymptotic results is given by
$`{\displaystyle \frac{1}{4n}}{\displaystyle \frac{u^2\pi ^2}{4n}}4n.`$ (57)
Because of the second inequality it is justified to neglect the term $`\beta \omega ^2x^2`$ which we will do in the remainder of this section.
By partial integration the expression (51) for the spectral density can be rewritten as
$`\rho \left(x\right)={\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle \frac{\beta xk^3\left[J_\nu ^2\left(kx\right)J_{\nu +1}\left(kx\right)J_{\nu 1}\left(kx\right)\right]}{4\mathrm{cosh}^2\frac{\beta }{2}\left(k^24n\omega \right)}}.`$ (58)
Using the Banks-Casher formula one finds that the chiral condensate depends on $`h`$. The leading order correction is given by
$`{\displaystyle \frac{\mathrm{\Sigma }\left(h\right)}{\mathrm{\Sigma }}}=1{\displaystyle \frac{\pi ^2}{96}}h^2+\mathrm{}.`$ (59)
The corresponding spectral density will be denoted by $`\overline{\rho }\left(h\right)2n\mathrm{\Sigma }\left(h\right)/\pi `$. The microscopic spectral density is then given by
$`\rho _s\left(u\right)`$ $`=`$ $`\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{2n\mathrm{\Sigma }\left(h\right)}}\rho \left({\displaystyle \frac{u}{2n\mathrm{\Sigma }\left(h\right)}}\right),`$ (60)
$`=`$ $`{\displaystyle \frac{\mathrm{\Sigma }}{\mathrm{\Sigma }\left(h\right)}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle \frac{k^2\rho _s^0\left(ku\mathrm{\Sigma }/\mathrm{\Sigma }\left(h\right)\right)}{h\mathrm{cosh}^2\frac{\left(k^21\right)}{h}}}.`$
where $`\rho _s^0`$ is the microscopic spectral density for the chGUE
$`\rho _s^0\left(u\right)={\displaystyle \frac{u}{2}}\left[J_\nu ^2\left(u\right)J_{\nu +1}\left(u\right)J_{\nu 1}\left(u\right)\right].`$ (61)
The interpretation is clear. The oscillations in the microscopic spectral density due to the Bessel functions are smeared out over a distance $`hu`$ by the integration over $`k`$. The oscillations are thus visible up to a distance of $`u1/h`$.
In the limit of small $`h`$ the main contribution to the integral comes from the region around the Fermi-surface. In this limit we can derive an approximate formula correct to order $`h^2`$ at fixed $`uh`$. To this end we change integration variables in (60) according to $`k=1+ht`$ neglecting terms that are sub-leading in $`h`$. This results in
$`\rho _s\left(u\right)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle \frac{\rho _s^0\left(u+hut\right)}{\mathrm{cosh}^2\left(2t\right)}}.`$ (62)
Next we Taylor expand the microscopic spectral density as follows
$`\rho _s^0\left(u+hut\right)`$ $`=`$ $`\rho _s^0\left(u\right)`$ (63)
$`+`$ $`{\displaystyle \frac{1}{2}}hut[J_\nu ^2+J_{\nu +1}J_{\nu 1})]`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\left(hut\right)^2}{2!}}\left[4J_\nu J_{\nu 1}+{\displaystyle \frac{f_1}{u}}\right]`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\left(hut\right)^3}{3!}}[4(J_{\nu 1}^2J_{\nu +1}^2)+{\displaystyle \frac{g_1}{u}}+{\displaystyle \frac{g_2}{u^2}})]`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\left(hut\right)^4}{4!}}[4^2J_\nu J_{\nu 1}+{\displaystyle \frac{h_1}{u^1}}+{\displaystyle \frac{h_2}{u^2}}+{\displaystyle \frac{h_3}{u^3}})],`$
where the terms that will be neglected are denoted by $`f_i`$, $`g_i`$, $`h_i`$, etc.. At small values of $`u`$ these terms are of order $`h^2u`$ whereas for large $`u`$ they are suppressed by order $`1/u^2`$ (notice the factor $`u`$ in (61)). By inspection one easily finds that the neglected terms are at most of order $`h^2`$ independent of the value of $`u`$. The leading order terms can be easily resummed to
$`\rho _s^0\left(u+hut\right)=\rho _s^0\left(u\right){\displaystyle \frac{1}{2}}J_\nu \left(u\right)J_{\nu 1}\left(u\right)\left[\mathrm{cos}\left(2hut\right)1\right]+O\left(h^2\right)+\mathrm{terms}\mathrm{odd}\mathrm{in}t.`$ (64)
The integral over $`t`$ in (62) can be performed analytically resulting in
$`\rho _s\left(u\right)=\rho _s^0\left(u\right){\displaystyle \frac{1}{2}}J_\nu \left(u\right)J_{\nu 1}\left(u\right)\left[{\displaystyle \frac{\pi hu}{2\mathrm{sinh}\left(\pi hu/2\right)}}1\right]+O\left(h^2\right).`$ (65)
In Fig. 1 we compare the exact expression for the microscopic spectral density (60) to this approximate formula. We observe that even for a value of $`h`$ as large as $`h=0.3`$ the two results are very close.
### 4.2 Two-Point Function
In this subsection we derive a semi-classical expression for the two-point correlation function. To this end the wave functions in the expression (46) for the connected two-point correlation function are replaced by Bessel functions,
$`R_2(x,y)=\left|{\displaystyle _0^{\mathrm{}}}k𝑑k{\displaystyle \frac{\sqrt{xy}J_\nu \left(kx\right)J_\nu \left(ky\right)}{1+e^{\beta k^24n\beta \omega }}}\right|^2.`$ (66)
By partial integration with respect to $`k`$ the correlation function can be expressed as
$`R_2(x,y)=\left|{\displaystyle \frac{\overline{\rho }}{h}}{\displaystyle _0^{\mathrm{}}}k^2𝑑k{\displaystyle \frac{K(x\pi \overline{\rho }k,y\pi \overline{\rho }k)}{\mathrm{cosh}^2\left(\left(k^21\right)/h\right)}}\right|^2,`$ (67)
where $`K(x,y)`$ is the two-point kernel for the chGUE given by
$`K(x,y)=\sqrt{xy}{\displaystyle \frac{xJ_{\nu +1}\left(x\right)J_\nu \left(y\right)yJ_\nu \left(x\right)J_{\nu +1}\left(y\right)}{x^2y^2}}.`$ (68)
To order $`h^2`$ the correlation function can be simplified to
$`R_2(x,y)=\left|\overline{\rho }{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle \frac{K(x\pi \overline{\rho }(1+ht),y\pi \overline{\rho }(1+ht)}{\mathrm{cosh}^22t}}\right|^2+O\left(h^2\right),`$ (69)
In the same way as for the one-point function we now will derive an approximate formula for the two-point function correct to order $`h^2`$ at fixed value of $`uh`$. This can be done conveniently by using the following summation formula
$`J_{\nu +1}\left(x+xh\right)J_\nu \left(y+yh\right)`$ $`\pm `$ $`J_\nu \left(x+xh\right)J_{\nu +1}\left(y+yh\right)`$ (70)
$``$ $`\left[J_{\nu +1}\left(x\right)J_\nu \left(y\right)\pm J_\nu \left(x\right)J_{\nu +1}\left(y\right)\right]\mathrm{cos}\left[\left(x+y\right)h\right]`$
$`+`$ $`\left[J_\nu \left(x\right)J_\nu \left(y\right)J_{\nu +1}\left(x\right)J_{\nu +1}\left(y\right)\right]\mathrm{sin}\left[\left(x+y\right)h\right].`$
This formula has been derived by means of a Taylor expansion and a subsequent resummation employing the following approximate derivative formulas
$`_k^{2n}[J_{\nu +1}\left(kx\right)J_\nu \left(ky\right)`$ $`\pm `$ $`J_\nu \left(kx\right)J_{\nu +1}\left(ky\right)]`$ (71)
$``$ $`\left(1\right)^n\left(x\pm y\right)^{2n}\left[J_{\nu +1}\left(kx\right)J_\nu \left(ky\right)\pm J_\nu \left(kx\right)J_{\nu +1}\left(ky\right)\right],`$
and
$`_k^{2n1}[J_{\nu +1}\left(kx\right)J_\nu \left(ky\right)`$ $`\pm `$ $`J_{\nu +1}\left(kx\right)J_\nu \left(ky\right)]`$ (72)
$``$ $`\left(1\right)^{n+1}\left(x\pm y\right)^{2n1}\left[J_\nu \left(x\right)J_\nu \left(y\right)J_{\nu +1}\left(x\right)J_{\nu +1}\left(y\right)\right].`$
They have been obtained by means of recursion relations for Bessel functions neglecting terms that are suppressed by order $`1/x`$ or $`1/y`$. One can easily show that the combined powers of $`x`$ and $`y`$ in the prefactor is always larger than the combined powers of $`x`$ and $`y`$ that have been neglected. Since only even terms in $`t`$ contribute to the integral in (69) our final result for the two-point function, correct to order $`h^2`$, is given by
$`R_2(x,y)=`$ $`\overline{\rho }^2|{\displaystyle \frac{\pi ^2h\overline{\rho }\sqrt{xy}}{8}}\left[{\displaystyle \frac{J_{\nu +1}\left(x\pi \overline{\rho }\right)J_\nu \left(y\pi \overline{\rho }\right)+J_\nu \left(x\pi \overline{\rho }\right)J_{\nu +1}\left(y\pi \overline{\rho }\right)}{\mathrm{sinh}\left(\left(x+y\right)\pi ^2h\overline{\rho }/4\right)}}\right]`$ (73)
$`+`$ $`{\displaystyle \frac{\pi ^2h\overline{\rho }\sqrt{xy}}{8}}\left[{\displaystyle \frac{J_{\nu +1}\left(x\pi \overline{\rho }\right)J_\nu \left(y\pi \overline{\rho }\right)J_\nu \left(x\pi \overline{\rho }\right)J_{\nu +1}\left(y\pi \overline{\rho }\right)}{\mathrm{sinh}\left(\left(xy\right)\pi ^2h\overline{\rho }/4\right)}}\right]|^2`$
In the limit $`x,y\left|xy\right|`$ the analytical result for the two-point function of the model of Moshe, Neuberger and Shapiro is recovered from the leading order asymptotic expansion of the Bessel functions. For unitary invariant ensembles, it can be shown that the result of for critical statistics and Wigner-Dyson statistics are the only two possibilities . At this moment it is not clear whether this argument can be extended to the chiral unitary ensembles as well.
The number variance of the eigenvalues near $`x=0`$ is obtained by integrating the two-point correlation function including the self-correlations
$`\mathrm{\Sigma }^2\left(L\right)={\displaystyle _0^{L/\overline{\rho }\left(h\right)}}𝑑x{\displaystyle _0^{L/\overline{\rho }\left(h\right)}}𝑑y\left[\delta \left(xy\right)\rho \left(x\right)+R_2(x,y)\right].`$ (74)
We study its asymptotic behavior in the limit $`L\mathrm{}`$. Starting from the expressions (58) and (66), $`\mathrm{\Sigma }^2\left(L\right)`$ can be simplified by means of an orthogonality relation for Bessel functions. To leading order in $`h`$ we find
$`\mathrm{\Sigma }^2\left(L\right)`$ $`=`$ $`{\displaystyle _0^{L/\overline{\rho }\left(h\right)}}x𝑑x{\displaystyle _0^{\mathrm{}}}k𝑑k{\displaystyle \frac{J_\nu ^2\left(kx\right)}{4\mathrm{cosh}^2\frac{\beta }{2}\left(k^24n\omega \right)}}`$ (75)
$`=`$ $`{\displaystyle \frac{L}{\overline{\rho }2\pi \beta \overline{k}}}={\displaystyle \frac{h}{4}}L.`$
The same asymptotic result can be derived from the analytical result (73) (In this case the average spectral density does not depend on $`h`$ (see 62).). Such linear term, first proposed in , is believed to be characteristic for universal critical statistics valid at the mobility edge and has been related to the multifractality of the wave functions .
In Fig. 2 we compare the number variance derived from the approximate analytical result (73) and from the exact result (66). Clearly, even for a value of $`h`$ as large as 0.3 the two curve are barely distinguishable.
The asymptotic linear behavior of the number variance seems to be contradicted by the sum rule
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑y\left[\delta \left(xy\right)\rho \left(x\right)+R_2(x,y)\right]=0.`$ (76)
The resolution of this paradox is probably best illustrated by considering the Poisson ensemble for $`n`$ uncorrelated eigenvalues with average spacing $`\overline{\rho }`$. To satisfy the sum-rule we have that $`R_2(x,y)=\overline{\rho }/n`$ instead of zero for uncorrelated eigenvalues resulting in the number variance $`\mathrm{\Sigma }_2\left(L\right)=LL^2/n`$. We conclude that an asymptotic linear behavior is possible if the thermodynamic limit is taken before the limit $`L\mathrm{}`$.
To make contact with the partially quenched effective partition function, for which the number of subsequent eigenvalues around zero that are correlated according to the chGUE scales as $`\sqrt{n}`$ , we have to scale $`h`$ as $`hh/\sqrt{n}`$. In this limit microscopic universality is recovered for the interpolating chiral unitary ensemble.
## 5 Comparison with Instanton Simulations
Spectral correlations have been studied in great detail for both lattice QCD simulations and instanton-liquid simulations (see for a recent review). In lattice QCD they were studied by means of the disconnected scalar susceptibility , and complete agreement with partially quenched chiral perturbation theory was found. In particular, it was shown that the number of subsequent eigenvalues around zero described by chRMT scales as $`F^2\sqrt{V}`$. In instanton simulations a weaker volume dependence of the number of such eigenvalues was observed suggesting an approach to a critical point similar to a localization transition. Indeed the multifractality index of the fermionic wave functions was found to be nonzero. We thus compare the instanton data with the model in previous section at fixed value of the parameter $`h`$. Results for the number variance, $`\mathrm{\Sigma }^2\left(L\right)`$ versus $`L`$ are shown in Fig. 3. The closed and open circles represent results for the eigenvalues of the Dirac operator with field configurations given by an ensemble of instantons and an equal number of anti-instantons with a total density of 1 $`fm^4`$. The total number of (anti-)instantons is given in the label of the figure. In the same figure we show the result for the chGUE (dotted curve) and results for the model (66) for $`h=0.18`$ (full curve) and $`h=0.23`$ (dashed curve). We observe that both the slope and the range of agreement with the chGUE curve only shows a week volume dependence. Outside this domain the data show a linear $`L`$-dependence. Both features are nicely reproduced by the critical Random Matrix Model. The slightly positive curvature of the instanton data might be a remnant of the $`L^2`$ dependence predicted for the Altshuler-Shklovsky domain .
## 6 Conclusions
We have analyzed a chiral random matrix model that interpolates between the chGUE and the chiral Poisson ensemble. This model is a generalization of a model originally proposed by Moshe, Neuberger and Shapiro . It has been mapped onto a gas of non-interacting fermions and was solved by means of statistical mechanics methods. To leading order in the deviation from the chGUE we have obtained compact analytical expressions for the microscopic spectral density and the two-point level correlation function. We have shown that this critical chiral random matrix model provides a good description of the level correlations of Dirac eigenvalues for gauge field configurations given by a liquid of instantons.
The number variance of the critical chiral random matrix model shows a linear $`L`$-dependence for a large $`L`$ (in units of the average level spacing) whereas it coincides with the chGUE result for small values of $`L`$. The characteristic feature is that the transition point between these two domains is stable in the thermodynamic limit. This situation is very different for a non-linear $`\sigma `$-model description of disordered systems where this transition point or the Thouless energy is determined by the competition between the mass term and the kinetic term. In that case one finds the scaling behavior $`E_cD/L_s^2`$ with $`D`$ the diffusion constant and $`L_s`$ the linear size of the sample. The theoretical reason for a scale independent dimensionless conductance (i.e. the Thouless energy in units of the average level spacing) is that the localization length diverges at a critical value of the disorder. In the approach to this limit the diffusion constant has to become scale dependent. If $`E_c`$, in units of the average level spacing, becomes scale independent the diffusion constant has to be scale dependent leading to a multi-fractal scaling of the wave functions.
The weak volume dependence and the linear number variance observed in correlations of eigenvalues of the QCD Dirac operator with instanton liquid gauge field configurations suggests that we are dealing with a critical system close to a localization transition. Indeed the same numerical simulations suggest a small nonzero multifractality index of the wave-functions. On the other hand, the dimensionless conductance found in lattice QCD simulations scales according to our expectations from chiral perturbation theory. At this moment we do not have a good explanation for this discrepancy. It could simply be that the expected scaling behavior is only recovered for much large volumes in instanton simulations. Indeed, a very slow approach to the thermodynamic limit has been found for other quantities such a quenched chiral logarithms. Clearly, more work has to be done to resolve this issue.
Acknowledgements
This work was partially supported by the US DOE grant DE-FG-88ER40388. One of us (A.M.G.) was supported by “laCaixa Fellowship Program”. J.J.M.V. thanks the Institute for Nuclear Theory at the University of Washington for its hospitality and partial support during the completion of this work. Dominique Toublan is thanked for a critical reading of the manuscript. |
warning/0003/cs0003048.html | ar5iv | text | # PAL: Pertinence Action Language
## General Information
The PAL system (standing for Pertinence Action Language) is an interpreter of a causal language for describing action domains and it has been proved in PCs under Linux and Sun workstations under SunOS. It is written in C (using GNU’s compiler gcc 2.7.2.1) and sized in 150K of source code, distributed in around 6500 lines. The executable is sized in 65/80K (depending on the platform). The programming effort can be estimated in about 1 man/years.
## Description of the System
The purpose of PAL is to show the effects of using the concept of Pertinence (formally introduced in (?)) in problems of Reasoning about Actions and Change, and to help in establishing semantic features about the logical formalization of this concept. The features of Pertinence have been extracted from a practical tool for Temporal Expert Systems (mainly for the medical domain(???)) called Medtool (?) which uses a formalism very close to the current one used in PAL.
The tool is initially intended for solving temporal projection problems, being planning problems reduced to small sized ones, but we expect to use it as a starting point for future efficiency improvements.
## Applying the System
The system can be tested in two different ways, both available in the web page
```
http://www.dc.fi.udc.es/ai/~cabalar/pal
```
The two options are:
1. For an easy to use way, a simple web-based interface (PalWeb) has been implemented. It is enough with writing a domain description in the input box (or just selecting one of the preinstalled examples) and pressing the Process button to see the results.
2. The other option is downloading the source code (file pal1.2.tar.gz) and compiling it in a local machine. To this aim, follow the next steps:
```
gunzip pal1.2.tar.gz
tar -xvf pal1.2.tar
cd pal1.2
make
```
Finally, for testing the system, try any of the files in the examples directory:
```
pal < examples/yale.pal
```
We may also use the executable in the following way:
```
pal examples/yale.pal
```
so that standard input remains open to accept new queries or new action executions, providing a rudimentary interpreter.
### Methodology
No general methodology has been developed by now, although due to the closeness to Medtool, many methodological aspects can be directly inherited from practical experience in Medtool expert systems design.
### Specifics
The use of logic is the main feature of PAL with respect to previous work done on Medtool expert systems design. The purpose of PAL is to help in the logical formalization for these practical systems so that new formalism features can be incorporated or new tasks like explanation or planning can be done in a coherent way.
As an action formalism, the main features of PAL are:
1. It is narrative-based. Each model of a PAL theory has the shape of a finite sequence of situations.
2. It is a causal formalism. Rules describing the domain behavior are causal rules. Pertinence acts in a similar way to Lin’s $`Caused`$ predicate or to occlusion, although there are fundamental differences (see (?)), especially on how rule conditions are interpreted.
3. It allows concurrent actions.
4. Actions and fluents are functional, though limited to finite domains and codomains. Arithmetic expressions relating actions and fluents are allowed.
5. It allows both “real” execution and hypothetical reasoning. There exists a real narrative which can be updated in an execution trace and whose (current and past) results can be consulted at any time. Hypothetical reasoning allows to test the results of an hypothetical execution, or to find plans that satisfy a given expression for a future situation.
Each high level domain description is translated into a set of ground rules using a propositional representation with two kinds of atoms: $`holds(Q,value,situation)`$ and $`pertinent(Q,situation)`$ (being $`Q`$ a fluent or an action). The system is intended for testing different semantics for these ground rules. To this aim, the rules behavior has been implemented in a separated module so that it can be easily changed. The first approach we have done is a version based on Well Founded Semantics, since it is the closest one to the way in which the practical tool currently works. This well founded version interprets the set of rules as an objective logic program, leaving default negation just for minimizing $`pertinent`$ atoms. The iterative algorithm for computing the well founded model has been specialized so that pertinence minimization and inertia axioms are, in fact, executed implicitly.
A second option that will be soon available is using as “semantics” the inference mechanism of Medtool (which was already implemented in a C library known as Medtool-Connection). Of course, this would be an operational or procedural solution, but it is useful for comparing the practical tool results with the different logic-based versions. We also expect to incorporate a future Stable Models version (calling to smodels) and a Completion based version (using some SAT prover, following a similar technique to CCALC).
The work could influence both areas of Reasoning about Actions and Temporal Expert Systems. The interest for the former could be theoretical (the introduction of pertinence in action domains) but also practical, since PAL formalization could eventually allow studying Medtool expert systems as action theories. The interest for Temporal Expert Systems is that PAL provides an underlying logical formalization. Other less related areas could also be influenced. For instance, a strong relationship between Medtool practical formalism and Systems Theory formalizations like DEVS (Discrete Events Systems Specification) has been established ((??)). Finally, we are also studying the relation between the use of the real narrative done in PAL and approaches based on Temporal Deductive Databases.
### Users and Useability
In order to use the high level language, the user does not need to be an expert in logic, although some notions about pertinence and logic-based systems is recommendable. The basic syntax is quite direct, following the style of Medtool formalism, on which PAL is inspired. In fact, Medtool formalism is currently being used by the experts (mainly physicians), so that they directly develop the expert systems, under some minimal assistance. However, in order to use PAL as a comparison tool for different semantic implementations, a deep knowledge about Nonmonotonic Reasoning and Logic Programming semantics will be required.
As an example of flexibility, although it is not one of the initial purposes, PAL can also be indirectly used as a Logic Programming tool under Well Founded semantics. An example which includes the transformation steps to this aim can be found in PAL distribution (file wf.pal).
The PAL system is currently under continuous development and has only been used by our research group until now. However, Medtool formalism, on which PAL is inspired, has been used for constructing several real Expert Systems(???), but also in the Financial domain.
## Evaluating the System
### Benchmarks
Due to the premature stage of PAL, little comparative evaluation work has been done yet. The kind of properties usually evaluated in actions formalizations are more theoretical, especially related to flexibility and elaboration tolerance, than practical, like memory usage or response time, although several approaches are already obtaining successful results for these last parameters. As explained in the description section, many usual action topics like, for instance, qualifications, defeasible rules or delayed effects, are not covered yet by PAL, although their introduction is a subject of current research. However, we think that the current version provides enough expressivity for a practical use, properly dealing, for instance, with the frame and ramification problems (as shown by standard examples included in files yale.pal and suitcase.pal).
Unfortunately, there are not too many benchmarks for testing action approaches yet. As an exception, we could mention the recent initiative of the Logic Modelling Workshop<sup>1</sup><sup>1</sup>1http://www.ida.liu.se/ext/etai/lmw/, whose examples will be studied under PAL formalization.
As for PAL time performance, we could mention that the system works acceptably well for temporal projection problems. For instance, the execution of 5 transitions for the domain example counter.pal, which contains a causal chain of 1000 fluents, is performed in 0.85 seconds<sup>2</sup><sup>2</sup>2All times obtained in a 133MHz Pentium with 32M of RAM and under Linux Slackware 2.0.29.. Also, some small planning problems can be performed in an reasonable way. For instance, it takes 0.42 seconds for finding all the 781 ways for opening Lin’s suitcase in 5 situations (allowing concurrent actions and no execution of any action). Planning may become excessively hard, since it is attached in a naive way, generating all the possible combinations of actions and allowing concurrent execution. A nonconcurrent actions option has been added for better planning performance in most usual domains. Thus, for instance, a solution of 11 situations long for the missionaries and cannibals problem is found in 0.13 seconds, and the 4 solutions for that length are obtained in 0.7 seconds.
With respect to user-friendliness, the main advantage is the comfortability of domain descriptions involving numerical variables (limited to finite and discrete domains). For instance, the missionaries and cannibals problem can be described using 9 high level rules. With respect to applying the system, the current version is mainly batch-oriented, although it can be used as a rudimentary interpreter by entering sentences from standard input. We plan to extend this capability to a real interpreter interface.
### Comparison
As many action approaches, the PAL system can compete with planning systems, particularly in the expressivity of their representational formalism. However, no serious comparison to usual planning benchmarks has been established yet.
### Problem Size
Temporal projection domains can be reasonably handled for thousands of fluents. Planning problems, however, are quite limited yet, specially when dealing with a small/medium number of situations and concurrent execution of actions.
The system must be considered as a prototype, but several efficiency improvements (goal-directed planning, taking benefit of particular structures in the rules, heuristics, etc) are currently under study and could allow handling larger planning problems.
## An example: the blocks world
In this last section, we include, as a piece of syntax example, a possible representation of the blocks world domain. We will not provide an exhaustive description of PAL syntax, but will emphasize instead those most relevant features of the formalism. A PAL domain description consists of two parts: declarations and sentences. In the declarations part, we define the sets of fluents, actions and rules, using perhaps some auxiliary definitions of constants and sets to this aim. For instance, we will declare the following sets:
```
sets
block = [1,4];
location = block + {table};
```
Intuitively, these lines define the set of blocks as the integers from 1 to 4, and the set of locations as any block or the table. A set is usually defined as a group of elements embraced by {…} , like for instance the singleton set {table}. Elements in a set can be both symbol names (an identifier with a lower-case initial, like table) or integer numbers. An alternative way of defining a set is using the interval notation, so that is equivalent to {1,2,3,4}. Set expressions can be constructed using binary operators +,-,\* standing for union, difference and intersection respectively.
Actions and fluents are considered to be functions. Thus, an action or fluent definition consists of the descriptions of its domain and codomain using the standard mathematical notation:
```
f: set1 x set2 x set3 -> set4
```
Note that, although x is used here as an operator, it can also be used as an identifier, depending on the grammar context. As an example of functional action, we define:
```
actions
carry: block -> location;
```
so that we carry a block to a unique location. Notice that, in this way, we implicitly specify that a block cannot be placed on two different locations, since the value of a function (i.e. each carry(B)) is unique. In the example, we will use the following fluents:
```
fluents
loc: block -> location;
free: block -> {true,false};
```
When a fluent is boolean, we can simply ignore the codomain definition in the following way:
```
free: block;
```
In the same way, we could define actions or fluents without domain. Imagine an action for just marking a unique block:
```
actions
mark: -> block;
```
or a fluent that points if we have made some mark:
```
fluents
markdone: -> {true,false};
```
This last case can be further abbreviated as:
```
fluents
markdone;
```
In the declarations section we can also define variables. In PAL, variables are identified by an upper-case initial, and vary on a given sort. They are used for defining rule schemata and for making complex queries. We will use the following two block variables:
```
vars
B,C : block;
```
Finally, the declarations section would contain the set of rules describing the system behavior:
```
rules
loc(B):=carry(B);
not free(C) if carry(B)=C;
free(B) if pert(carry(C)) and prev(loc(C))=B;
false if pert(carry(B)) and not prev(free(B));
false if carry(B)=C and not prev(free(C));
```
Rules have a head of shape:
```
<fluent>:=expr
```
and an optional condition preceeded by the conditional operator if. In a rule head we allow the abbreviations:
```
<fluent>
not <fluent>
```
standing for:
```
<fluent>:=true
<fluent>:=false
```
respectively, and we allow constraint rules with a false head.
Coming back to our example, the first rule says that whatever the block B we carry, its current location loc(B) will be the value given by the carry(B) action. The second rule asserts that a block C becomes not free when we carry some B on top of it. The third rule says that if carry(C) is pertinent, that is, we carry C without regarding to which location, and the previous location of C was B then B becomes free. The fourth rule asserts that we cannot perform a carry(B) when B was not free, whereas the fifth one, avoids carrying B to a nonfree block C.
The sentences part may contain declarations about the actual narrative or queries. The actual narrative is first specified by providing the initial situation:
```
initially
loc(B):=table,free(B);
```
This means that all the blocks are on the table and are free. Performing actions in the actual narrative can be done in the following way:
```
do {carry(1):=2;}
```
This would carry block 1 on top of block 2. The output would be the following one:
```
1)
carry(1):=2
loc(1):=2
free(2):=false
```
The output shows only the pertinent facts, just pointing out the performed action and its derived effects (all the rest has persisted). We can perform several actions in a sequence:
```
do {carry(1):=table;carry(2):=3; carry(1):=2;}
```
which would have as output:
```
2)
carry(1):=table
loc(1):=table
free(2):=true
3)
carry(2):=3
loc(2):=3
free(3):=false
4)
carry(1):=2
loc(1):=2
free(2):=false
```
Notice how performing actions is an incremental task: the execution of the sequence generates situations from 2 to 4, since it was performed taking 1 as the initial situation (that is, the one we had obtained previously). We can reinitialize (“resume”) the actual narrative by providing a new initially clause. Besides, actions can also be performed concurrently, or we can even perform no action:
```
initially
loc(B):=table,free(B);
do {carry(1):=2,carry(3):=4; carry(1):=3; ; }
```
that is, we first simultaneously move 1 on top of 2 and 3 on top of 4, in the next situation, we move 1 to 3, and finally we perform no action. The output would be:
```
Resume
1)
carry(1):=2
carry(3):=4
loc(1):=2
loc(3):=4
free(2):=false
free(4):=false
2)
carry(1):=3
loc(1):=3
free(2):=true
free(3):=false
3)
```
Finally, we will briefly comment the other kind of sentences: queries. The simplest kind of query allows asking about the current state. For instance, we can check that block 2 is free and it is on table:
```
query
free(2) and loc(2)=table?
yes
```
Variables can be used in a similar way to Prolog queries. For instance, we can find all the blocks that are not free:
```
query
not free(B)?
B=3
B=4
2 solutions
```
The general shape of queries allows asking properties about an hypothetical future. Answers will have the shape of plans, that is, sequences of actions to be performed for satisfying the query. For making some tests, we will add the assumption of non-concurrency of actions (otherwise, we should add a rule for specifying that there is not space for two blocks on top of a third block). The non-concurrency assumption is specified in the declarations part, for instance, as:
```
options
not concurrent;
```
Non-concurrency also assumes that an (unique) action is mandatorily performed. Now, using the initial state:
```
initially
loc(B):=table,free(B);
```
we can try to make block 3 not free in the next situation:
```
query
true;not free(3) ?
\noindent obtaining 4 possible answers:
\begin{verbatim}
Resume
Solution 1:
1)
carry(1):=3
loc(1):=3
free(3):=false
Solution 2:
1)
carry(2):=3
loc(2):=3
free(3):=false
Solution 3:
1)
carry(3):=3
loc(3):=3
free(3):=false
Solution 4:
1)
carry(4):=3
loc(4):=3
free(3):=false
4 solutions
```
Notice that in the current (present) situation we do not require anything and we use the true expression. When this happens, we can abbreviate it as an empty expression like in:
```
query ; ; ; not free(3)?
```
that would look for the ways of getting block 3 not free after 3 steps (this generates 1072 solutions!). We can fix the number of solutions we wish by another option:
```
options
solutions=1;
```
It is also possible to repeat the same expression a number of times along several situations. For instance, we can ask the queries:
```
query
...{3} not free(3) ?
free(1) ...{3} not free(3)?
```
so that the first one replaces our immediately previous example, and the second one tries to find ways of occupying block 3 in 3 steps but maintaining block 1 free.
#### Acknowledgments.
This research is partially supported by the Government of Spain, grant PB97-0228. |
warning/0003/cond-mat0003258.html | ar5iv | text | # Wavefronts may move upstream in semiconductor superlattices
## I Introduction
Current instabilities in doped semiconductor superlattices (SL) have been an active subject of research during this decade. For strongly coupled SL, Bloch oscillations and Wannier-Stark hopping produce negative differential conductivity (NDC) at high electric fields. This may result in self-sustained oscillations of the current due to recycling of charge dipole domains as in the Gunn effect of bulk n-GaAs . For weakly coupled SL, sequential tunneling is the main mechanism of vertical transport. Under dc voltage bias conditions, stationary electric field domains may form if doping is large enough . Below a critical doping value, the existing charge inside the SL may not be able to pin domain walls, and current self-oscillations appear . These oscillations may be due to recycling of charge monopoles (domain walls) or dipoles depending on the boundary condition at the injecting contact region (in a typical n<sup>+</sup>-n-n<sup>+</sup> configuration with the SL imbedded between highly doped regions, the doping at the emitter region is crucial) . Driven chaotic oscillations have also been predicted and observed in experiments . Lastly, there are ways to tune the charge inside the SL (and therefore obtain stationary domains or self-oscillations) without replacing it by a different one. For example, by applying a transverse magnetic field or by photoexciting the SL .
Transport in weakly coupled SL can be described by simple rate equation models for electron densities and average fields in the wells, . Many of the effects related above have been explained by means of a simple discrete drift model . In this model, the tunneling current between two adjacent wells, $`J_{ii+1}`$, equals the 2D electron charge density at well $`i`$ times a drift velocity, which depends on the electric field at the same well. By starting from a microscopic sequential tunneling model, it has been shown that the discrete drift model is a good approximation at low temperatures and for fields above the first plateau of the SL current-voltage characteristic . For low dc voltages on the first plateau, a discrete diffusion (which is a nonlinear function of the field) should be added. This term contains the contribution to $`J_{ii+1}`$ of the tunneling from well $`i+1`$ back to well $`i`$ (which vanishes for large enough electric fields) . In this paper we report an interesting consequence of electron diffusivity at low fields: if the current is sufficiently high, and so is the doping, a domain wall (monopole wave) which connects two domains may travel in a direction opposite to the flow direction for electrons (i.e., upstream, in the positive current direction!). This striking phenomenon is contrary to the usual situation: a monopole either moves downstream (in the direction of the flow of electrons), or it remains stationary, . We substantiate our claim both by numerical simulations of the discrete drift-diffusion model and by rigorous mathematical analysis based upon a comparison principle . Mathematical analysis yields useful bounds for critical values of current and well doping, and for monopole velocity.
There are related fields for which differential-difference equations (similar to discrete drift-diffusion models) model the systems of interest. Well-known are propagation of nerve impulses along myelinated fibers, modelled by discrete FitzHugh-Nagumo equations , motion of dislocations and sliding charge density waves , modelled by variants of the Frenkel-Kontorova model , etc. The theory of wavefront propagation has been developed for some of these models, which are simpler than ours: convection is typically absent from them and diffusion is purely linear .
The rest of the paper is as follows. We write the drift-diffusion model with appropriate boundary conditions in Section II. There we render these equations dimensionless and explain the results of numerical simulations on a current biased infinitely long SL. Furthermore, we find by numerical simulations that our results for infinite SL may be realized in finite SL with appropriate boundary conditions under constant current bias. The theoretical analysis based on the comparison principle is presented in Section III. Section IV contains our conclusions. Finally some material of a more technical nature is relegated to the Appendices.
## II Discrete drift-diffusion model
### A Equations and boundary conditions
At low enough temperatures (much less than a typical Fermi energy of a SL well measured from the first subband, say 20 meV or 232 K), the following discrete drift-diffusion equations model sequential vertical transport in a weakly doped SL :
$`{\displaystyle \frac{\epsilon }{e}}{\displaystyle \frac{dF_i}{dt}}+{\displaystyle \frac{n_iv(F_i)}{d+w}}D(F_i){\displaystyle \frac{n_{i+1}n_i}{(d+w)^2}}=J(t),`$ (1)
$`F_iF_{i1}={\displaystyle \frac{e}{\epsilon }}(n_iN_D^w).`$ (2)
Eq. (1) is Ampère’s law establishig that the total current density, $`eJ`$, is sum of displacement and tunneling currents. The latter consists of a drift term, $`en_iv(F_i)/(d+w)`$, and a diffusion term, $`eD(F_i)(n_{i+1}n_i)/(d+w)^2`$. We have adopted the convention (usual in this field) that the current density has the same direction as the flow of electrons. Eq. (1) holds for $`i=1,\mathrm{},N1`$. Eq. (2) is the Poisson equation, and it holds for $`i=1,\mathrm{},N`$. $`n_i`$ is the 2D electron number density at well $`i`$, which is singularly concentrated on a plane located at the end of the well. $`F_i`$ is minus an average electric field on a SL period comprising the $`i`$th well and the $`i`$th barrier (well $`i`$ lies between barriers $`i1`$ and $`i`$: barriers 0 and $`N`$ separate the SL from the emitter and collector contact regions, respectively). Parameters $`\epsilon `$, $`d`$, $`w`$, and $`N_D^w`$ are well permittivity, barrier width, well width and 2D doping in the wells, respectively.
Drift velocity and diffusion coefficient are depicted in Fig. 1 for the 9nmGaAs/4nmAlAs SL of Ref. . We have obtained them from microscopic calculations presented in Ref. (which is appropriate for these sample parameters ) by setting $`v(F)=J(N_D^w,N_D^w,F)(d+w)/N_D^w`$ and $`D(F)=(J(N_D^w,N_D^w,F)/n_{i+1})(d+w)^2`$. Here $`eJ(n_i,n_{i+1},F_i)`$ is the tunneling current between wells $`i`$ and $`i+1`$, $`J_{ii+1}`$. We assume that the tunneling current is a function of the average field at the $`i`$th SL period, $`F_i=F`$, and of the 2D electron densities at wells $`i`$ and $`i+1`$, $`n_i`$ and $`n_{i+1}`$, respectively. Notice that our model for the tunneling current,
$`eJ(n_i,n_{i+1},F_i)={\displaystyle \frac{en_iv(F_i)}{d+w}}eD(F_i){\displaystyle \frac{n_{i+1}n_i}{(d+w)^2}}`$ (3)
$`{\displaystyle \frac{en_iv^{(f)}(F_i)en_{i+1}v^{(b)}(F_i)}{d+w}},`$ (4)
is reasonable for temperatures much lower than a typical Fermi energy in the wells measured from the first subband (say 20 meV), . The tunneling current density should change sign if we reverse the electric field and exchange the electron densities at wells $`i`$ and $`i+1`$: $`J(n_i,n_{i+1},F_i)=J(n_{i+1},n_i,F_i)`$. This inversion symmetry implies
$$v^{(f)}(F)=v^{(b)}(F)\text{and}v(F)=v(F),$$
where $`v^{(b)}(F)=D(F)/(d+w)`$ and $`v^{(f)}(F)=v(F)+v^{(b)}(F)`$. See Figure 1(d).
Equations (1) and (2) should be supplemented with appropriate bias, initial and boundary conditions. Among possible bias conditions, we shall consider the extreme cases of current bias ($`J(t)`$ specified) and voltage bias:
$`(d+w){\displaystyle \underset{i=1}{\overset{N}{}}}F_i=V,`$ (5)
with specified $`V=V(t)`$. Using (5) ignores potential drops at the contact regions and at barrier 0, and it overestimates the contribution of barrier $`N`$ by a factor $`1+w/d`$ . These contributions are negligible for long SL ($`N=40`$ or larger), so that we shall adopt the simpler expression (5). Appropriate boundary conditions have been derived under the same approximations as in (1) . They are
$`{\displaystyle \frac{\epsilon }{e}}{\displaystyle \frac{dF_0}{dt}}+j_e^{(f)}(F_0){\displaystyle \frac{n_1w^{(b)}(F_0)}{d+w}}=J(t),`$ (6)
$`{\displaystyle \frac{\epsilon }{e}}{\displaystyle \frac{dF_N}{dt}}+{\displaystyle \frac{n_Nw^{(f)}(F_N)}{d+w}}=J(t),`$ (7)
where the emitter current density, $`ej_e^{(f)}(F)`$, the emitter backward velocity, $`w^{(b)}(F)`$, and the collector forward velocity, $`w^{(f)}(F)`$ are functions of the electric field depicted in Fig. 3 of Ref. for contact regions similar to those used in experiments .
To analyze the discrete drift-diffusion model, it is convenient to render all equations dimensionless. Let $`v(F)`$ reach its first positive maximum at $`(F_M,v_M)`$. We adopt $`F_M`$, $`N_D^w`$, $`v_M`$, $`v_M(d+w)`$, $`eN_D^wv_M/(d+w)`$ and $`\epsilon F_M(d+w)/(eN_D^wv_M)`$ as the units of $`F_i`$, $`n_i`$, $`v(F)`$, $`D(F)`$, $`eJ`$ and $`t`$, respectively. For the first plateau of the 9/4 SL of Ref. , we find $`F_M=6.92`$ kV/cm, $`N_D^w=1.5\times 10^{11}`$ cm<sup>-2</sup>, $`v_M=156`$ cm/s, $`v_M(d+w)=2.03\times 10^4`$ cm$`{}_{}{}^{2}/`$s and $`eN_D^wv_M/(d+w)=2.88`$ A/cm<sup>2</sup>. The units of current and time are 0.326 mA and 2.76 ns, respectively. Then (1) to (5) become
$`{\displaystyle \frac{dE_i}{dt}}+v(E_i)n_iD(E_i)(n_{i+1}n_i)=J,`$ (8)
$`E_iE_{i1}=\nu (n_i1),`$ (9)
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}E_i=\varphi .`$ (10)
Here we have used the same symbols for dimensional and dimensionless quantities except for the electric field ($`F`$ dimensional, $`E`$ dimensionless). The parameters $`\nu =eN_D^w/(\epsilon F_M)`$ and $`\varphi =V/[F_MN(d+w)]`$ are dimensionless doping and average electric field (bias), respectively. For the 9/4 SL, $`\nu 3`$. We recall that $`i=1,\mathrm{},N1`$ in (8) and $`i=1,\mathrm{},N`$ in (9). The boundary conditions (6) and (7) become
$`{\displaystyle \frac{dE_0}{dt}}+J_e(E_0)w_e(E_0)n_1=J,`$ (11)
$`{\displaystyle \frac{dE_N}{dt}}+w_c(E_N)n_N=J,`$ (12)
where
$`J_e(E_0)={\displaystyle \frac{j_e^{(f)}(F_ME_0)(d+w)}{N_D^wv_M}},`$ (13)
$`w_e(E_0)={\displaystyle \frac{w^{(b)}(F_ME_0)}{v_M}},`$ (14)
$`w_c(E_N)={\displaystyle \frac{w^{(f)}(F_ME_N)}{v_M}}.`$ (15)
Figure 2 shows $`J_e`$, $`w_e`$ and $`w_c`$ as functions of the electric field. They are dimensionless versions of the curves plotted in Figure 3 of Ref. .
### B Numerical simulations
Simple solutions of the drift-diffusion equations (8) - (9) under constant current bias are stationary or moving monopole wavefronts connecting two electric field domains. Let us consider monopole solutions with profiles $`\{E_i\}`$ which are increasing functions of $`i`$, for they are compatible with realistic boundary conditions in which the emitter region is highly doped . We have simulated numerically on a large SL,
$`{\displaystyle \frac{dE_i}{dt}}{\displaystyle \frac{D(E_i)+v(E_i)}{\nu }}(E_{i1}E_i)`$ (16)
$`{\displaystyle \frac{D(E_i)}{\nu }}(E_{i+1}E_i)=Jv(E_i),`$ (17)
with fixed $`J`$, which is equivalent to (8) - (9). Let $`E^{(1)}(J)<E^{(2)}(J)<E^{(3)}(J)`$, be the three solutions of $`v(E)=J`$ for $`v_m<J<1`$, where $`(E_m,v_m)`$ is the minimum of $`v(E)`$ for $`E>1`$. For the 9/4 SL of Fig. 1, $`E_m=9.8571`$, $`v_m=0.02192`$. We have simulated (17) for different values of $`\nu >0`$ and of $`J(v_m,1)`$. The initial condition was chosen so that $`E_iE^{(1)}(J)`$ as $`i\mathrm{}`$, and $`E_iE^{(3)}(J)`$ as $`i\mathrm{}`$. We observed that, after a short transient, a variety of initial conditions sharing these features evolved towards either a stationary or moving monopole. For systematic numerical studies, we therefore adopted an initial step like profile, with $`E_i=E^{(1)}(J)`$ for $`i<0`$, $`E_i=E^{(3)}(J)`$ for $`i>0`$ and $`E_0=E^{(2)}(J)`$. The boundary data were taken to be $`E_N=E^{(1)}(J)`$, $`E_N=E^{(3)}(J)`$ with $`N`$ large.
Our results show that the dimensionless doping $`\nu `$ determines the type of solution of (17) which is stable. There are two important values of $`\nu `$, $`\nu _1<\nu _2`$.
* For $`0<\nu <\nu _1`$ and each fixed $`J(v_m,1)`$, only traveling monopole fronts moving downstream (to the right) were observed. For $`\nu >\nu _1`$, stationary monopoles were found. According to the arguments of Wacker et al for the discrete drift model with $`D(E)=0`$, stationary monopoles exist for dimensionless doping larger than a critical value. An upper bound for this critical doping is
$`\nu _c=v_m{\displaystyle \frac{E_m1}{1v_m}},`$ (18)
which equals $`\nu _c=0.198`$ for our numerical example. We have found that $`\nu _1=0.16`$. This agreement with results obtained assuming $`D(E)=0`$ is not surprising: we shall prove in Section III that (18) holds as well for the model (8) - (9) with nonzero diffusivity.
* For $`\nu _1<\nu <\nu _2`$, traveling fronts moving downstream exist only if $`J(v_m,J_1(\nu ))`$, where $`J_1(\nu )<1`$ is a critical value of the current. If $`J(J_1(\nu ),1)`$, the stable solutions are steady fronts (stationary monopoles). We have found that $`\nu _2=0.33`$.
* New solutions are observed for $`\nu >\nu _2`$. As before, there are traveling fronts moving downstream if $`J(v_m,J_1(\nu ))`$, and stationary monopoles if $`J(J_1(\nu ),J_2(\nu ))`$, $`J_2(\nu )<1`$ is a new critical current. For $`J_2(\nu )<J<1`$, the stable solutions of (17) are monopoles traveling upstream (to the left). As $`\nu `$ increases, $`J_1(\nu )`$ and $`J_2(\nu )`$ approach $`v_m`$ and 1, respectively. Thus stationary solutions are found for most values of $`J`$ if $`\nu `$ is large enough.
Figure 3 depicts $`J_1(\nu )`$ and $`J_2(\nu )`$ as functions of $`\nu `$. Notice that $`J_1`$ decreases from $`J_1=1`$ to $`J_1=v_m`$ as $`\nu `$ increases from $`\nu _1`$. Similarly, $`J_2`$ decreases from $`J_2=1`$ to a minimum value $`J_20.53`$ and then increases back to $`J_2=1`$ as $`\nu `$ increases. Monopole velocity as a function of current has been depicted in Figure 4 for four different doping values, $`\nu =0.5`$, $`\nu =1`$, $`\nu =3`$ and $`\nu =10`$. For larger $`\nu `$, the interval of $`J`$ for which stationary solutions exist becomes wider again, trying to span the whole interval $`(v_m,1)`$ as $`\nu \mathrm{}`$. For very large $`\nu `$, the velocities of downstream and upstream moving monopoles become extremely small in absolute value.
Notice that if we use the complete sequential tunneling current instead of the drift-diffusion approximation (4) in Eq. (1), the situation is the same. Figure 5 depicts monopole velocity versus current for well doping corresponding to the 9/4 SL of Ref. . Results obtained with the complete sequential tunneling current or with approximation (4) (corresponding to Fig. 4 with $`\nu =3`$) are compared. Both velocity curves are similar, and their quantitative discrepancies are irrelevant in view of the uncertainties involved in a theoretical calculation of the tunneling current (typically the off-resonance current is larger than the theoretical prediction).
Once different stable monopole solutions (moving either downstream or upstream, stationary) have been identified, we raise the natural question of whether they are compatible with boundary conditions. Another series of numerical simulations was carried out to answer this. We solved numerically (17) for a current biased finite SL ($`N=40`$) with boundary conditions (11) - (15). Our results are depicted in Figure 6 for realistic doping at the contact layers. We observe that the emitter boundary condition results in the creation of a charge accumulation layer near this contact. A charge depletion layer is formed near the collector contact as a result of the corresponding boundary condition. Except for these layers, existence and configuration of monopoles moving downstream, upstream or remaining stationary, agrees with the previous simulations (corresponding to an infinitely long current-biased SL with a monopole-like initial condition).
## III Mathematical analysis of traveling monopoles and stationary solutions
In this Section, we study theoretically moving or stationary monopoles on an infinitely long, current-biased SL. Our findings will confirm the picture suggested by the numerical simulations of the previous Section for any doped weakly coupled SL. Furthermore, we shall prove stability of the different monopole solutions and find bounds for the critical values of $`\nu `$ and $`J_i`$. Our results are based upon and extend ideas first proposed by J.P. Keener for discrete FitzHugh-Nagumo equations, corresponding to signal transmision in myelinated neurons . Mathematically analogous problems arise in models of propagation of defects in crystals . These problems have the following structure,
$`{\displaystyle \frac{dE_i}{dt}}d(E_{i+1}2E_i+E_{i1})=Jv(E_i),`$ (19)
which is much simpler than (17). Here the parameter $`d>0`$ is a constant diffusion coefficient, and $`v(E)`$ a ‘cubic’ function with three branches as the electron drift velocity of Fig. 1.
For (19), there are critical values of $`J`$, $`J_1`$ and $`J_2`$, characterizing wavefront behavior . For $`J>J_2(d)`$, there exist wavefront solutions of (19) moving upstream (to the left). For $`J<J_1(d)`$, there are wavefronts moving downstream (to the right), whereas for $`J_1(d)<J<J_2(d)`$, stationary fronts exist. The width of the interval $`(J_1(d),J_2(d))`$ is an increasing function of $`d`$.
### A Propagation failure and stationary solutions
In Appendix A we state and prove a comparison principle for (17). As a consequence, if our initial field profile is monopole-like \[monotone increasing with well index, and sandwiched between $`E^{(1)}(J)`$ and $`E^{(3)}(J)`$\], so is the electric field profile for any later time $`t>0`$; see Appendix A:
$`\{E_i(0)\}\text{increasing with }i`$ (20)
$`E^{(1)}(J)<E_i(t)<E_{i+1}(t)<E^{(3)}(J),i,t>0.`$ (21)
We now obtain sufficient conditions for an initial monopole not to propagate upstream or downstream. Under these conditions, the monopole may remain stationary or move downstream or upstream, respectively. Let us start with a condition pinning the left tail of a monopole. As $`E_{i1}<E_i`$ and $`E_{i+1}<E^{(3)}(J)`$, we have:
$`{\displaystyle \frac{dE_i}{dt}}`$ $`=`$ $`{\displaystyle \frac{D(E_i)}{\nu }}(E_{i+1}E_i)`$ (22)
$`+`$ $`{\displaystyle \frac{D(E_i)+v(E_i)}{\nu }}(E_{i1}E_i)+Jv(E_i)`$ (23)
$``$ $`{\displaystyle \frac{D(E_i)}{\nu }}[E^{(3)}(J)E_i]+Jv(E_i)0,`$ (24)
provided there exist $`a_l<b_l`$ such that
$`{\displaystyle \frac{D(E)}{\nu }}[EE^{(3)}(J)]Jv(E),E(a_l,b_l),`$ (25)
and then we choose some initial field, $`E_i(0)(a_l,b_l)`$. The previous inequality then implies $`E_i(t)(E^{(1)}(J),b_l)`$ for all $`t>0`$. This in turn forbids a monopole to move upstream (to the left). We say that condition (25) pins the left tail of the monopole. Whether such $`(a_l,b_l)`$ exist, depends on the parameters $`\nu `$ and $`J`$; see Figure 7.
Let us now pin the right tail of a monopole. As $`E_{i+1}>E_i`$ and $`E_{i1}>E^{(1)}(J)`$, we have
$`{\displaystyle \frac{dE_i}{dt}}`$ $`=`$ $`{\displaystyle \frac{D(E_i)}{\nu }}(E_{i+1}E_i)`$ (26)
$`+`$ $`{\displaystyle \frac{D(E_i)+v(E_i)}{\nu }}(E_{i1}E_i)+Jv(E_i)`$ (27)
$``$ $`{\displaystyle \frac{D(E_i)+v(E_i)}{\nu }}[E^{(1)}(J)E_i]+Jv(E_i)0,`$ (28)
provided there exist $`a_r<b_r`$ such that
$`{\displaystyle \frac{D(E)+v(E)}{\nu }}[EE^{(1)}(J)]Jv(E),`$ (29)
$`E(a_r,b_r),`$ (30)
and we choose some initial field, $`E_i(0)(a_r,b_r)`$. The previous inequality then implies $`E_i(t)(a_r,E^{(3)}(J))`$ for all $`t>0`$. A monopole cannot then move downstream (to the right), and we say that its right tail is pinned. Figure 8 illustrates our arguments: for fields larger than $`a_r`$, the $`E_i`$’s tend to increase above $`a_r`$ toward $`E^{(3)}(J)`$. Then the monopole cannot move downstream. For $`E_i<b_l`$, the fields tend to $`E^{(1)}(J)`$, and the monopole cannot move upstream. As before, the existence of $`(a_r,b_r)`$ depends on the values of $`\nu `$ and $`J`$; see Figure 7.
Figures 7 show the curves $`Jv(E)`$, $`D(E)(EE^{(3)})/\nu `$ and $`[D(E)+v(E)](EE^{(1)})/\nu `$ for $`\nu =3`$ and different values of $`J`$. At $`J=0.08`$, Figure 7(a) shows that there is an interval $`(a_l,b_l)`$ as in (25), but no interval $`(a_r,b_r)`$ as in (30) exist. Then the left tail of a monopole is pinned, but its right tail is free. In this conditions, a monopole may move downstream. Figure 7(b) shows a monopole with both its left and right tails pinned for $`J=0.2`$. Then our theory implies that wavefront propagation fails and a monopole-like stationary solution is stable. Numerical simulations show that there are stationary solutions when $`J[0.09,0.53]`$. Finally, Figure 7(c) shows that, if $`J=0.34`$, the right tail of a monopole is pinned, but not its left tail. Under these conditions a monopole may move upstream. For larger $`\nu `$, the estimates become sharper. For instance, when $`\nu =10`$, (25) and (30) hold for $`J[0.05,0.45]`$. Direct numerical simulations show that stationary solutions exist for $`J[0.04,0.55]`$. Systematic use of these criteria allows us to estimate the critical doping values $`\nu _j`$ and critical current values $`J_i(\nu )`$, $`i=1,2`$ defined in the previous Section; see Appendix B. Instead of looking for $`J_1(\nu )`$ and $`J_2(\nu )`$, it is more convenient to look for their inverse functions, which we may call $`\nu _1(J)`$ and $`\nu _2^\pm (J)`$. According to Figure 3, the inverse function of $`J_2(\nu )`$ is two-valued, and its two branches are $`\nu _2^{}(J)<\nu _2^+(J)`$. We have found the following upper bounds $`\nu _{1b}(J)`$ and $`\nu _{2b}^+(J)`$ for $`\nu _1(J)`$ and $`\nu _2^+(J)`$, respectively:
$`\nu _{1b}(J)=v_m{\displaystyle \frac{E_mE^{(1)}(J)}{Jv_m}},`$ (31)
$`\nu _{2b}^+(J)=D(1){\displaystyle \frac{E^{(3)}(J)1}{1J}}.`$ (32)
If $`\nu >\nu _{1b}(J)`$, the right tail of the monopole is pinned, whereas the left tail of the monopole is pinned if $`\nu >\nu _{2b}^+(J)`$; see Appendix B. Notice that $`\nu _{1b}(J)`$ is a decreasing function of $`J`$. Therefore the critical value $`\nu _1`$ (above which there are stationary solutions) is smaller than $`\nu _{1b}(1)`$, which is exactly Wacker et al’s bound, (18). This explains why the bound (18) gives surprisingly good results even for the first plateau of the SL current-voltage characteristics \[despite having been obtained under the assumption $`D(E)0`$\] . Notice that the bound (32) is reasonable for large dopings and currents $`J1`$. See Figure 3 for a comparison between the critical curves $`J_1(\nu )`$ and $`J_2(\nu )`$ and the bounds (31) and (32).
### B Propagation: traveling fronts
Having shown that only one tail of a monopole is pinned suggests that the monopole may move in the opposite direction. Direct simulations show that this is often the case, and we will prove this now.
An upstream traveling wave solution of (17) may have the form
$`E_i(t)=w(i+ct),c>0.`$ (33)
We will look for an electric field profile $`w(z)`$, $`z=i+ct`$, which is not an exact solution of (17), but instead it satisfies
$`c{\displaystyle \frac{dw}{dz}}{\displaystyle \frac{D(w(z))+v(w(z))}{\nu }}[w(z1)w(z)]`$ (34)
$`{\displaystyle \frac{D(w(z))}{\nu }}[w(z+1)w(z)]+v(w(z))J0.`$ (35)
If this subsolution is initially below an initial field profile, i.e. $`w(i)<E_i(0)`$, for all $`i`$, then the comparison theorem of Appendix A guarantees that $`E_i(t)>w(i+ct)`$ for later times. As $`w(i+ct)`$ moves upstream, so does $`E_i(t)`$, and the electric field profile corresponds to a monopole moving upstream with velocity at least $`c`$. See Figure 9: a subsolution “pushes” the monopole upstream, whereas a supersolution (defined below) “pushes” the monopole downstream.
How do we find a reasonable subsolution? An idea is to try a piecewise continuous solution which equals $`E^{(1)}(J)`$ for $`z<z_0`$ and a larger constant $`A`$, $`A(E^{(2)}(J),E^{(3)}(J))`$ \[and therefore $`v(A)J0`$\], for $`z>z_1`$, with $`z_1>z_0`$. For $`z_0<z<z_1`$, $`w(z)`$ is an unspecified smooth increasing function with $`w(z_0)=E^{(1)}(J)`$ and $`w(z_1)=A`$. Now we shall select conveniently the numbers $`z_0`$, $`z_1`$, $`c`$ and $`A`$ so that (35) holds. Clearly, (35) holds for $`z+1<z_0`$ and for $`z1>z_1`$. Suppose that $`0<z_1z_0<1`$. Then there are five possibilities:
1. $`z<z_0`$, $`z+1>z_1`$. Then $`w(z1)=w(z)=E^{(1)}(J)`$, $`w(z+1)=A`$, which inserted in (35) yields $`D(E^{(1)})(AE^{(1)})/\nu 0`$ (obviously true).
2. $`z1<z_0`$, $`z>z_1`$. Then $`w(z1)=E^{(1)}(J)`$, $`w(z)=w(z+1)=A`$, which inserted in (35) yields
$`Jv(A){\displaystyle \frac{D(A)+v(A)}{\nu }}[AE^{(1)}(J)].`$ (36)
3. $`z_0<z1<z_1`$. Then $`E^{(1)}(J)<w(z1)<A`$ and $`w(z)=w(z+1)=A`$, which yields $`Jv(A)[D(A)+v(A)][Aw(z1)]/\nu `$. This inequality holds if (36) does.
4. $`z_0<z<z_1`$. Then $`w(z1)=E^{(1)}(J)`$, $`E^{(1)}(J)<w(z)<A`$ and $`w(z+1)=A`$. Inserting this in (35), we find
$`c{\displaystyle \frac{dw}{dz}}{\displaystyle \frac{D(w(z))+v(w(z))}{\nu }}[E^{(1)}w(z)]`$ (37)
$`+{\displaystyle \frac{D(w(z))}{\nu }}[Aw(z)]+Jv(w(z)).`$ (38)
Let us now assume that we can select $`A(E^{(2)},E^{(3)})`$ such that the right hand side of this expression is positive, say
$`{\displaystyle \frac{D(w)+v(w)}{\nu }}[E^{(1)}(J)w]+{\displaystyle \frac{D(w)(Aw)}{\nu }}`$ (39)
$`+Jv(w)\delta >0,E^{(1)}<w<A,`$ (40)
and that we choose $`c`$ so that $`cdw/dz<\delta `$. Then (35) holds.
5. $`z_0<z+1<z_1`$. Then $`w(z1)=w(z)=E^{(1)}(J)`$, and $`E^{(1)}(J)<w(z+1)<A`$, which inserted in (35) yields $`D(E^{(1)})[w(z+1)E^{(1)}]/\nu 0`$ (obviously true).
Summarizing the previous arguments, provided (36) and (40) hold, $`w(z)`$ is a subsolution obeying (35). The parameter $`A`$ can be found graphically. First of all, we depict the functions $`Jv(E)`$ and $`f_1(E;J)[D(E)+v(E)][EE^{(1)}(J)]/\nu `$. Possible values of $`A`$ are those $`E`$ for which $`Jv(E)f_1(E;J)`$. For such A, we may plot the left side of (40),
$`f_2(E;J,A){\displaystyle \frac{D(E)+v(E)}{\nu }}[E^{(1)}(J)E]`$ (41)
$`+{\displaystyle \frac{D(E)(AE)}{\nu }}+Jv(E).`$ (42)
If $`f_2(E;J,A)>0`$ for $`E(E^{(1)}(J),A)`$, then the selected value of $`A`$ allows us to construct the sought subsolution. See Figure 10 for a practical realization of this graphical construction.
We have proved rigorously that monopoles may move upstream under favorable circumstances. Our proof using subsolutions may yield a very practical additional bonus: an upper bound, $`c^{}`$, for the velocity of the monopole. Let us choose $`\delta (J,A)=`$ min$`{}_{E^{(1)}<E<A}{}^{}f_{2}^{}(E;J,A)`$, $`z_1z_0=1`$, and $`w(z)=[AE^{(1)}(J)](zz_0)`$ for $`z_0<z<z_1`$. Then $`c^{}=\delta (J,A)/[AE^{(1)}(J)]`$. In Figure 4, $`c^{}`$ is represented by a line of thick dots for doping $`\nu =3`$ corresponding to the 9/4 SL.
In a similar vein, we can construct supersolutions which push the monopole field profile to the right; see Figure 9. Now we start from a monopole profile moving downstream,
$`E_i(t)=w(ict),c>0.`$ (43)
The electric field profile $`w(z)`$, $`z=ict`$ should satisfy
$`c{\displaystyle \frac{dw}{dz}}+{\displaystyle \frac{D(w(z))+v(w(z))}{\nu }}[w(z1)w(z)]`$ (44)
$`+{\displaystyle \frac{D(w(z))}{\nu }}[w(z+1)w(z)]+Jv(w(z))0.`$ (45)
We seek a piecewise continuous supersolution which equals a constant, $`A`$, $`A(E^{(1)}(J),E^{(2)}(J))`$, for $`z<z_0`$, and $`w(z)=E^{(3)}(J)`$ for $`z>z_1`$, with $`z_1>z_0`$. For $`z_0<z<z_1`$, $`w(z)`$ is an unspecified smooth increasing function with $`w(z_0)=A`$, and $`w(z_1)=E^{(3)}(J)`$. As for subsolutions, we now select conveniently the numbers $`z_0`$, $`z_1`$, $`c`$ and $`A`$ so that (45) holds. Clearly, (45) holds for $`z+1<z_0`$ and for $`z1>z_1`$. Suppose that $`0<z_1z_0<1`$. An analysis of the remaining five possibilities yields the following criteria to hold for $`w(ict)`$ to be a supersolution:
$`Jv(A){\displaystyle \frac{D(A)}{\nu }}[E^{(3)}(J)A]f_3(A;J),`$ (46)
$`f_4(w;J,A){\displaystyle \frac{D(w)+v(w)}{\nu }}(wA)`$ (47)
$`+{\displaystyle \frac{D(w)}{\nu }}[E^{(3)}(J)w]+Jv(w)\delta ,`$ (48)
$`\text{for }AwE^{(3)}(J),`$ (49)
$`c{\displaystyle \frac{dw}{dz}}\delta .`$ (50)
Provided such $`w(ict)`$ is found, solutions $`E_i(t)`$ of (17) with $`E_i(0)<w(i+\tau )`$ will satisfy $`E_i(t)<w(ict+\tau )`$ and propagate to the right with speed larger than $`c`$. $`\tau `$ is a constant which can be conveniently chosen to keep the monopole profile below the supersolution. Figure 11 illustrates the graphical construction of the supersolution by checking that (46) and (48) hold for particular values of $`J`$ and $`\nu `$.
As in the subsolution case, an upper bound $`c^{}`$ for the monopole velocity $`c`$ is estimated by choosing $`\delta (J,A)=`$ max$`{}_{A<E<E^{(3)}}{}^{}f_{4}^{}(E;J,A)`$, $`z_1z_0=1`$, and $`w(z)=[E^{(3)}(J)A](zz_0)`$ for $`z_0<z<z_1`$. Then $`c^{}=\delta (J,A)/[E^{(3)}(J)A]`$. In Figure 4, $`c^{}`$ is represented by a line of thick dots for doping $`\nu =3`$ corresponding to the 9/4 SL.
We can now summarize the results obtained from sub and supersolutions; see Figures 10 and 11. We find reasonably good upper bounds for the absolute value of monopole velocity. Furthermore, for $`\nu =3`$, conditions (36) and (40) hold for $`J=0.6`$ and $`A=12`$, whereas conditions (46) and (48) hold for $`J=0.05`$ and $`A=3`$. Therefore, monopoles move downstream for $`J0.05`$ and they move upstream for $`J0.6`$. Direct numerical simulations show that: (i) the estimate $`J_1=0.05`$ for the first critical current can be improved to $`J_1=0.08`$; and (ii) $`J_2=0.6`$ for the second critical current can be improved to $`J_2=0.54`$.
## IV Conclusions and final comments
We have presented a theory of monopoles moving downstream or upstream on an infinitely long doped, current biased superlattice when the fields are on the first plateau of the current–voltage characteristic. This theory has been corroborated with numerical evidence, which sharpens our results. Furthermore, we have simulated a 40-well 9nmGaAs/4nmAlAs SL under doping and contact conditions similar to experimental ones , but under constant current bias conditions. This situation is different from the usual case of voltage bias conditions. We have obtained that it is possible to observe monopole wavefronts moving upstream when the current is kept at large enough levels. Together with our theoretical bounds for critical currents and dopings, this numerical prediction could be used to set up an experiment to observe this striking phenomenon. For this purpose, we would need an initial condition corresponding to a monopole separating two electric field domains at high enough current. In an ideal world, this situation could be obtained by first fixing a low dc voltage for the 9/4 sample at a value near the top of one of the first branches of the current - voltage characteristics. Then we could switch from voltage to current bias conditions. The outcome would be a monopole moving upstream until the emitter region is reached. Presumably an idea of the field distribution corresponding to this situation could be obtained by time-resolved photoluminescence measurements .
There are technical problems that must be overcome if one wants to observe these features in real experiments: when we switch, there will always be a Faraday-like inductive pulse which will probably perturb the state of the system in an uncontrolled way. There are other possible biases we could think of. Under dc voltage bias, upstream moving monopoles are probably created for a short time during relocation experiment . In these experiments, one has a doped SL with a current-voltage characteristics correponding to multiple stationary monopole solution branches. Voltage is set at a particular value near the end of a branch, so that the field profile is that of a monopole layer connecting a low to a high field domain. Let the monopole layer be located at well $`i`$ (counted from the emitter contact). Then the voltage is suddenly and appropriately increased. After a certain time, the field profile settles to a new situation corresponding to a monopole layer centered at well $`i1`$ . This could be an indication of a monopole moving upstream, albeit for a short time. To increase this time, we could try to set a hybrid bias (between current and voltage bias) by including a finite series resistance in our external circuit. Additional theoretical and numerical work is needed to explore these possibilities.
## ACKNOWLEDGMENTS
We thank A. Amann, G. Platero and D. Sánchez for fruitful discussions and collaboration on the discrete drift-diffusion model. LLB thanks S.W. Teitsworth for a critical reading of the manuscript and helpful comments. This work was supported by the Spanish DGES through grant PB98-0142-C04-01, and by the European Union TMR contracts ERB FMRX-CT96-0033 and ERB FMRX-CT97-0157.
## A Comparison principle
The main theorem which we use to prove our results in Section III is the following comparison principle:
Theorem A.1 Let $`U_i(t)`$ and $`L_i(t)`$, $`iZ`$, be differentiable sequences such that
$`{\displaystyle \frac{dU_i}{dt}}d_1(U_i)[U_{i+1}U_i]`$ (A1)
$`d_2(U_i)[U_{i1}U_i]f(U_i)`$ (A2)
$`{\displaystyle \frac{dL_i}{dt}}d_1(L_i)[L_{i+1}L_i]`$ (A3)
$`d_2(L_i)[L_{i1}L_i]f(L_i),`$ (A4)
$`U_i(0)>L_i(0).`$ (A5)
where $`f`$, $`d_1>0`$ and $`d_2>0`$ are Lipschitz continuous functions. Then,
$$U_i(t)>L_i(t),t>0,iZ$$
In our discrete drift-diffusion model,
$`d_1(E)={\displaystyle \frac{D(E)}{\nu }},d_2(E)={\displaystyle \frac{D(E)+v(E)}{\nu }},`$ (A6)
$`f(E)=Jv(E).`$ (A7)
Proof: The proof is by contradiction. Set $`W_i(t)=U_i(t)L_i(t)`$. At $`t=0`$, $`W_i(0)>0`$ for all $`i`$. Let us assume that $`W_i`$ changes sign after a certain minimum time $`t_1>0`$, at some value of $`i`$, $`i=k`$. Thus $`W_k(t_1)=0`$ and $`dW_k/dt0`$, as $`tt_1`$. We shall show that this is contradictory. At $`t=t_1`$, there must be an index $`m`$ (equal or different from $`k`$) such that $`W_m(t_1)=0`$, while its next neighbor $`W_{m+j}(t_1)>0`$ ($`j`$ is either 1 or -1), and $`W_i(t_1)=0`$ for all indices between $`k`$ and $`m`$. For otherwise $`W_k`$ should be identically 0 for all $`k`$. Equation (A4) implies
$`{\displaystyle \frac{dW_m}{dt}}(t_1)d_1(U_m(t_1))W_{m+1}(t_1)`$ (A8)
$`+d_2(U_m(t_1))W_{m1}(t_1)>0.`$ (A9)
This contradicts the fact that $`dW_m/dt`$ should have been nonpositive as $`tt_1`$, for $`W_m(t_1)`$ to have become zero in the first place.
Corollary A.1 Any solution $`E_i(t)`$ of (17) with initial data $`E_i(0)(E^{(1)}(J),E^{(3)}(J))`$ satisfies $`E_i(t)(E^{(1)}(J),E^{(3)}(J))`$ for $`t>0`$.
Proof: Apply Theorem A.1 first with $`L_i=E^{(1)}(J)`$ and $`U_i=E_i`$, then with $`L_i=E_i`$ and $`U_i=E^{(3)}(J)`$.
Corollary A.2 If $`E_i(0)`$ is monotone increasing, that is, $`E_i(0)<E_{i+1}(0)`$, then, $`E_i(t)`$ is also monotone increasing, i.e., $`E_i(t)<E_{i+1}(t)`$ for $`t>0`$.
Proof: Apply Theorem A.1 with $`L_i=E_i(t)`$ and $`U_i=E_{i+1}(t)`$.
Remark. Strict inequalities in these theorems can be replaced by inequalities and the corresponding statements still hold. However the proofs become rather more technical and involved.
## B Bounds for critical doping values
We want to estimate the curves $`\nu _{1b}(J)`$ and $`\nu _{2b}^\pm (J)`$ defined in Section III. To estimate $`\nu _{2b}^+(J)`$, assume that $`J1`$ and $`\nu `$ is large. The left tail of a monopole is pinned if Eq. (25) holds. For large currents, (25) certainly holds if the curve corresponding to the left side of the inequality is above that of the right hand side, for $`E=1`$ (this is possible because $`D(E)`$ decreases rapidly to zero as the field increases). Setting $`E=1`$ in (25), we obtain
$$\frac{D(1)[1E^{(3)}(J)]}{\nu }>J1.$$
In turn, this implies $`\nu >\nu _{2b}^+(J)`$, defined in (32). This argument fails for the small values of $`J`$ used to draw Figure 7. We believe that quite different reasoning is needed to estimate $`\nu _{2b}^{}(J)`$.
The same argument yields our estimate $`\nu _{1b}(J)`$ of (31). For (30) to hold, the curve corresponding to the left side of the inequality should be below that of the right hand side for $`E=E_m`$. As $`D(E_m)0`$, we obtain
$$\frac{v(E_m)}{\nu }[E_mE^{(1)}(J)]<Jv(E_m),$$
which yields (31). Fig. 3 shows that the bound (31) is reasonably good for all eligible values of $`\nu `$ and $`J`$. |
warning/0003/cond-mat0003345.html | ar5iv | text | # Enhancement of GMR due to spin-mixing in magnetic multilayers with a superconducting contact
## Abstract
We study the Giant Magnetoresistance (GMR) ratio in magnetic multiayers with a single superconducting contact in the presence of spin-mixing processes. It has been recently shown that the GMR ratio of magnetic multilayers is strongly suppressed by the presence of a superconducting contact when spin-flipping is not allowed. In this Letter we demonstrate that the GMR ratio can be dramatically enhanced by spin-orbit interaction and/or non-collinear magnetic moments. The system is described using a tight-binding model with either $`s`$-$`p`$-$`d`$ or $`s`$-$`d`$ atomic orbitals per site.
PACS numbers: 75.70.Pa, 74.80.Dm
Hybrid nanostructures form a fascinating melting pot for studying the interplay between fundamental quantum phenomena, often revealing new and unexpected physics. One recently-recognized class of such structures, involving the coexistence of superconducting contacts and ferromagnetic domains, has led to the identification of a number of fundamental issues , several of which are currently unresolved. In this Letter, we examine one such issue, posed by experiments on giant magnetoresistance (GMR) in magnetic (M) multilayers with superconducting (S) contacts and current-perpendicular-to-the-plane (CPP). Recognizing that the sub-gap conductance of such structures is mediated by Andreev scattering, it was recently noted that in the absence of spin-flip processes, the conductance of a metallic (i.e. diffusive) multilayer in the presence of aligned magnetic moments is almost identical to that of the multilayer when adjacent moments are anti-aligned and therefore the conventional GMR ratio should be strongly suppressed. Since large GMR ratios are observed experimentally , it is clear that even a qualitative understanding of transport in such structures must incorporate the effects of spin-mixing. The aim of this Letter is to present the first theoretical description of CPP GMR in M-multilayers with S-contacts, which incorporate spin-flip scattering. As sources of spin-mixing we consider both spin-orbit (SO) coupling and non-collinear magnetizations in adjacent magnetic layers.
The system under consideration is a disordered magnetic multilayer consisting of an alternating sequence of magnetic layers each of length $`l_M`$ and non-magnetic layers (N) of length $`l_N`$. The building block of the magnetic structure is the bilayer \[M/N\] of length $`l_B=l_N+l_M`$. The magnetic moments of even-numbered M-layers make an angle $`\theta `$ relative to those of odd-numbered M-layers. Experimentally, $`\theta `$ can be varied by applying an external magnetic field with antiparallel (AP) alignment ($`\theta =\pi `$) typically occurring at zero field and parallel (P) alignment ($`\theta =0`$) at large enough fields. The current flows perpendicular to the planes of the multilayer, which makes contact with a metallic normal lead on the left-hand side of the multilayer and a superconducting lead on the right-hand side. GMR is the drastic increase in electrical conductance $`G(\theta )`$ that occurs when the system switches from the AP to the P alignment with the conventional GMR ratio defined by: $`\rho =\frac{G(0)G(\pi )}{G(\pi )}`$.
Following , the multilayer and leads are modelled using a tight-binding Hamiltonian on a cubic lattice with hoppings to nearest neighbours. Lattice imperfections and impurities are simulated by adding to the on-site energies a random number in the range $`[\frac{W}{2},+\frac{W}{2}]`$. The on-site Hamiltonian has the following structure:
$$H=\left(\begin{array}{cccc}H^p& \mu _{xy}& \mathrm{\Delta }& 0\\ \mu _{xy}^{}& H^p& 0& \mathrm{\Delta }\\ \mathrm{\Delta }^{}& 0& H^h& \mu _{xy}^{}\\ 0& \mathrm{\Delta }^{}& \mu _{xy}& H^h\end{array}\right)$$
(1)
where $`H^{p()}`$ is the Hamiltonian for up (down)-spin particles ($`s`$ and $`d`$ bands), $`H^{h()}=H^{p()}`$ is the Hamiltonian for up(down)-spin holes and $`\mathrm{\Delta }`$ is the superconducting order parameter. Here $`\mu _{xy}=\mu _x+i\mu _y`$, where $`\mu _{x(y)}`$ is the $`x(y)`$-component of the exchange field $`\stackrel{}{\mu }`$. Note that $`\stackrel{}{\mu }`$ is non-zero only for electrons in the $`d`$-band in the M-layers and $`\mathrm{\Delta }`$ is non-zero only in the right-hand-side superconducting lead. Within the tight-binding formulation, SO interaction can be included by adding to the Hamiltonian the following term:
$$V_{SO}=V_1\underset{i,j,\alpha ,s}{}\stackrel{}{\sigma }\stackrel{}{R}_{i,j}c_{\alpha ,i}^\sigma c_{\alpha ,j}^\sigma $$
(2)
where $`V_1`$ is a constant which determines the interaction strength, $`\stackrel{}{\sigma }`$ is a vector of Pauli matrices, $`\stackrel{}{R}_{i,j}`$ is the unit vector which connects site $`i`$ with the neighbouring site $`j`$. $`c_{\alpha ,i}^\sigma `$ is the annihilation operator for electrons of spin $`\sigma `$ in the $`\alpha `$ ($`s`$, $`d`$)-band on site $`i`$. In the presence of disorder (2) produces spin-flip scattering since it couples electrons with different spin on neighbouring sites.
In the presence of disorder, to study the largest possible sample cross sections, we consider 2 orbitals per site, which is the minimal model capable of reproducing scattering potential at the N/M interface and interband scattering . The tight-binding parameters are chosen to reproduce the GMR ratio and conductances obtained from an ab initio material specific calculation for Cu/Co multilayers .
In the presence of spin-flip scattering the two-spin-fluid approximation does not hold and the zero-temperature, zero-bias, normal-state Landauer formula takes the form:
$$G^{NN}=\frac{e^2}{h}\underset{\sigma \sigma ^{}}{}\text{Tr}\left\{t^{\sigma \sigma ^{}}t^{\sigma \sigma ^{}}\right\}$$
(3)
where $`t^{\sigma \sigma ^{}}`$ is the matrix of transmission amplitudes for injected $`\sigma ^{}`$-spin electrons in the left-hand lead into $`\sigma `$-spin electrons in the right-hand lead. When the right-hand lead is in the superconducting state, the conductance is given by
$$G^{NS}=\frac{e^2}{h}2\underset{\sigma \sigma ^{}}{}\text{Tr}\left\{r_a^{\sigma \sigma ^{}}r_a^{\sigma \sigma ^{}}\right\}$$
(4)
where $`r_a^{\sigma \sigma ^{}}`$ is the Andreev reflection matrix for injected $`\sigma ^{}`$-spin electrons in the left-hand lead to be reflected into $`\sigma `$-spin holes. In what follows, the scattering amplitudes are calculated exactly by solving the Bogoliubov-de Gennes equation using an efficient recursive Green’s function technique .
We shall now turn to the central results of this Letter, namely that in the presence of strong-enough spin-mixing, either produced by SO coupling or non-collinear moments, GMR in the presence of a S-contact approaches that of two normal contacts, with values of $`\rho `$ of the order of 100 %. First consider the effect of SO coupling within the multilayer. Fig. 1 shows the conventional GMR ratio as a function of the SO interaction strength $`V_1`$ for a disordered multilayer of 40 bilayers, with $`l_M=15`$ and $`l_N=8`$. As expected, in the NN case $`\rho `$ decreases monotonically from $`200`$ % at $`V_1=0`$, to zero at large $`V_1`$ ($`0.17`$ eV). In contrast for the NS case, $`\rho `$ initially increases with increasing $`V_1`$, eventually joining the NN curve at $`V_10.08`$ eV. As a second source of spin-mixing, consider the effect of non-collinear magnetic moments when $`V_1=0`$. Fig. 2 shows the $`\theta `$-dependence of the GMR ratio defined as $`\rho (\theta )=\frac{G(\theta )G(\pi )}{G(\pi )}`$. Whereas in the NN case $`\rho (\theta )`$ decreases monotonically with increasing $`\theta `$, in the NS case $`\rho (\theta )`$ exhibits a pronounced maximum around $`\theta =\pi /8`$.
To understand these results, first consider the case of non-collinear moments. In the presence of two normal-metallic contacts, the conductance $`G(\theta )`$ has been theoretically studied in where it is predicted that $`G(\theta )G(\pi )`$ tends monotonically to zero as $`\theta `$ varies from $`0`$ to $`\pi `$. In addition, the dependence of the resistance on the angle $`\theta `$ has been experimentally found to contain a term proportional to cos$`{}_{}{}^{2}(\theta /2)`$ and a second term proportional to cos$`{}_{}{}^{4}(\theta /2)`$. In the presence of a S-contact, where $`G(0)G(\pi )`$, this behaviour is drastically changed by the presence of an extremum which occurs at some intermediate angle $`\theta _c`$, the value of which depends on the interplay between competing effects. Since $`\theta (H)`$ is a function of the applied magnetic field $`H`$, the presence of a S-contact introduces a new characteristic field $`H_c`$ for which $`\theta (H_c)=\theta _c`$. For a disordered multilayer of 22 bilayers, with $`l_M=30`$, $`l_N=16`$, the insert in Fig. 2 shows the $`\theta `$-dependence of the conductance divided by the number of open channels in the normal lead. As expected $`G^{NN}(\theta )`$ is a monotonic function of $`\theta `$, whereas $`G^{NS}(\theta )`$ possesses an extremum at $`\theta _c\pi /8`$. To understand why the extremum is a maximum, recall that for $`\theta =0`$ or $`\theta =\pi `$, when spin is conserved, current flows when a right-going (spin $`\sigma `$) electron passes through the multilayer, Andreev reflects as a left-going (spin $`\sigma `$) hole, which retraverses the multilayer. A M-layer whose moment is aligned with the spin of the incident electron is anti-aligned with the spin of the outgoing hole and consequently the number of aligned and anti-aligned M-layers encountered by a given quasi-particle is the same for both $`\theta =0`$ and $`\theta =\pi `$ (only the order differs). When the elastic mean free path is comparable with the total multilayer length, the resistance of traversed layers add in series and therefore, apart from small differences due to interference effects , $`G^{NS}(0)G^{NS}(\pi )`$. Furthermore, since a quasi-particle must necessarily traverse regions in which it is a minority spin, both $`G^{NS}(0)`$ and $`G^{NS}(\pi )`$ are low-conductance states. In contrast, as $`\theta `$ increases from zero, this conductance bottleneck is removed, because an Andreev reflected minority hole can spin-convert to a majority hole, thereby avoiding anti-aligned moments on its return journey. Of course this initial increase in $`G^{NS}(\theta )`$ is eventually overcome by the usual GMR effect which decreases $`G^{NS}(\theta )`$ as $`\theta \pi `$, thereby producing an overall maximum.
In the absence of disorder, the nature of the extremum is determined by interface scattering and band structure. To illustrate this consider a clean multilayer which is perfectly periodic and therefore the variation of the conductance with $`\theta `$ arises from tuning of the ballistic spin-filtering by the structure. Fig. 3 shows the conductance divided by the number of open channels in the left-hand-side normal lead. As expected, $`G^{NN}(\theta )`$ is a monotonic function of $`\theta `$, whereas $`G^{NS}(\theta )`$ exhibits a minimum around $`\pi /2`$ and then increases. (In this case, translational invariance in the transverse direction allowed us to use a full ab initio, $`spd`$ Hamiltonian to obtain the results of Fig. 3.) In the NN case, the dependence of the multilayer resistance on $`\theta `$ predicted by our model is in good agreement with experiment . In Ref. the ratio between the resistance at a given $`\theta `$ and the resistance with AP alignment has been found to fit the following function:
$$\frac{R(\theta )}{R(\pi )}=1a\mathrm{cos}^2(\theta /2)+b\mathrm{cos}^4(\theta /2)$$
(5)
where $`a`$ and $`b`$ are fitting constants. In Fig. 4 we show the plot of such a ratio for the disordered multilayer considered above, along with the best fit to function (5). In addition we also checked that this ratio cannot be fitted with the same accuracy assuming a pure dependence on cos$`{}_{}{}^{2}(\theta /2)`$ (i.e. with $`b=0`$). For $`G^{NS}(\theta )`$ however, no such analytic results currently exist.
Let us now turn attention to the effect of SO coupling. Figs. 5a and 5b show the conductances as a function of the SO strength $`V_1`$ for, respectively, the NN and the NS case. In the NN case (Fig. 5a) $`G_P^{NN}`$ decreases as $`V_1`$ increases and eventually joins the curve for $`G_{AP}^{NN}`$. This can be understood in terms of the heuristic model presented in Ref. , because, as the SO strength increases, the average length required for a spin to flip (spin relaxation length $`\lambda _{sf}`$) gets shorter. Therefore in the P alignment, an injected majority electron travels through the multilayer for a length $`\lambda _{sf}`$ before being scattered into a minority spin, thereby producing a decrease in the conductance. This suggests that the value of $`V_1`$ for which $`G_P^{NN}G_{AP}^{NN}`$ corresponds to a spin relaxation length $`\lambda _{sf}`$ close to the period $`l_B`$ of the multilayer. We have carried out a range of simulations which show that this value of $`V_1`$ does not depend on the overall length of the multilayer, but decreases with increasing $`l_B`$. As expected, the conductance with AP alignment does not change significantly with $`V_1`$.
In the NS case (Fig. 5b) the conductance in the P aligned state rapidly increases with $`V_1`$, reaching a maximum and thereafter decreases, eventually joining the curve for the AP configuration. Clearly the enhancement in $`G_P^{NS}`$ is produced by the onset of spin-flip scattering. The abrupt increase is understandable, since even a relatively small probability for spin flipping opens a highly conductive “channel” if the spin-flip events take place in the vicinity of the interface. As one can see in the insert of Fig. 1 for larger values of $`V_1`$, the conductance $`G_P^{NS}`$ joins $`G_P^{NN}`$ and together they decrease thereafter. As in the normal case $`G_{AP}^{NS}`$ depends weakly on $`V_1`$. The value $`V_10.08`$, at which $`G^{NS}`$ is maximum, corresponds to a spin relaxation length close to the total length of the multilayer and, as expected, separate simulations show that this value of $`V_1`$ decreases with increasing total length. Similarly the value of $`V_1`$ at which the GMR ratio (of Fig. 1) vanishes corresponds to a spin-relaxation length of the order the bilayer thickness $`l_B`$ and is independent of the total length of the system.
In conclusion, we have demonstrated that spin-mixing plays a crucial rôle in determining both the qualitative and quantitative features of GMR in magnetic multilayers with a S-contact. In contrast with the normal case, where spin-mixing suppresses GMR, we find that the GMR ratio can be dramatically enhanced by the presence of spin-orbit interactions and/or non-collinear magnetic moments. In experiments carried out to-date, the presence of large spin-orbit scattering presumably masks the mechanism shown in Fig. 2, which is predicted to be a generic feature in the absence of other spin-mixing processes. This suggests that lighter metals and superconductors would be more appropriate for observing the new extrema predicted in this Letter. Finally we note that for the future it would be of interest to examine spin-mixing in non-diffusive NS structures such as clean spin-valves , where the GMR ratio can be non-zero or negative, even in the absence of spin-flip processes. |
warning/0003/cond-mat0003063.html | ar5iv | text | # Three-dimensional vortex dynamics in Bose-Einstein condensates
## Abstract
We simulate in the mean-field limit the effects of rotationally stirring a three-dimensional trapped Bose-Einstein condensate with a Gaussian laser beam. A single vortex cycling regime is found for a range of trap geometries, and is well described as coherent cycling between the ground and the first excited vortex states. The critical angular speed of stirring for vortex formation is quantitatively predicted by a simple model. We report preliminary results for the collisions of vortex lines, in which sections may be exchanged.
The formation and properties of vortices in Bose-Einstein condensates are attracting intense current interest. A number of schemes for producing vortices in trapped condensates have been proposed theoretically, and the first experimental creation of a topological vortex in a two-component condensate has been reported , in agreement with a numerical model . Simulations of flow of a condensate past a spherical object have shown vortex ring solutions , and numerical studies of vortex formation in rotating traps have shown the formation of vortex arrays . Very recently the first report has been made of the observation of vortices in a single-component condensate .
In an earlier paper , we showed that a rotationally stirred Bose-Einstein condensate, simulated in two dimensions, exhibits a simple single vortex cycling regime, a behaviour we described as nonlinear Rabi cycling. In the present paper we demonstrate that the single vortex cycling regime is also found in full three-dimensional simulations of condensate stirring (as shown in Fig. 1), in a parameter regime accessible to current experiments. The single vortex regime is studied in a variety of trap geometries (oblate to prolate) as is the dynamical stability of central vortices. The line character of a vortex in three dimensions also gives rise to some qualitatively new behaviour that has no analogue in two dimensions. We show that in a spherical trap, stirring of the condensate can produce two vortex lines which collide and exchange sections.
As in the earlier work, our treatment is based on the time-dependent Gross-Pitaevskii (GP) equation,
$$i\frac{\psi (𝐫,t)}{t}=\left[^2+V(𝐫,t)+C|\psi (𝐫,t)|^2\right]\psi (𝐫,t),$$
(1)
in dimensionless units . The total external potential $`V(𝐫,t)`$ is given by $`(x^2+y^2+\lambda ^2z^2)/4+W(𝐫,t)`$, where the first term represents a cylindrically symmetric trap with trap anisotropy parameter $`\lambda `$. The second term $`W(𝐫,t)`$ is the contribution of the stirrer, a far-blue detuned Gaussian laser beam cylindrically symmetric about a line parallel to the $`z`$ axis,
$`W(𝐫,t)=W_s(t)\mathrm{exp}\left[\left({\displaystyle \frac{|𝝆𝝆_s(t)|}{w_s/2}}\right)^2\right],`$ (2)
where $`𝝆`$ is the projection of $`𝐫`$ into the $`z=0`$ plane, and $`𝝆_s`$ is the centre of the Gaussian stirrer in that plane.
The condensate is stirred in a manner similar to our earlier work , with the stirrer moving in a circle of radius $`\rho _s`$ at a constant angular velocity $`\omega _f`$. Here, however, we begin with the ground state of the harmonic trap (with no stirrer present), and minimise transient effects by increasing the stirrer amplitude $`W_s`$ linearly from zero to a final value of $`W_0`$ between $`t=0`$ and $`t=\pi `$. With an appropriate choice of parameters, as in Fig. 1, a single vortex with positive circulation is drawn close to the centre of the condensate, and for continuous stirring cycles in and out. This is exactly the behaviour found previously in a two-dimensional system , and in fact the parameter choice for the single vortex cycling behaviour in three dimensions is guided by the same considerations outlined there, as we discuss below. Figure 2, where we plot the phase of the wavefunction, shows that the central feature of Fig. 1 is a vortex of positive circulation. In Fig. 3 the evolution of the angular momentum expectation value (which in our dimensionless units is 0 for the ground state and 1 for the central vortex state) is presented and confirms that for continuous stirring (thin line), the vortex cycles in and out of the condensate.
In contrast to the two-dimensional case, the three-dimensional simulation can be unambiguously related to a realistic current experimental scenario. The $`C=1000`$ case shown in Figs. 13 corresponds to $`N=1.8\times 10^4`$ atoms of <sup>87</sup>Rb in the $`|F=2,m_f=2`$ hyperfine state (for which we use an s-wave scattering length $`a=5.29`$ nm) in a time orbiting potential (TOP) trap with radial trap frequency $`\omega _r=2\pi \times 15`$ Hz and $`\lambda =\omega _z/\omega _r=\sqrt{8}`$. For this example, our computational units of time and length correspond to $`10.6`$ ms and $`1.97`$ $`\mu `$m respectively, and the Thomas-Fermi diameter of the condensate in the radial direction of 11.7 corresponds to about 23 $`\mu `$m. The stirrer Gaussian $`1/e`$ diameter $`w_s=4`$ corresponds to a Gaussian $`1/e^2`$ beam diameter of approximately 11 $`\mu `$m. The required intensity of the stirring laser is detuning dependent, but can be easily calculated from the light shift potential $`\mathrm{}\mathrm{\Omega }^2/4\mathrm{\Delta }`$, where $`\mathrm{\Omega }`$ and $`\mathrm{\Delta }`$ are the Rabi frequency and detuning respectively of the atom-laser interaction. For example, if the laser beam is 50 nm blue-detuned from the 780 nm atomic transition, $`\mathrm{\Delta }=1.65\times 10^{14}`$ s<sup>-1</sup>, and a laser power of 300 nW is required for the 3 nK height stirrer we describe. The first angular momentum peak in Fig. 3 occurs at 0.26 s, which is much less than the lifetime of the condensate. The occurrence of the single vortex cycling behaviour in this experimentally accessible parameter regime suggests a method of preparing a central vortex state; stir the condensate as we have described, then remove the stirrer at the time when the angular momentum peak is expected. We illustrate the implementation of this method in Fig. 3 with the thick curve, for which the stirrer is linearly withdrawn between $`t=7\pi `$ and $`t=8\pi `$. Once the stirrer is removed, the vortex remains in the condensate, circling about and close to the centre, and $`L_z`$ is constant.
The vortex cycling behaviour illustrated by the example in Figs. 13 indicates that the two-state model described in our earlier work (for the two-dimensional case ) also has application in three-dimensions. Here we briefly outline the modifications required to extend the two-state model into three dimensions. We represent the condensate by an axisymmetric component $`\varphi _s`$, and a vortex-like component $`\varphi _ve^{i\theta }`$ which has a circulation about the $`z`$ axis, where $`\varphi _s`$ and $`\varphi _v`$ are real and nonnegative, so
$$\psi (𝐫,t)=a_s(t)\varphi _s(\rho ,z,n_v)+a_v(t)\varphi _v(\rho ,z,n_v)e^{i\theta },$$
(3)
with $`\rho =\sqrt{x^2+y^2}`$. We choose $`\varphi _s`$ and $`\varphi _v`$ to satisfy the following time-independent coupled equations,
$`\mu _s\varphi _s`$ $`=`$ $`[{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{d}{d\rho }}\left(\rho {\displaystyle \frac{d}{d\rho }}\right){\displaystyle \frac{d^2}{dz^2}}`$ (6)
$`+{\displaystyle \frac{\rho ^2+\lambda ^2z^2}{4}}+C(n_s\varphi _s^2+2n_v\varphi _v^2)]\varphi _s,`$
$`\mu _v\varphi _v`$ $`=`$ $`[{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{d}{d\rho }}\left(\rho {\displaystyle \frac{d}{d\rho }}\right){\displaystyle \frac{d^2}{dz^2}}+{\displaystyle \frac{1}{\rho ^2}}`$ (8)
$`+{\displaystyle \frac{\rho ^2+\lambda ^2z^2}{4}}+C(n_v\varphi _v^2+2n_s\varphi _s^2)]\varphi _v,`$
where the eigenvalues $`\mu _s`$ and $`\mu _v`$ are found by solving Eqs. (8) for a given choice of the vortex fraction $`n_v=|a_v|^2`$. The quantity $`n_s=|a_s|^2=1n_v`$. We note that when $`n_v=0`$, Eqs. (3) reduce to Eq. (6), the time-independent GP equation for the ground state, while if $`n_v=1`$, Eqs. (3) reduce to Eq. (8), the time-independent GP equation for the first excited vortex. The resulting nonlinear Rabi equations for the two-state system \[obtained by substituting Eq. (3) into Eq. (1)\] are then identical in form to Eqs. (5–6) of , but with the wavefunctions $`\varphi _s`$ and $`\varphi _v`$ now defined over $`\rho `$ and $`z`$ instead of $`r`$.
The underlying assumption of our two-state model is that the ground state (or more generally, $`\varphi _s`$) is coupled primarily to the vortex state of lowest energy with axis parallel to the stirrer axis. In two dimensions, there are a limited number of excited condensate states in the appropriate energy range that can be coupled to the symmetric state by the stirring potential, and we have previously shown that the two state model works well within a specified validity range. In three-dimensional cylindrically symmetric traps, many more condensate eigenstates exist, and although symmetry considerations limit the possible coupling from the symmetric state, a two state model places more constraint on the description of the system than in the two-dimensional case. We do not aim therefore to provide a detailed representation of the stirring behaviour with our two state model, instead we use it primarily to provide conceptual understanding of the cycling behaviour. Nevertheless we also find it gives a quantitatively accurate prediction of the critical frequency of rotation $`\omega _c`$, the stirring frequency which causes a single vortex to cycle right to the centre of the condensate. We have examined the single vortex cycling behaviour using the full GP solution, for a variety of trap geometries including the oblate ($`\lambda =\sqrt{8}`$), spherical ($`\lambda =1`$), and prolate ($`\lambda =1/3`$) cases, with a condensate of $`C=1000`$, and a stirring potential of the same size and position as in Fig. 1. We find that the critical rotational frequency for single vortex cycling is accurately predicted by the two state model as the difference in energy between the ground and vortex states of the condensate (modified by the finite size of the stirrer, see Eq. (7) of ). As in , a sufficiently large Rabi frequency is required to distort the energy barrier between ground and vortex state enough to permit cycling. When the critical frequency of rotation is exceeded, additional vortices penetrate the condensate.
It is worth noting that as the $`z`$ component of the trapping potential become weaker (that is, as $`\lambda `$ becomes smaller), then in addition to the vortex state, the stirring potential increasingly excites the dipole (centre of mass) oscillation of the condensate. This can be understood in terms of the Ehrenfest theorem: for our choice of parameters the mean force $`W`$ of the stirrer increases as the condensate radial size decreases because the condensate is increasingly concentrated near the steepest part of the stirring potential. The centre of mass motion contributes angular momentum $`L_z_{COM}`$ to the condensate angular momentum, where
$$L_z_{COM}=xP_yyP_x.$$
(9)
We find that $`L_zL_z_{COM}`$ (which is a measure of the angular momentum caused by the presence of the vortex) for all three trap geometries cycles between 0 and 1 when $`\omega _f`$ is slightly less than $`\omega _c`$ (the single vortex cycling regime) and oscillates about 1 when $`\omega _f`$ just exceeds $`\omega _c`$ (and multiple vortices penetrate the condensate). $`L_z_{COM}`$ is small for the $`\lambda =\sqrt{8}`$ case shown in Fig. 3.
If the single stirrer used in Fig. 1 is replaced by two stirrers of half the height but on opposite sides of the $`z`$ axis, no vortex cycling occurs, and $`L_z`$ exhibits only small oscillations of amplitude less than $`0.04`$. This is easily understood from our two state model, because the Rabi frequency $`\mathrm{\Omega }`$ (see Eq. (6b) of ) is zero for this stirrer configuration, as it is for any stirrer configuration of even symmetry. It is worth pointing out that our scenario of coherent cycling should be distinguished from vortex production by damping in a trap with a symmetric rotating distortion, such as the recently reported experiment , or numerical damping schemes such as imaginary-time propagation .
Our three-dimensional simulations also allow us to test (in the mean-field limit) the stability of pure vortex eigenstates. Using as an initial state a central vortex (an eigenstate of the time-independent GP equation) for $`C=1000`$, a perturbation was applied by linearly inserting and withdrawing (between $`t=0`$ and $`t=\pi `$) a stirrer of the same size as in Fig. 1 but at a fixed position in the laboratory frame. We find that for a range of trap geometries ($`\lambda =\sqrt{8}`$, $`1`$, or $`1/3`$), the $`L_z=1`$ vortex is dynamically stable, and remains near the centre of the condensate. We can understand this result in terms of Rabi cycling. The stationary stirrer ($`\omega _f=0`$) provides a potential which is far-detuned from resonance ($`\omega _f=\omega _c`$), and thus there is a low probability of transfer out of the initial (vortex) state. Indeed if the stirrer remains fixed in the condensate after insertion and is not withdrawn, we find that the vortex still remains near the centre of the condensate and follows a small closed path displaced slightly away from the stirrer. By contrast, for these geometries, the $`L_z=2`$ central vortex eigenstate immediately dissociates into two unit vortices when perturbed. We note that these stability results refer to dynamical stability in the GP limit (that is, $`T=0`$), as opposed to thermodynamic stability .
In addition to confirming the basic features of the two state model of vortex cycling, our three-dimensional simulations also reveal some new phenomena. The vortices produced in three dimensions by stirring are line vortices. In the $`\lambda =\sqrt{8}`$ configuration, the lines are nearly straight and parallel to the stirrer axis, and when two lines of the same circulation form, we find they repel each other. For a spherical trap ($`\lambda =1`$) however, vortex lines along any direction are degenerate (in the absence of the stirrer) with the vortex eigenstate about the $`z`$ axis. Accordingly, the vortex lines can become more curved than in the oblate case, as we find in our simulations below the critical frequency, where the vortex that forms can have appreciable curvature. Above the critical frequency, when multiple vortices are drawn into the condensate, there is a striking consequence of this potential for vortex lines to deform. In a collision between two such curved vortex lines, sections of the lines can be exchanged, as we show in the collision sequence in Fig. 4. The condensate is stirred somewhat above the critical frequency in a spherical trap, and by $`t=6.79`$ two vortex lines (both of positive sense ) have formed. The figure shows only the vortex lines (detected numerically), which have been shaded light and dark to distinguish one from the other. In Fig. 4(a) the dark vortex line approaches the light vortex line from behind, and because of their mutual repulsion, the light line has developed a bulge. In Fig. 4(b) the dark line is now close to the light line, causing the bulge to become recurved, with two small portions now nearly antiparallel to the dark line. In Fig. 4(c) the two lines have completed a collision and exchanged their central sections.
In conclusion, we have shown that the single vortex cycling regime we discovered previously in simulations of two-dimensional condensates is also present in three-dimensional condensates, for a range of trap geometries. We have related our simulations to current experimental configurations, and have demonstrated how the cycling behaviour could be used to generate a central vortex state. We have investigated the dynamical stability of central vortex states, and have also given initial evidence of the rich dynamical behaviour of vortex lines.
This work was supported by the Marsden Fund of New Zealand under contract PVT902. |
warning/0003/astro-ph0003135.html | ar5iv | text | # Extended X-ray emission in the radio loud galaxy 3C382
## 1 Introduction
Soft X-ray excess emission above a simple extrapolation of the hard energy spectrum is found in a considerable number of AGNs, mainly radio-quiet sources (see Mushotzky, Done, & Pounds 1993 and references therein). There is a growing body of evidence for its spectral ubiquity below $``$2 keV or so in Seyferts (e.g. Pounds et al. 1994) as in highly luminous quasars (e.g. Saxton et al. 1993). Most previous studies before ROSAT converge toward the idea that the soft excess is a rather common feature among radio-quiet quasars, whereas it is almost absent in their radio-loud counterparts. ROSAT/PSPC data of radio loud sources have shown that a soft excess component is also present in radio loud sources (e.g. Buehler et al. 1995; Prieto 1996; Siebert et al. 1998).
Possible interpretations for the soft X-ray excess in AGNs include thermal emission from the inner regions of an accretion disk, scattering by highly ionized material in its vicinity (Pounds et al. 1986; Ross & Fabian 1993), or thermal emission due to shock-heated gas in the close vicinity of the nucleus (Viegas & Contini 1994). The poor spatial resolution of the ROSAT/PSPC makes difficult the separation between possible components of the observed emission. Indeed, the large PSPC resolution beam, $``$ 25 arcsec at 1 keV, makes plausible that an important part of the observed emission to be due to an extended gas component surrounding the AGN. In the particular case of radio-loud galaxies which are characterized by large radio sizes, a hot surrounding medium becomes a necessary component for providing the working surface for the radio emission. In the analysis of the 3CRR sample by Prieto, a first attempt to fit the PSPC spectra of sources with extended emission –mostly in Fanaroff & Riley (1974) type I sources (FRI)– with a single power-law leaded to extreme step spectral index, the reason being due to the dominant contribution of the gaseous medium in which those sources usually reside. In the case of FRII, a single power-law fit provided a fair representation of the PSPC spectrum but with average spectral index about -1.1, and so above the extrapolation of the canonical hard-energy spectrum into the soft X-rays. Clustering of galaxies about FRII sources is less common than in FRI, in particular at low redshift; yet, FRII could contain their own extended gaseous atmosphere which may directly translate into a steepening of the PSPC spectrum. This component however may prove to be elusive with present X-ray instrumentation.
This paper presents deep ROSAT/HRI observations of the powerful X-ray radio-loud source 3C 382. This is one of the few nearby broad-line galaxies ($`z=`$0.0578) that show extremely bright and broad permitted lines (FWZI$`>25000kms1`$; Tadhunter, Perez & Fosbury 1986) and a strong continuum, with with a X-ray luminosity in 0.2-2.4 keV band of $`L_\mathrm{x}7.10^{44}`$ erg s<sup>-1</sup> (Prieto 1996), and a radio power at 178 MHz of $`L_{178\mathrm{M}\mathrm{H}\mathrm{z}}3.10^{33}`$ erg s$`{}_{}{}^{1}Hz^1`$ (Laing et al. 1983).
EXOSAT monitoring of the source (1983–1985) tightly constrains the high-energy (above 2 keV) spectral index of 3C382 to $`\alpha =0.7\pm 0.1`$ (Ghosh & Soundararajaperumal 1992). This is also confirmed by more recent ASCA data (Wozniak et al. 1998). However, the ROSAT/PSPC spectral analysis of the source shows compatible with a power-law model with spectral index $`\alpha =1.2\pm 0.3`$, and absorbed by a column density, $`N(H)=0.78\times 10^{21}cm^2`$, that is in agreement with the Galactic value. Thus, 3C382 shows a soft excess emission below $``$ 2keV (Prieto 1996). Independently, the presence of a soft excess is also inferred from the analysis of the EXOSAT (Ghosh & Soundararajaperumal), ASCA (Wozniak et al.) and Ginga (Kaastra et al 1991) data.
Extended soft X-ray emission associated with this source is detected in the ROSAT/HRI data. On the basis of that new component a re-evaluation of the PSPC spectrum is presented.
Throughout this work $`H_0`$ = 50 km s<sup>-1</sup> Mpc<sup>-1</sup>. 1 arcsecond corresponds to $``$ 1.7 kpc at the source.
## 2 Analysis of HRI data
The HRI observations of 3C 382 were conducted in 1996 October and 1997 April (WG900720H and WG900720H-1 datasets respectively) The corresponding total accepted times were 4514 s and 13310 s, respectively. The counts are integrated from channels 1 to 8 which enclose most of the energy accumulated in the ROSAT band, yielding for both dataset count rates of $``$1 cts s<sup>-1</sup>.
The nominal resolution of the ROSAT/HRI is $``$ 5–6 arcsec (FWHM). Residual errors in the ROSAT aspect solution are known to give rise to elongated images (David et al. 1996), the shape of the surface brightness profile dramatically departing from the expected point response function (PRF). In the case of very bright sources, improvement of the HRI spatial resolution becomes feasible by using speckle interferometric techniques such as the “shift-and-add” method. If one constructs images over short time intervals, the PRF becomes symmetric and therefore the elongation in the image appears as an apparent residual motion of the X-ray source in the sky. Such residual motion can be corrected for by de-speckling.
3C 382 is particularly suitable for that technique as it shows very bright in the ROSAT band, with $``$ 2 cts s<sup>-1</sup> in the PSPC and $`1`$ cts s<sup>-1</sup> in the HRI.
The procedure used follows the same criteria and approach originally presented in Schmitt, Güdel & Predehl (1994). Basically, each event file is divided into time bins of typically 50 s. For each bin, the apparent X-ray position of the source (i.e., RA and $`\delta `$) is determined as a function of the observing time. A spline function is fitted to these data points, all recorded photons being then corrected with the appropriate time-dependent correction in RA and $`\delta `$.
To validate the correction, new measurements of the source centroid are repeated on the corrected event file. The correction is considered as satisfactory if the new centroid positions cluster during the time period of the observation about an average constant value. The uncertainty in the final source centroid is about 2.5 arcsec in RA and 1.5 arcsec in $`\delta `$.
The X-ray spatial analysis was then performed on the corrected HRI event files. Because of the much higher statistic of the April event file, reliable results from the de-speckle procedure are only found from that dataset. Thus, the following analysis focuses only on this dataset.
## 3 The extended X-ray component
An X-ray contour image of 3C 382 obtained after the de-speckle procedure is shown in Figure 1a. The image is background-subtracted and smoothed with a Gaussian filter of FWHM $``$12 arcsec. The background level is estimated from different regions in the image within the central 5 arcmin. The emission is dominated by a central peak component, and a surrounding halo slightly more asymmetric towards the North-East side. Some of the morphological features are also apparent in the shortest exposed 1996 October event file -not shown. Overall, the X-ray emission extends out to about 100 arcsec from the center, $``$170 kpc at the distance of the source. This is about the size of the radio structure at 8.3 GHz (Black et al. 1992). The slight North-East asymmetry in the X-ray image is virtually coinciding with the direction of the radio axis and with that of the faint filamentary regions and bright companion galaxy seen in the HST/WFPC2 image (Fig. 1b).
The corresponding surface-brightness profile is presented in Figures 2. For comparison, the new improved ROSAT/HRI PSF profile by Predhel (1998) —currently upgrade in the EXSAS package— which includes the effect of the ROSAT mirror scattering and has been fit to Capella and Sirius, is shown superimposed. To illustrate the effect of the de-speckle procedure on the data, the surface brightness profile as derived from the original data is shown is Fig 2a, and that after de-speckling in Fig. 2b. The departure from the ROSAT/HRI PRF is clearly evidenced in Fig. 2a, where larger residuals are seen all over the profile. The de-speckle procedure produced a sharper profile which virtually fills the core of the PRF (Fig. 2b). Yet, some systematic residuals still remain at radius beyond $``$10–15 arcsec from the center, which we assume to be related to an extended emission component. This component is roughly 10% of the total emission.
In an attempt to get a better characterization of the observed profile, a combination of an unresolved component represented by the HRI/PSF, and an extended component represented by a $`\beta `$-model are fit to the data. The functional form of the $`\beta `$-model follows King’s approximation (1972) for gas confined in an isothermal sphere: $`S(r)(1+(r/r_c)^2)^{3\beta +1/2}`$, with S(r) the surface brightness a radius r, $`r_c`$ the core radius and $`\beta `$ the slope parameter. Given the uncertainty inherent with this approximation for the extended component, a simple addition of the PSF and the beta model is fit to the data.
Figure 2c shows the resulting fit corresponding to the composite model: PRF plus a $`\beta `$ model. Fig 2c shows a clear fit improvement to the observed profile as the systematic residual trend seen in Fig 2b is now removed; yet, the composite model still overpredicts the observed emission at radius beyond 10 arcsec. Overall, the $`\beta `$-model does not provide a statistically acceptable fit with reduced $`\chi ^2`$ exceeding unity. The main reason for that may reside in the validity of the beta-model for the specific case of the extended gas emission in 3C382. Besides, there is the fact that the dominant contribution of the emission is the unresolved component, which makes the modeling of the extended diffuse emission difficult. Also, there is the rather asymmetric morphology of the emission which may clearly depart from the simple model used here.
Nevertheless, the fit results from the composite model leaded to a tight range for both the core radius and the $`\beta `$ value. The minimum reduced $`\chi ^2`$, $`<1.52>`$, was obtained for a core radius $`r_c`$ 20-30 arcsec (30 - 50 kpc at the distance of 3C 382) and $`\beta 0.70.8`$. A much simpler composite model represented by the PRF plus a Gaussian model for the extended component, yielded FWHM for the extended component in the 30 - 40 arcsec range.
### 3.1 Re-evaluating the PSPC spectrum
A single power-law model with spectral index -1.2 provides a fair fit of the PSPC spectrum of 3C382, with the derived N(H) in good agreement with the Galactic value. However, the detection of an extended component in the HRI data along with the reported evidence for soft excess emission in this source from EXOSAT and ASCA data besides PSPC (cf. sect. 1) prompt to a re-evaluation of the PSPC data. Because the extended nature of the new component, its association with thermal emission arises as the most natural explanation.
Due to the short energy range and relatively low spectral resolution of the PSPC, a complex analysis of the PSPC spectrum involving a combination of several models is difficult. However, the evidence for extended emission and the fact that the hard X-ray spectral index of this source remains constant about the 0.7 value are additional inputs that can be used for forcing the PSPC fit with a more complex modeling.
Accordingly, a combination of a power-law, with spectral index fixed at $`\alpha `$= –0.7, and a thermal model are fitted to the PSPC data. In this composite model, $`N`$(H) is fixed at the Galactic value following Prieto’s results (1996). Thus, the only parameters that vary freely are the normalizations of the respective thermal and power-law components and the temperature of the thermal component.
The new composite fit (Fig. 3) provides a fair representation of the PSPC spectrum, with a reduced $`\chi ^2`$1.2 and constrained values for all the free parameters. A gas temperature of $`0.6_{0.1}^{+0.4}`$ keV is derived. Alternative fits, letting free the index of the power-law spectrum or the N(H), lead to larger $`\chi ^2`$ and unconstrained fit parameters. Fitting a single bremsstrahlung model produces also an acceptable fit but the derived N(H) value is found lower than the Galactic value.
The results from the composite model can be compared with those derived from the single power-law model and from the single thermal model in Table 1. Given the PSPC spectral limitations together with the additional complexity of the composite model, the derived fit values may be subjected to larger uncertainties. In this sense, the temperature could be largely affected whereas integrated fluxes are the least dependent on the adopted model. The total flux contribution from the thermal and power-law components in the composite model compares with that derived from the single power-law model. The reduced $`\chi ^2`$ is slighter better in the case of a simple power-law model; but we consider the difference marginal. The lowest reduced $`\chi ^2`$ is obtained with the thermal model; yet the derived N(H) in this case is inferior than the Galactic value.
### 3.2 How do the PSPC fluxes compare with those derived from the HRI data?
Fluxes estimate from the HRI data for the extended and unresolved component are primarily derived on the basis of HRI spatial analysis (§3). The total number of counts within the unresolved and extended component leads to count rates of $``$ 1 cts s<sup>-1</sup> and about 0.1 cts s<sup>-1</sup> respectively.
The lack of spectral resolution of the HRI hampers any direct spectral modeling of the X-ray data. Yet, the availability of the PSPC data allows us to use the same model and fit parameters as derived from the PSPC fit (§3.1) to convert the HRI counts to fluxes. These are given in Table 1
The unresolved HRI component, modeled with a power-law index $`\alpha `$= –0.7, yields a flux about a factor 2 larger than that measured by the PSPC. 3C 382 is known to be a variable source, with reported maximum-to-minimum variations of up to 120% as measured by EXOSAT (Ghosh & Soundararajaperumal 1992). Regarding the ROSAT observations, previous HRI observations of the source in 1992 March revealed a drop in the total number of counts by a factor 1.8 with respect to the values measured in the present observations taken 5 years later. This drop is thus compatible with the still lower flux measured by the PSPC in 1990. Thus, the difference found between the HRI and PSPC fluxes for the unresolved component appears compatible with the observed variability level of the source.
Regarding the extended component, bremsstrahlung models with temperatures within the range 0.6-3 keV lead to corresponding HRI fluxes $`106.\times 10^{12}`$ erg cm$`{}_{}{}^{2}s_{}^{1}`$ respectively, i.e., a factor 2 to 3.5 smaller than that derived from the PSPC for the same component. We note however that as neither the Gaussian approximation nor the $`\beta `$-model provide a statically acceptable fit to that component, the corresponding number of HRI counts for the extended component was not derived from the model but from the difference between the total integrated number of counts and that integrated within the unresolved component represented by the PRF. Given the simplicity of the method, in particular taking into account that the dominant contribution of the total emission is the unresolved component, the derived PSPC and HRI fluxes can be considered consistent with each other within the order of magnitude. Besides, it is also possible that the emission is more extended than what is detected –at least beyond the radio structure–, but the surface brightness there is too low for being detected with the HRI.
To summarize, the inferred luminosities for the extended and unresolved components estimated from the HRI data appear compatible in order of magnitude with the respective thermal and power-law luminosities derived from PSPC spectrum of the source.
## 4 Discussion
Extended X-ray emission associated with the very bright, nearby radio-loud galaxy 3C 382 is detected. The analysis of the ROSAT/HRI data shows that about 10% of the total 0.2 -2.4 keV emission is compatible with the presence of an extended component. Assuming that component to be due to hot gas emitting via bremsstrahlung, and allowing for the contribution of such thermal component into the PSPC spectrum, it is found that the non-thermal component of 3C 382 emission becomes consistent with the extrapolation of the “well established” 3C 382 high-energy power-law spectrum —above 2 keV— into the soft X-ray regime.
The temperature of the gas component as formally derived from the PSPC fit is $`0.6_{0.1}^{+0.4}`$ keV. This low temperature contrasts with the high luminosity of the gas component, $`3\times 10^{44}ergsec^1`$. Taken together the spectral limitation of the PSPC and the complexity of the model fit, the uncertainty in the temperature could however be larger. There is furthermore the possibility that a temperature gradient dominates the gas emission — a cooling flow process. In this case, the PSPC spectrum may be mostly sampling the central gas region, where the coolest and more dense gas is located. There are however other factors that could also have lead to an overestimation of the gas luminosity. The analysis by Markevitch (1998) on clusters with strong cooling flows indicates moderate temperature increase of up to 20% but luminosity decrease of up to 40% for the cluster gas after excising the cooling flow regions. On the other hand, if cooling by metals is considered – for sake of simplicity, pure bremsstrahlung has been assumed – the luminosity of the gas could decrease by about 40%, assuming a Raymond-Smith model with metal abundance Z=0.35.
Still, the derived thermal luminosity in 3C382 is about two order of magnitude larger than that found in isolated, normal elliptical galaxies (Canizares, Fabbiano & Trincheri 1987), and in low-power radio galaxies (Worral & Birkinshaw, 1994) but it is in the range found in powerful radio sources (Worrall et al. 1994; O’Dea et al. 1996; Hardcastle et al. 1999; Crawford et al. 1999). Also, the estimated core radius for 3C382, $``$50 kpc, is within the range found by Crawford et al. and Hardcastle et al. in their respective samples of 3CRR radio-loud sources.
Large X-ray halos are often seen in FRI sources, those being associated with the cluster environment in which they often reside (e.g. M 87, Perseus, 3C465). These halos largely dominate the ROSAT emission from these sources. Besides the outstanding case of Cygnus A, evidence for clustering is less obvious in classical double radio sources, particularly at low redshift (cf. Hill & Lilly 1991; Miller et al. 1999). Unambiguous extended X-ray emission in powerful FRII radio galaxies and quasars has mostly being found in sources with redshift larger than 0.1 (Hardcastle and Worall, 1999; Crawford et al 1999; O’Dea et al. 1996); yet, a few low-redshift FRII are reported to show extended X-ray emission (cf. Hardcastle and Worral). In most of these cases, the large X-ray luminosities are found compatible with thermal emission from a moderately rich cluster environment.
Comparing with Cyg A, the archetypal double radio source at z= 0.0574, 3C382 is also one of the few very bright double sources at low redshift with extended X-ray emission. Contrarily to Cyg A which presents an optical narrow line spectrum, 3C382 presents an extreme, in width and strength, broad permitted line spectrum. If this difference is interpreted as due to obscuration of the AGN region in Cyg A, it may explain why the dominant X-ray feature in Cyg A is emission from a hot diffuse gas –the AGN component is obscured at X-ray waves– whereas in 3C382, the unresolved X-ray nuclear component –presumably associated with the AGN– dominates the total X-ray emission, making more difficult the detection of any extended gas component.
The X-ray luminosity of 3C382, of about $`10^{44}ergs^1`$, compares with that of rich Abell clusters (this is also the case of Cyg A; yet Cyg A is at the center of a poor cluster of galaxies). Longair & Seldner (1979) derived however a rather poor environment in the vicinity of 3C 382 on the basis of their cross-correlation analysis between the radio position and galaxy counts. HST/WFPC2 images of 3C382 collected in parallel mode show an elliptical galaxy with a very bright unresolved nucleus and a halo very smooth (Martel et al. 1999). Yet, within the 2.5 arcminutes field of view (Fig 1b), several small galaxies can easily be distinguished in the 300 seconds exposure; the WFPC2 images also show a bright galaxy at 85 arcsec Northeast from 3C382, presenting two at least extended gaseosus tails of material in the direction of 3C382; two additional difuse regions located close to 3C382 and in the direction of the bright galaxy are also apparent. Judging from the HST images, 3C382 may be residing in a relatively poor cluster environment; also, it may be in interaction with that gas-rich galaxy companion. Such interaction could have brought plenty of gas into 3C382.
The luminosity of the halo component in 3C382 would imply a large mass of gas, of about $`10^{11}`$Mo, assuming it concentrated in a sphere of about 50 kpc –the estimate core radius derived from the HRI spatial analysis– and a temperature in the 0.6 -1 keV range. An alternative to the cluster environment is that 3C382 it may consist of a self-contained gravitational potential deep enough to restrain such large amount of gas. This could also be the case of the radio-louds 3C48 and 3C273, for which extended X-ray emission is found but the evidence for a cluster environment from optical images is minor (Crawford et al. 1999). Evidence for a massive dark halo in 3C382 comes from the velocity measurements on the extended ionized gas surrounding this galaxy. Tadhunter et al. (1986) detected ionized gas up to 25 kpc from the galaxy center. The gas follow a a rotation curve which extend flat up to those distances with velocities of about 400 km/s relative to the systemic velocity. Assuming a spherical potential, the estimated mass within a 25 kpc radius would be $`8\times 10^{11}Mo`$. This is about the gravitating mass needed to keep the X-ray gas binded to the galaxy. Following Fornan, Jones and Tucker (1985) formalism, the total gravitating mass within a 50 kpc radius is estimated between $`815\times 10^{11}Mo`$ for gas temperatures between 0.6 and 1 keV, which is in the order of magnitude of the mass derived from the ionized gas kinematics.
The results so far derived show compatible with a cooling flow process being dominating the gas emission. If the extended gas emission is modeled as that of a uniform sphere of hot gas emitting via bremsstrahlung, for a maximum radius of about 170 kpc (the size of the radio structure) and a temperature in the 0.6 - 1 keV range, the implied density would be $`simeq6\times 10^3cm^3`$. This yields a cooling time of about $`4\times 10^9yr`$, considerable smaller than the Hubble time. Thus, a cooling flow process could be operating in 3C 382.
A better characterization of the extended X-ray emission in 3C382 would demand much larger spatial resolution but also deeper observations.
Acknowledgments:
It is a pleasure to thank Peter Predehl for his support regarding the de-speckling procedure of the HRI data and Günther Hasinger and Paddy Leahy for critical reading of early versions of the manuscript.
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Figure captions
Figure 1a:
HRI contour image of 3C 382 extracted from channels 1 to 8. It is background-subtracted and smoothed with a Gaussian filter with a FWHM $``$12 arcsec. Contours are $`10^3ctss^1arcmin^2\times `$ (2.8, 4.4, 6, 7.6, 9.2, 10, 16, 40,4000); the first contour is about the 2$`\sigma `$ level measured on the background-subtracted image. 1 sky-pixel is 0.5 arcsec.
Figure 1b:
Broad band HST/WFPC2 image of 3C382 (data set u27l6w01-2) at 7200 A an equivalent exposure of 300 seconds. 3C382 is filling the upper left CCD. The bright companion galaxy (upper right CCD) is North-East of 3C382. The image covers 2.5 arcminutes.
Figure 2a:
HRI surface brightness profile of 3C 382 before de-speckling. Points with error bars (Poissonian noise) are the data. Dashed line is the HRI PRF; the continuum line is the PRF plus background. The background level is measured in an annulus between 3 to 5 arcmin from the center (flat part of the profile). The residuals between the data and the PRF+background are shown below. Note the large departure from a point source profile due to incorrect attitude reconstruction of the ROSAT data (see text).
Figure 2b:
HRI surface brightness profile of 3C 382 after de-speckling. The interpretation of the figure is as in Fig. 2a.
Figure 2c: The same as in Fig 2b but in this case the continuum line is the model fit to the surface brightness profile using a combination of an unresolved component represented by the HRI PRF (dash) and an extended component represented by a $`\beta `$ model (dots). The reduced $`\chi ^2`$ of the fit is 1.7-2 over 35 degrees of freedom (d.o.f.).
Figure 3:
Observed and best-model fit to the PSPC spectrum of 3C 382. PSPC data are from Prieto (1996). The dotted line represents the thermal 0.6-keV component; the dashed line, the power-law $`\alpha =0.7`$ component. Residuals represent 1-$`\sigma `$ errors. The first two bins corresponding to energies $`0.1keV`$ are not included in the fit. |
warning/0003/gr-qc0003040.html | ar5iv | text | # Group Averaging and Refined Algebraic Quantization
## 1 Introduction
Refined Algebraic Quantization (RAQ) is an attempt (amongst others) to concretize Dirac’s program for the quantization of constrained systems within a generally applicable, well defined mathematical framework. It was first formulated as a general scheme in and recently developed further in . Here I wish to report on these recent developments.
The method itself has already been used (partly implicitly) earlier in some successful applications to quantum gravity in specialized situations, like linearized gravity on symmetric backgrounds or various minisuperspace models . These suggested the program to find a general scheme of which these cases are just special cases.
Any scheme that deals with constrained systems needs to interpret the phrase ‘solving the constraints’. Here the guiding idea of RAQ is to work from the onset within an auxiliary Hilbert space, $`_{\mathrm{aux}}`$, by means of which a $``$-algebra of observables, $`𝒜_{\mathrm{obs}}`$, is constructed before the constraints are ‘solved’. The $``$-operation on $`𝒜_{\mathrm{obs}}`$ derives from the adjoint-operation $``$ on $`_{\mathrm{aux}}`$, which allows to connect the auxiliary inner product with the inner product on the physical Hilbert space, $`_{\mathrm{phys}}`$, since the latter is required to support a $``$-representation of $`𝒜_{\mathrm{obs}}`$. A possibly non-trivial limitation derives from the fact that the constraint operators on $`_{\mathrm{aux}}`$ are required to be self-adjoint, which may not be possible in case they do not form a Lie-algebra (i.e. ‘close’ with structure-functions only). This difficulty clearly does not arise if the constraints derive from the action of a Lie-group $`G`$. In the following we shall restrict attention to such cases. More precisely, we consider situations where a finite dimensional Lie group $`G`$ acts by some unitary representation $`U`$ on $`_{\mathrm{aux}}`$.
## 2 Concretizing the Dirac Procedure
From the RAQ point of view, the concretization of the Dirac procedure starts form an auxiliary Hilbert space $`_{\mathrm{aux}}`$ with inner product $`|_{\mathrm{aux}}`$ and the set of unitary operators $`\{U(g)gG\}`$. The naive reading of Dirac’s prescription is to identify the space of physical states with those elements in $`_{\mathrm{aux}}`$ that are fixed by $`G`$:
$$U(g)|\psi =|\psi ,gG,$$
(1)
which just says that one should pick the trivial subrepresentation of $`U`$. But clearly this statement needs not be well defined since $`U`$ might simply not contain the trivial representation as sub-representation. This will happen if the operators $`U(g)`$ do not all have the value ‘one’ in the discrete part of their spectrum. Another difficulty comes in if $`G`$ is not unimodular. Then it has been convincingly argued using methods of geometric quantization (see and references therein) that (1) is simply not the right condition, but that (1) should formally be replaced with
$$U(g)|\psi =\mathrm{\Delta }^{1/2}(g)|\psi ,$$
(2)
where $`\mathrm{\Delta }(g):=\mathrm{det}[\mathrm{Ad}_\mathrm{g}]`$ is the modular function on $`G`$ (a one-dimensional, real representation of $`G`$). We said ‘formally’ since (2) cannot hold in $`_{\mathrm{aux}}`$, for this would mean that the unitary operator $`U(g)`$ can change lengths by $`\mathrm{\Delta }^{1/2}(g)1`$, a plain contradiction. This is connected with the first problem since it means that $`\mathrm{\Delta }^{1/2}(g)`$ cannot be a discrete spectral value of $`U(g)`$. In physical terms, $`|\psi `$ in (2) is not normalizable. Hence one must read these equations in the appropriate sense. It is useful to break up the further development into several steps.
Step 1: We denote by $``$ the adjoint map on operators on $`_{\mathrm{aux}}`$ with respect to $`|_{\mathrm{aux}}`$. Choose a dense linear subspace $`\mathrm{\Phi }_{\mathrm{aux}}`$ which is left invariant (as set) by $`G`$’s action. This choice of $`\mathrm{\Phi }`$ is an important step and will generally require some physical input. Let $`\mathrm{\Phi }^{}`$ be the algebraic dual (linear functionals) of $`\mathrm{\Phi }`$. Put on $`\mathrm{\Phi }^{}`$ the topology of pointwise convergence: $`f_nf`$ in $`\mathrm{\Phi }^{}`$ iff $`f_n[\varphi ]f[\varphi ]`$ as real numbers for all $`\varphi \mathrm{\Phi }`$. Note that, conversely, each $`\varphi \mathrm{\Phi }`$ defines a continuous linear functional on $`\mathrm{\Phi }^{}`$ by setting $`\varphi (f):=f[\varphi ]`$. Hence $`\mathrm{\Phi }`$ embeds in the topological dual (continuous linear functionals) of $`\mathrm{\Phi }^{}`$.
Step 2: Consider the set $`:=\{A_{\mathrm{aux}}D(A)_{\mathrm{aux}}\}`$ of linear operators, where $`D(A)`$ denotes the domain of $`A`$. It is not an algebra due to mismatches of ranges and domains. We define a subset of $``$ by
$`𝒜:=\{A\mathrm{\Phi }D(A)D(A^{})`$
$`\mathrm{\Phi }A(\mathrm{\Phi }),\mathrm{\Phi }A^{}(\mathrm{\Phi })\},`$ (3)
and make it into an algebra by restricting the action of each $`A𝒜`$ to $`\mathrm{\Phi }`$. Since $``$ also restricts to $`𝒜`$, we have a $``$-algebra which – without indicating the restriction to $`\mathrm{\Phi }`$ – we continue to call $`𝒜`$. Note that, by definition of $`\mathrm{\Phi }`$, $`𝒜`$ contains the operators $`U(g)`$. From now on, the $``$-operation of this algebra is the only trace left by $`|_{\mathrm{aux}}`$. Finally, we define a sub-$``$-algebra $`𝒜_{\mathrm{obs}}𝒜`$ as the ‘commutant’ in $`𝒜`$ of the set $`\{U(g)gG\}`$:
$$𝒜_{\mathrm{obs}}:=\{A𝒜U(g)A=AU(g),gG\}.$$
(4)
Step 3: Each $`A𝒜`$ acts as continuous linear map on $`\mathrm{\Phi }^{}`$ via the ‘adjoint’ action:
$$Af:=fA^{}.$$
(5)
Hence the constraint operators also act on $`\mathrm{\Phi }^{}`$ and we can define the solution set, $`𝒱`$, by
$$𝒱:=\{f\mathrm{\Phi }^{}U(g)f=\mathrm{\Delta }^{1/2}(g)f\},$$
(6)
which, by construction, carries an anti-$``$-representation of $`𝒜_{\mathrm{obs}}`$ by continuous maps. (The ‘anti’ is a consequence of the $``$ in (5) and does no harm)
Step 4: The Hilbert space of physical states is now to be found within $`𝒱`$. Hence we seek a subspace $`_{\mathrm{phys}}𝒱`$ with inner product $`|_{\mathrm{phys}}`$ that makes $`_{\mathrm{phys}}`$ into a Hilbert space. Generally one cannot turn all of $`𝒱`$ into a Hilbert space, since this would imply that all $`A𝒜_{\mathrm{obs}}`$ were defined everywhere on $`_{\mathrm{phys}}`$ and hence bounded, which is too restrictive. (Note: Merely having the whole Hilbert space as domain does not yet imply boundedness of a linear operator $`A`$. But if $`A`$ and its adjoint are defined everywhere boundedness follows.) Hence $`_{\mathrm{phys}}`$ will in general be a proper subset of $`𝒱`$ and the topology of $`_{\mathrm{phys}}`$ induced by $`|_{\mathrm{phys}}`$ must be finer than that it inherits from $`\mathrm{\Phi }^{}`$, since the former must be closed. This also means that operators in $`𝒜_{\mathrm{obs}}`$ will not necessarily be bounded on $`_{\mathrm{phys}}`$ and hence we do not have an anti-$``$-representation of $`𝒜_{\mathrm{obs}}`$ on $`_{\mathrm{phys}}`$ since domains and ranges might not match. Hence we proceed as usual by assuming that there is a dense subspace $`\mathrm{\Phi }_{\mathrm{phys}}_{\mathrm{phys}}`$ on which we have such a representation. But besides these technicalities the relevant condition on $`|_{\mathrm{phys}}`$ is this: the physical inner product is to be chosen such that for all $`A𝒜_{\mathrm{obs}}`$ we have $`A^{}=A^{}`$ on $`A^{}`$’s domain in $`_{\mathrm{phys}}`$, where $``$ now denotes the adjoint operation with respect to $`|_{\mathrm{phys}}`$. This is how $`|_{\mathrm{phys}}`$ is influenced (but generally not determined) by $`|_{\mathrm{aux}}`$. It is with respect to this adjoint operation that we speak of an anti-$``$-representation of $`𝒜_{\mathrm{obs}}`$ on $`\mathrm{\Phi }_{\mathrm{phys}}`$.
## 3 RAQ and the $`\eta `$-Map
RAQ concretizes step 4 of the Dirac procedure just outlined. It aims to construct $`_{\mathrm{phys}}`$ by finding a so-called $`\eta `$-map (or ‘rigging map’), which is an antilinear map $`\eta :\mathrm{\Phi }\mathrm{\Phi }^{}`$ such that the image of $`\eta `$ consists entirely of solutions, i.e., $`\eta (\mathrm{\Phi })𝒱`$. Further conditions on this map are
1. $`\eta `$ is real: $`\eta (\varphi _1)[\varphi _2]=\overline{\eta (\varphi _2)[\varphi _1]}`$
2. $`\eta `$ is positive: $`\eta (\varphi )[\varphi ]0`$
3. $`\eta `$’s image is an invariant domain for $`𝒜_{\mathrm{obs}}`$ and $`\eta `$ intertwines the representations of $`𝒜_{\mathrm{obs}}`$ on $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$:
$$A\eta (\varphi )=\eta (A\varphi )$$
(7)
Given such an $`\eta `$, one defines $`|_{\mathrm{phys}}`$ on its image by
$$\eta (\varphi _1)|\eta (\varphi _2)_{\mathrm{phys}}:=\eta (\varphi _2)[\varphi _1]$$
(8)
and then defines $`_{\mathrm{phys}}`$ as the completion of $`\eta `$’s image with respect to the uniform structure (Cauchy sequences) defined by this inner product. By construction, $`|_{\mathrm{phys}}`$ satisfies the condition that the $``$-operation on $`𝒜_{\mathrm{obs}}`$ is the adjoint with respect to $`|_{\mathrm{phys}}`$.
We need only check two points in order to see that the $`_{\mathrm{phys}}`$ so defined satisfies the conditions of step 4: First, we wanted $`_{\mathrm{phys}}`$ to be a subset of $`𝒱`$, so we need to check that the completion just mentioned does not add points outside $`𝒱`$, i.e., that $`_{\mathrm{phys}}`$ is indeed a subset of $`𝒱`$ with finer intrinsic topology. To see this, we consider the map $`\sigma :_{\mathrm{phys}}\mathrm{\Phi }^{}`$, defined by
$$\sigma (f)[\varphi ]:=\eta (\varphi )|f_{\mathrm{phys}},$$
(9)
and note immediately that $`\sigma (f)`$ vanishes iff $`f`$ is orthogonal to all elements in the images of $`\eta `$. But since the image is dense in $`_{\mathrm{phys}}`$ by construction $`f`$ must itself vanish, hence $`\sigma `$ is injective. Next we prove that $`\sigma `$ is continuous: if $`f_nf`$ in $`_{\mathrm{phys}}`$ then $`\eta (\varphi )|f_n_{\mathrm{phys}}\eta (\varphi )|f_{\mathrm{phys}}`$ for all $`\varphi `$, since $`|\eta (\varphi )|_{\mathrm{phys}}`$ is a continuous linear form on $`_{\mathrm{phys}}`$. Hence $`\sigma (f_n)[\varphi ]\sigma (f)[\varphi ]`$ for all $`\varphi `$ which, since $`\mathrm{\Phi }^{}`$ carries the topology of pointwise convergence, implies $`\sigma (f_n)\sigma (f)`$ and therefore continuity of $`\sigma `$.
The second point to check is that we indeed have an anti-$``$-representation of $`𝒜_{\mathrm{obs}}`$ on a dense subspace $`\mathrm{\Phi }_{\mathrm{phys}}`$ in $`_{\mathrm{phys}}`$. To show this, we simply identify $`\mathrm{\Phi }_{\mathrm{phys}}`$ with the image of $`\eta `$ and note that by (5) and (7) we have for all $`\varphi \mathrm{\Phi }`$:
$`\sigma (Af)[\varphi ]`$ $`=`$ $`\eta (\varphi )|Af_{\mathrm{phys}}=\eta (A^{}\varphi )|f_{\mathrm{phys}}`$ (10)
$`=`$ $`\sigma (f)[A^{}\varphi ]=A\sigma (f)[\varphi ].`$
Hence $`\sigma :\mathrm{\Phi }_{\mathrm{phys}}\text{Image}(\sigma )\mathrm{\Phi }^{}`$ is an isomorphism of anti-$``$-representations of $`𝒜_{\mathrm{obs}}`$.
An interesting question is how general this method of constructing $`_{\mathrm{phys}}`$ via an $`\eta `$-map actually is; that is, given $`_{\mathrm{phys}}`$ as in step 4 above, is there always an $`\eta `$ map whose image is $`_{\mathrm{phys}}`$? Well, remember that each $`\varphi \mathrm{\Phi }`$ defines a continuous linear functional on $`\mathrm{\Phi }^{}`$. Restriction to the linear subspace $`_{\mathrm{phys}}`$ of $`\mathrm{\Phi }^{}`$ yields a continuous linear functional on $`_{\mathrm{phys}}`$ with respect to its intrinsic topology, since the latter is finer than that induced by $`\mathrm{\Phi }^{}`$. Hence, by Riesz’ theorem, there is a a unique vector $`\eta ^{}(\varphi )_{\mathrm{phys}}`$ which satisfies
$$\varphi (f):=f[\varphi ]=\eta ^{}(\varphi )|f_{\mathrm{phys}}$$
(11)
for each $`f_{\mathrm{phys}}`$. The map $`\eta ^{}:\mathrm{\Phi }_{\mathrm{phys}}`$ is obviously antilinear and has a non-trivial image, since $`\text{kernel}(\eta ^{})=_{f𝒱}\text{kernel}(f)\mathrm{\Phi }`$ for $`𝒱\{0\}`$. We can now define $`\mathrm{\Phi }_{}^{}{}_{\mathrm{phys}}{}^{}`$ to be the image of this $`\eta ^{}`$ and $`_{}^{}{}_{\mathrm{phys}}{}^{}_{\mathrm{phys}}`$ its completion. This almost proves that we can at least reproduce the subspace $`_{}^{}{}_{\mathrm{phys}}{}^{}`$ by an $`\eta `$-map if $`\eta ^{}`$ were indeed an $`\eta `$-map in the technical sense. But it might fail condition 3 above that it intertwines with $`𝒜_{\mathrm{obs}}`$; the reason being that so far nothing ensures $`\mathrm{\Phi }_{}^{}{}_{\mathrm{phys}}{}^{}`$ to be contained in $`\mathrm{\Phi }_{\mathrm{phys}}`$. Hence operators in $`𝒜_{\mathrm{obs}}`$ might not act on $`\mathrm{\Phi }_{}^{}{}_{\mathrm{phys}}{}^{}`$ at al! However, if $`\mathrm{\Phi }_{}^{}{}_{\mathrm{phys}}{}^{}\mathrm{\Phi }_{\mathrm{phys}}`$ then $`\eta ^{}`$ is indeed an $`\eta `$-map and we can at least claim to be able to reconstruct some sector $`_{}^{}{}_{\mathrm{phys}}{}^{}_{\mathrm{phys}}`$.
Given that, we may finally wonder what additional assumptions would guarantee that we can construct all of $`_{\mathrm{phys}}`$ by RAQ. One such condition is the following : If $`_{\mathrm{phys}}`$ constructed by the Dirac procedure decomposes into a direct sum of superselection sectors, then each sector can be separately constructed by the Dirac procedure. If this is satisfied we argue as follows: suppose $`_{}^{}{}_{\mathrm{phys}}{}^{}`$ is a maximal subsector of $`_{\mathrm{phys}}`$ which can be constructed by RAQ. Then its orthogonal complement is also a sector which, by hypothesis, is separately constructible by the Dirac procedure. But then, again by hypothesis, we can reconstruct a subsector of this orthogonal complement, which contradicts the assumption that $`_{}^{}{}_{\mathrm{phys}}{}^{}`$ was maximal.
## 4 Group Averaging
Whereas the $`\eta `$-map is a specific way to build $`_{\mathrm{phys}}`$ for step 4 in the Dirac procedure, Group Averaging is in turn a specific way to construct an $`\eta `$-map. In Dirac’s $`\text{bra}|\text{ket}`$-terminology, the idea is simply to define (and make sense of)
$$\eta |\varphi :=_G𝑑\mu (g)\varphi |U(g),$$
(12)
where $`d\mu `$ is some appropriately chosen measure on $`G`$. If $`G`$ were unimodular, the right- and left-invariant Haar measures coincide and are the correct choice for $`d\mu `$ in order for $`\eta `$’s image to formally solve (1) and make $`\eta `$ real in the sense of condition 1 above. Note that this reality condition requires the measure $`d\mu (g)`$ to be invariant under the inversion map $`I:gg^1`$. For non-unimodular groups neither the left- nor the right-invariant measure is $`I`$-invariant. Instead one has to take the ‘symmetric’ measure, defined by
$`d\mu _0(g):`$ $`=`$ $`\mathrm{\Delta }^{1/2}(g)d\mu _L(g)`$ (13)
$`=`$ $`\mathrm{\Delta }^{1/2}(g)d\mu _R(g),`$
where $`d\mu _L,d\mu _R`$ are the left- and right-invariant measures respectively. Moreover, this measure also implies that $`\eta `$’s image (formally) satisfies the modified equation (2). We refer to for a more extensive discussion of the modified Dirac condition and how it can be derived by Group Averaging from an appropriate adaptation of the ‘unimodularization’ technique originally developed in geometric quantization .
The physical inner product according to Group Averaging is given by
$$\eta (\varphi _2)|\eta (\varphi _1)_{\mathrm{phys}}:=_G𝑑\mu _0(g)\varphi _1|U(g)\varphi _2_{\mathrm{aux}}$$
(14)
so that the Group Averaging procedure only makes sense for states $`\varphi _{1,2}`$ for which this integral converges absolutely. But it turns out that a more restricted choice of $`\mathrm{\Phi }`$ is appropriate: we say that $`\varphi `$ is an $`L^1`$ state, iff
$$_g𝑑\mu _n\varphi |U(g)\varphi _{\mathrm{aux}}$$
(15)
converges absolutely for all integers $`n`$, where $`d\mu _n(g):=\mathrm{\Delta }^{n/2}(g)d\mu _0`$. One can then also construct the group algebra $`𝒜_G`$ of functions on $`G`$ which are $`L^1`$ with respect to $`d\mu _n`$ for all $`n`$, and prove that its action on $`L^1`$ states results again in $`L^1`$ states . This fact allows to prove a uniqueness theorem in the following form (see for more details):
###### Theorem 1
Suppose $`\mathrm{\Phi }`$ is an $`L^1`$ subspace of $`_{\mathrm{aux}}`$ which is invariant under $`𝒜_G`$. Then, up to overall scale, any $`\eta `$-map is of the form (12) with $`d\mu =d\mu _0`$.
## 5 Summary
We outlined two technical devices which, to a certain extent, concretize the Dirac approach to quantizing constrained systems along the sequence \[Group Averaging\] $``$ \[$`\eta `$-map\] $``$ \[Dirac approach\]. There are fairly strong uniqueness theorems regarding these devices, provided Group Averaging converges sufficiently rapidly. But there is no general characterization when this can actually be achieved. Physical as well as mathematical inputs are required, perhaps most importantly in the choice of $`\mathrm{\Phi }`$, and results are likely to delicately depend on this choice. Recent investigations of convergence properties in contexts of specific models confirm this general expectation.
Finally we mention the obvious, namely that the restriction to systems whose first class constraints close with structure constants (i.e. form a Lie algebra) should be lifted. After all, GR is not in this form. We plan to investigate this in the future. |
warning/0003/cond-mat0003061.html | ar5iv | text | # Magnetotransport in an array of magnetic antidots
## Abstract
Classical transport properties of an electron, moving in plain, in an array of magnetic antidots has been calculated. Homogeneous magnetic field $`\stackrel{}{B}=B\stackrel{}{e}_z`$ fills the whole space except of cylinders of radius $`r_0`$, which are arranged in a square lattice with a lattice constant $`a`$. The magnetoresistance shows additional peaks and minima by varying the strength of the magnetic field, according to pinned orbits at antidots and to propagating orbits in transport direction, respectively.
The dynamics of charged particles in spatially modulated magnetic field give rise to a variety of interesting phenomena. The motion of ballistic electrons in modulated magnetic field is also believed to be closely related to the motion of composite fermions in a density modulated 2DEG in the fractional quantum Hall regime . Commensurability effects in weak periodic magnetic fields were predicted and observed for modulation in one and two directions experimentally. Such periodic magnetic field can be generated by using either patterned superconducting or ferromagnetic gates ,. Closely related to our system are antidot arrays, which consist of periodically arranged voids in an electron gas. Charged particles move in this two dimensional potential landscape under the influence of perpendicular magnetic field and are scattered on the antidots. Additional peaks in the low field magnetoresistance have been observed in these antidot arrays and explained by a pinning mechanism of classical circular orbits enclosing 1,2,4,9 and 21 antidots in terms of nonlinear dynamics. The nature of magnetoresistance peaks are caused by islands of regular motion due to nonlinear resonances.
In the present paper we study the magnetotransport in an array of magnetic antidots. Instead of antidot scatterers we introduce circular regions into the system with zero magnetic field inside. Inside the magnetic antidots the electron moves in a field-free region and outside is its motion described by the Lorentz force. As demostrated in Fig. 1 the inhomogenity of the z-componetnt of the magnetic field for one magnetic antidot is given by
$$B(r)\stackrel{}{e}_z=\{\begin{array}{ccc}B\stackrel{}{e}_z& \text{for}& r>r_0\\ 0& \text{else}& \end{array}$$
and these antidots are arranged in a rectangular array of lattice constant $`a`$. In the recent experiments on modulated magnetic fields, the period of modulation lies by 700 nm and is larger then the Fermi wave length, but smaller then the elastic mean free path (10 $`\mu `$m) and thus the dynamics of the wave packet approaches the classical limit. Magnetic antidot array could be realized using supeconducting materials, which would be maped on two dimensional electron gas in a rectangular array arrangement. By applying magnetic field, in circular region of superconducting material, the applied magnetic field would be swept out and the electrons would move in a system described in this paper. In the calculation we will use a system, that would be produced by an ideal vortex with a constant magnetic field outside and zero magnetic field inside, instead of more realistic exponential distribution of the magnetic. We consider also one elctron approximation.
The classical trajectory is found as the solution to Newtons’ equation of motion with the force given by the Lorentz expression $`\stackrel{}{F}=e\stackrel{}{v}\times \stackrel{}{B}`$ for a particle of charge $`e`$. It consist, as well known of straight line segments inside the antidot and an arc of a circle outside, with the radius of the curvature given by the cyclotron radius $`r_c=\frac{v}{\omega _c}`$ with $`v`$ being the particle velocity and $`\omega _c=\frac{eB}{m}`$ the cyclotron frequency. There exist also undisturbed circular orbits, which do not intersect any antidot. The classical approximation for the dynamics of an electron wave packet in a modulated magnetic field is described by the Hamiltonian for one magnetic antidot
$$H(x,y,p_x,p_y)=\{\begin{array}{ccc}[(p_x+eBy/2)^2+(p_yeBx/2)^2]/2m& \text{for}& r>r_0\\ \stackrel{}{p}^2/2m& \text{else}& \end{array}$$
where $`m`$ is the effective mass of the electron. First we introduce some characteristic parameters of the system and express all quantities in dimensionless units: the coordinate $`x=x/a`$, $`y=y/a`$,$`t=t/\tau _0`$ and $`H=H/ϵ_F`$ with $`ϵ_F`$, the Fermi energy, $`\tau _0=(ϵ_F/ma^2)^{1/2}`$ . The magnetic field is scaled by $`B_0`$ at which the cyclotron radius corresponds to $`a/2`$. Calculating the Poisson brackets and omitting the primes we obtain following equations of motion for one magnetic antidot
$$\dot{x}=v_x,\dot{y}=v_y$$
$$\dot{v}_x=\{\begin{array}{ccc}B/B_0v_y& \text{for}& r>r_0\\ 0& \text{else}& \end{array}$$
$$\dot{v}_y=\{\begin{array}{ccc}B/B_0v_x& \text{for}& r>r_0\\ 0& \text{else}& \end{array}$$
Choosing $`r_0=0.3a`$ and using a numerical integration method it is possible to calculate the trajectories of the electron for different values of magnetic field. We investigate the motion in phase space ($`x,y,v_x,v_y`$) by means of Poincaré surface of section at ($`x=0\text{mod}a`$). Since the energy $`E=v_x^2/2m+v_y^2/2m`$ of the system is conserved, we have to consider the surface of section ($`y,v_y`$) for $`v_x>0`$ and $`v_x<0`$ separately. The initial conditions ($`y_i,v_{yi}`$) lead for $`v_x>0`$ and $`v_x<0`$ to different trajectories and the corresponding two surfaces of section have mirror symmetry. Fig.2 shows the surface of section for $`v_x>0`$ for three different values of magnetic field $`B/B_0=0.5,0.6\text{and}1.0`$, the corresponding cyclotron radius is $`r_c=a,0.83a,a/2`$ . The whole phase space consits at $`B/B_0=0.5`$ of chaotic sea as shown in Fig.2 (c), one of possible chaotic trajectories in coordinate space is shown in Fig.3 (IV). At $`B/B_0=0.6`$ we found periodic orbits, which intersect two magnetic antidots as shown in Fig.3 (III). These orbits occupy a relative small portion in the phase space as shown in Fig.2 (b), label(III). An other type of periodic orbits, which intersect four magnetic antidots is shown in Fig.3 (V). The fingerprint of this orbit in the surface of section is shown in Fig.2 (a), label (V). When the cyclotron radius approaches half of the lattice constant, a different type of periodic orbit becomes dominat. These orbits are pinned around one magnetic antidot as shown in Fig.3 trajectory (II) and build in the surface of section a family of eliptic fixed points in a circular region at $`(y,v_y)=(0.5,0)`$ of diameter $`r_0`$ as shown in Fig.2 (a), region (II). These orbits can also intersect one magnetic antidot for higher magnetic fileds as shown in Fig.3 (I) and dominate the dynamics of the system. The orbits of type (I),(II) and (V) are present in a wide sector of magnetic field strength $`B/B_0]0.8,1.7]`$. For lower values of magnetic field $`B/B_0=0.3`$ periodic orbits, which intersect eight magnetic antidots occupy a relative small region (10%) in phase space. For $`B/B_0>2.3`$ all possible orbits are pinned, and no classical chaotic trajectories exist anymore. The inset of Fig.4 (a) shows the portion $`p`$ of regular orbits in dependence of magnetic field.
At low temperatures elastic impurity scatterings with mean scattering time $`\tau `$ have to be considered. Using the classical Kubo formula which includes these scattering, we can calculate the conductivity $`\sigma _{ij}`$ in dependence of magnetic field.
$$\sigma _{ij}(0,\tau )=\frac{ne^2}{k_BT}_0^{\mathrm{}}𝑑te^{1/\tau }<v_i(t)v_j(0)>_{\mathrm{\Gamma }_c}$$
$`<v_i(t)v_j(0)>`$ is the velocity correlation function averaged over chaotic part of the phase space $`\mathrm{\Gamma }_c`$, $`k_B`$ is the Boltzmann constant, $`T`$ temperature and $`n`$ the 2D electron density. The magnetoresistance can be obtain by inversion of the conductivity tensor $`R_{xx}=\sigma _{xx}/(\sigma _{xx}^2+\sigma _{xy}^2)`$ and $`R_{xy}=\sigma _{xy}/(\sigma _{xx}^2+\sigma _{xy}^2)`$. We have calculated the magnetoresistance $`R_{xx}`$ and the Hall resistance $`R_{xy}`$ in dependence of magnetic field for a typical value of $`\tau /\tau _0=30`$. The magnetoresistance shown in Fig.4 (a) displays a dominant peak at $`B/B_01`$ which corresponds to pinned orbits around one antidot. According to the portion of pinned orbits, see the inset of Fig.4 (a) we conclude, that this peak is caused by these periodic orbits. We should expect also a peak at $`B/B_00.3\text{and}0.6`$ as in the case of antidot lattices, but instead we observe a minimum at $`B/B_0=0.8`$. The main difference to antidot lattices, where electrons are scattered on antidots, is that the sign of the angular momentum of an electron, moving in a magnetic antidot array remains conserved. The magnetoresistance is not only caused by islands of regular motion in the phase space (pinned orbits), but also by orbits, which propagate in the direction of transport. For $`B/B_0=0.8`$ we found orbits, which propagate in transport direction, as shows in Fig. 3 (b) and explain the minimun in the magnetoresistance at this strength of magnetic field. We obtain even $`R_{xx}/R_{xx}(0)<1`$, the resistance lies lower then, for the case of zero magnetic field. The Hall resistance as shown in Fig. 4 (b) shows steps at magnetic field values for which the magnetoresistance $`R_{xx}`$ lies in minimum.
In conclusion, classical transport properties of electrons, moving in a magnetic antidot array has been explained by pinned orbits at antidots and by propagating orbits in transport direction.
I thank Prof.B.Kramer for the initial ideas of this work.
Figure captions |
warning/0003/hep-ex0003009.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The study of Bose-Einstein correlations between pairs of identical hadrons is an essential tool to obtain information on the space-time evolution of the extended hadron sources created in heavy ion collisions . In particular, a strong correlation between the momenta and the space-time production points of the particles suggests expanding sources as predicted by hydrodynamic models . The dynamical evolution of such systems can then be studied with interferometry via selection on the transverse momenta and rapidity of the correlated particle pairs.
In this paper, we present the analysis of two-particle correlations of identified $`\pi ^{}`$ measured in the WA98 experiment for central 158 AGeV Pb+Pb collisions at the CERN SPS.
## 2 Experimental setup and data processing
The WA98 experiment shown in Fig. 1 combined large acceptance photon detectors with a two arm charged particle tracking spectrometer. The incident 158 AGeV Pb beam interacted with a Pb target near the entrance of a large dipole magnet. Non-interacting beam nuclei, or beam fragments were detected in a forward calorimeter located at zero degree. A mid-rapidity calorimeter measured the total transverse energy in the rapidity region 3.2 $``$ $`\eta `$ $``$ 5.4, which was also used in the trigger for online centrality selection. The Plastic-Ball calorimeter measured the fragmentation of the target, and silicon detectors were used to measure the charged particle multiplicity. The photon detectors consisted of a large area photon multiplicity detector and a high granularity lead-glass calorimeter for single photon, $`\pi ^0`$, and $`\eta `$ physics .
The charged particle spectrometer made use of a 1.6 Tm dipole magnet with a 2.4$`\times `$1.6 m<sup>2</sup> air gap for magnetic deflection of the charged particles in the horizontal plane. The results presented in this paper are taken from the 1995 WA98 data set obtained with the negative particle tracking arm of the charged particle spectrometer. The second tracking arm was added to the spectrometer in 1996 to measure positive particles . The first tracking arm consisted of six multistep avalanche chambers with optical readout located downstream of the magnet. The active area of the first chamber was 1.2$`\times `$0.8 m<sup>2</sup>, while that of the other five was 1.6$`\times `$1.2 m<sup>2</sup>. The chambers contained a photoemissive vapour (TEA) which produced UV photons along the path of traversal of the charged particles. These were converted into visible light via wavelength shifter plates. On exit the light was reflected by mirrors at 45 to CCD cameras equipped with two image intensifiers. Each pixel of a CCD viewed a 3.1$`\times `$3.1 mm<sup>2</sup> area of a chamber. In addition, a 4$`\times `$1.9 m<sup>2</sup> Time of Flight wall positioned behind the chambers at a distance of 16.5 m from the target allowed for particle identification with a time resolution better than 120 ps.
Fig. 3 shows the Monte Carlo generated $`p_T`$-rapidity acceptance for $`\pi ^{}`$. The acceptance ranges from $`y`$=2.1 to 3.1 with an average at 2.70. The momentum resolution of the spectrometer was $`\mathrm{\Delta }p/p`$=0.005 at $`p`$=1.5 GeV/c, resulting in an average accuracy better than or equal to 10 MeV/c for all the Q variables used in the interferometry analysis and defined in section 5: $`\sigma (Q_{inv})`$=7 MeV/c, $`\sigma (Q_{TO})`$=10 MeV/c, $`\sigma (Q_{TS})`$=5 MeV/c, $`\sigma (Q_L)`$=3 MeV/c, $`\sigma (Q_T)`$=8 MeV/c, $`\sigma (Q_0)`$=5 MeV/c.
The analysis of the complete 1995 data set is presented here. These data have been taken with the most central triggers corresponding to about 10% of the minimum bias cross section of 6190 mb. Severe track quality cuts were applied at the expense of statistics resulting in final samples of 4.2$`\times 10^6`$ $`\pi ^{}`$ for the correlation analysis and 4.6$`\times 10^5`$ $`\pi ^{}`$ for the single particle spectrum.
## 3 Single particle spectra
The $`m_T`$=$`\sqrt{m_\pi ^2+p_T^2}`$ distribution of identified $`\pi ^{}`$, averaged over the rapidity acceptance, is shown in Fig. 3. The data were corrected for geometrical acceptance and efficiency of the chamber-camera-Time of Flight system using a full simulation of the experimental setup. The parameters of the simulation were optimized in an iterative way by comparing various distributions with the real data. The simulated data were then treated exactly like the real data. The measured $`1/m_TdN/dm_Tdy`$ distribution was then fitted to the form $`C`$exp$`(m_T/T)`$, expected for a source in thermal equilibrium . Such fits were applied to the data for different ranges of $`m_T`$, such as the one shown in Fig. 3. These fits do not reproduce the overall concave shape of the data, which is partly due to particles originating in resonance decays and could also be an indication of transverse flow . The shape of the $`\pi ^{}`$ $`m_T`$ distribution was found to be in good agreement with that of $`\pi ^0`$ obtained in the lead-glass calorimeter .
## 4 One-dimensional interferometry analysis
For the Bose-Einstein correlation studies, the data were Coulomb corrected in an iterative way . The Gamow correction was abandoned as it overcorrects the data for $`Q_{inv}`$ in the range of 0.1 to 0.3 GeV/c. A fit of the form $`1+\lambda \text{exp}[Q_{inv}^2R_{inv}^2]`$ was made to the $`Q_{inv}`$ correlation function yielding $`R_{inv}=6.83\pm 0.10`$ fm and $`\lambda =0.307\pm 0.008`$. An expanded view of the correlation distribution (Fig. 5) shows that the Gaussian fit used (full line) is not perfect, especially in the $`Q_{inv}`$ range of 40 to 80 MeV/c where the tail of the experimental distribution shows an excess which is not well reproduced by the fit. In addition to this Gaussian fit made over the whole range of $`Q_{inv}`$, Fig. 5 shows also different Gaussian fits using data in the ranges 25 MeV/c $``$ $`Q_{inv}`$ $``$ 200 MeV/c (dashed line) and 40 MeV/c $``$ $`Q_{inv}`$ $``$ 200 MeV/c (dotted line). These fits do not coincide. Different radii are then obtained for different starting points of the fit because the shape of the distribution is not Gaussian. This effect is independent of the severity of the track selection, and is therefore not due to spurious tracks. This is summarized in Fig. 7 where $`R_{inv}`$ and the corresponding $`\lambda `$ are plotted as a function of the lower bound of the fit. There is a statistically significant drop when using a Gaussian fit. A similar behaviour is observed when, instead of $`Q_{inv}`$, $`Q_{space}=\sqrt{Q_x^2+Q_y^2+Q_z^2}`$ is used, calculated in the longitudinally comoving system (LCMS) and fitted with 1+$`\lambda `$exp\[-$`Q_3^2R_3^2]`$. This method of fitting in varying ranges has a good sensitivity to the shape. It has been repeated by replacing the Gaussian fit by an exponential fit of the form 1+$`\lambda _e`$exp\[-2$`Q_{inv}R_e]`$ where the factor 2 is added to make the radius $`R_e`$ more comparable with $`R_{inv}`$. The results (Fig. 7) show that the stability is better with the exponential fit. Fig. 5 directly compares the Gaussian and exponential fits for $`Q_{inv}`$. Although the Gaussian fit still gives an acceptable $`\chi ^2/`$d.o.f., the exponential fit is better everywhere. A similar conclusion is reached when the first data point is excluded from the fit. This result is not based on the first bins which might be more affected by systematics due to large Coulomb correction or noise correlated with the true track signals in the chambers. It is rather based on the high statistics tail of the distribution which contributes in a different way to a Gaussian or an exponential fit. This quasi exponential behaviour is expected by different models including resonance decays . As a consequence small acceptance experiments may obtain a larger radius if a Gaussian fit is used because they are less sensitive to the tail. On the contrary large acceptance experiments have higher statistics at large $`Q`$-values, and the Gaussian fit will yield lower values of the radius.
## 5 Multi-dimensional interferometry analysis
The multi-dimensional analysis has been done with Gaussian fits to allow comparison with other experiments. Two different parameterizations have been used in the LCMS:
a) The standard fit in the 3-dimensional space of momentum differences $`Q_{TS}`$ (perpendicular to the beam axis and to the transverse momentum of the pair), $`Q_{TO}`$ (perpendicular to the beam axis and parallel to the transverse momentum of the pair), and $`Q_L`$ (parallel to the beam axis) . The fitted formula
$$C_2=1+\lambda \mathrm{exp}[Q_{TS}^2R_{TS}^2Q_{TO}^2R_{TO}^2Q_L^2R_L^22Q_{TO}Q_LR_{outlong}^2]$$
includes a cross term in $`Q_{TO}Q_L`$ as predicted .
b) The generalized Yano-Koonin (GYK) fit in the $`Q_0`$ (energy difference of the pair), $`Q_T`$,$`Q_L`$ space according to
$$C_2=1+\lambda \mathrm{exp}[Q_T^2R_T^2+(Q_0^2Q_L^2)R_4^2(QU)^2(R_0^2+R_4^2)]$$
where $`U=\gamma (1,0,0,v_L)`$, $`\gamma =1/\sqrt{1v_L^2}`$ with $`v_L`$ in units of c=1.
In the GYK approach, the radius parameters remain invariant under longitudinal Lorentz boost, the parameter $`v_L`$ connecting the “arbitrary” measurement frame (LCMS) to the Yano-Koonin frame. In addition, the extraction of the duration of emission, $`R_0`$, is straightforward.
The consequence of the finite resolution in the measurement of the $`Q`$ variables is an underestimate of the radii and $`\lambda `$ parameters. Morever, as the resolution is different for each $`Q`$ variable, this causes a bias which varies from parameter to parameter, leading to errors in the interpretation of the results in a multi-dimensional analysis. It is therefore essential to take into account the effect of the resolution in the fitting procedure. One way to do this is to replace the formula $`C_2(\stackrel{}{Q})`$ used to fit the data by
$$C_2^{rc}(\stackrel{}{Q})=r(\stackrel{}{Q},\stackrel{}{Q^{}})C_2(\stackrel{}{Q^{}})𝑑\stackrel{}{Q^{}}$$
which is the convolution of $`C_2(\stackrel{}{Q})`$ with the resolution function $`r(\stackrel{}{Q},\stackrel{}{Q^{}})`$. The resolution function is chosen to be Gaussian:
$$r(\stackrel{}{Q},\stackrel{}{Q^{}})=1/(2\pi )^{3/2}\mathrm{\hspace{0.25em}1}/|V|^{1/2}\mathrm{exp}[1/2(\stackrel{}{Q}\stackrel{}{Q^{}})^TV^1(\stackrel{}{Q}\stackrel{}{Q^{}})]$$
The diagonal elements of the covariance matrix $`V`$ are equal to the square of the resolution of the different $`Q`$ variables and are estimated separately as a function of $`k_T=|\stackrel{}{p}_{T1}+\stackrel{}{p}_{T2}|/2`$ of the pairs. The non-diagonal elements are neglected. For the one-dimensional Gaussian fit case with $`\stackrel{}{Q}=Q_{inv}`$, the resolution corrected values of the fitted parameters are $`R_{inv}=7.30\pm 0.12`$ fm and $`\lambda =0.328\pm 0.009`$.
The results of the multi-dimensional fits are presented in Table 1 for the full 1995 data sample. A multi-dimensional analysis as a function of $`k_T`$, both with the standard 5-parameter fit and with the GYK fit is shown in Figs. 8, 9, and 10.
The $`R_{TS}`$ and $`R_L`$ parameters from the standard fit are found to be compatible respectively with $`R_T`$ and $`R_4`$ from the GYK fit. The cross term $`R_{outlong}^2`$ from the standard fit and $`v_L`$ from the GYK fit are compatible with 0. In a source undergoing a boost invariant expansion, the local rest frame coincides with the LCMS. Both the cross term and $`v_L`$ estimated in the LCMS are then expected to vanish . As this is the case, it suggests that the source seen within the acceptance of the experiment has a strictly boost invariant expansion. The strong decrease of the longitudinal radius $`R_L`$ or $`R_4`$ with $`k_T`$ compared to the behaviour of the transverse radii $`R_T`$, $`R_{TS}`$, $`R_{TO}`$ suggests a longitudinal expansion larger than the lateral expansion. The longitudinal radius $`R_L`$ is shown with a fit of the form 1/$`\sqrt{m_T}`$ with $`m_T`$=$`\sqrt{m_\pi ^2+k_T^2}`$ inspired by the hydrodynamical expansion model. Using $`R_L=\tau _0\sqrt{T_0/m_T}`$ with a freeze out temperature $`T_0`$ of 120-170 MeV/c, we may extract a freeze out time $`\tau _0`$ in the range of 7.5-8.9 fm/c. Finally, the $`R_0`$ parameter from the GYK fit, which reflects the duration of emission, is compatible with 0 for all $`k_T`$ bins, excluding a long-lived intermediate phase.
Two other experiments, NA49 and NA44, have studied charged particle interferometry in Pb+Pb collisions at CERN energies. The WA98 analysis is in good agreement with the NA49 results, when the comparison is made for the same mean $`y`$ range of 2.70, although WA98 has used identified $`\pi ^{}`$ while the NA49 analysis used unidentified negative particles. Only the $`R_0`$ parameter tends to be smaller in WA98. The direct comparison with the NA44 experiment is not possible because NA44 and WA98 do not have the same $`y`$ range. The smaller radii measured by NA44 can be explained by the larger $`y`$ range of its acceptance (3.1$`<y<`$4.1).
## 6 Conclusion
The analysis of the two-particle correlation of identified $`\pi ^{}`$ from central Pb+Pb collisions at 158 AGeV gives fitted radii of about 7 fm. This should be compared to the equivalent rms radius of the initial Pb nucleus of 3.2 fm, which indicates a large final state emission volume.
The one-dimensional correlation functions analyzed in terms of $`Q_{inv}`$ or $`Q_{space}`$ are not Gaussian. They are better represented by exponentials. This study is based on the tail of the distributions and not on the first bins which might be subject to systematic effects. One possible explanation is that this behaviour is due to resonance effects. Fitting Gaussians to these correlation functions may produce different results depending on the acceptance of the experimental setup.
The generalized Yano-Koonin analysis gives similar results to within the error bars as the standard 3-dimensional analysis in the LCMS.
The cross term $`R_{outlong}^2`$ is found to be compatible with 0 in the LCMS and the same is true of $`v_L`$ in the GYK fit. This suggests that the source undergoes a boost invariant expansion.
A clear dependence of the longitudinal radius parameter on $`k_T`$ is observed, suggesting a larger longitudinal than transverse expansion of the source. In addition the short duration of emission disfavours any long-lived intermediate phase.
Acknowledgements
We would like to thank the CERN-SPS accelerator crew for providing an excellent lead beam and the Laboratoire National Saturne for the loan of the magnet Goliath.
This work was supported jointly by the German BMBF and DFG, the U.S. DOE, the Swedish NFR, the Dutch Stichting FOM, the Swiss National Fund, the Humboldt Foundation, the Stiftung für deutsch-polnische Zusammenarbeit, the Department of Atomic Energy, the Department of Science and Technology and the University Grants Commission of the Government of India, the Indo-FRG Exchange Programme, the PPE division of CERN, the INTAS under contract INTAS-97-0158, the Polish KBN under the grant 2P03B16815, and ORISE. ORNL is managed by Lockheed Martin Energy Research Corporation under contract DE-AC05-96OR22464 with the U.S. Department of Energy. |
warning/0003/nucl-th0003042.html | ar5iv | text | # On dipole compression modes in nuclei
## Abstract
Isoscalar dipole strength distributions in spherical medium- and heavy-mass nuclei are calculated within random phase approximation (RPA) or quasiparticle RPA. Different Skyrme-type interactions corresponding to incompressibilities in the range 200 - 250 MeV are used. The results are discussed in comparison with existing data on isoscalar giant dipole resonances. Two main issues are raised, firstly the calculated giant resonance energies are somewhat higher than the observed ones, and secondly a sizable fraction of strength is predicted below 20 MeV which needs to be experimentally confirmed.
<sup>a</sup> Dipartimento di Fisica, Università degli Studi, and INFN, Via Celoria 16, I-20133 Milano (Italy)
<sup>b</sup> Institut de Physique Nucléaire, IN2P3-CNRS, F-91406 Orsay (France)
<sup>c</sup> Dipartimento di Fisica Teorica, Università degli Studi, and INFN, via P. Giuria 1, I-10125 Torino (Italy)
PACS: 24.30.Cz, 21.60.Jz
Keywords: giant resonances, nuclear incompressibility, HF-RPA,
quasiparticle RPA.
The isoscalar giant dipole resonance (ISGDR) is a compressional mode, like the well known isoscalar giant monopole resonance (ISGMR), and its energy is related to the nuclear incompressibility $`K_{\mathrm{}}`$. For this reason it has been studied both experimentally and theoretically for many years, as reviewed in . It can be associated with the operator
$$\widehat{D}=\underset{i=1}{\overset{A}{}}r_i^3Y_{1\mu }(\widehat{r}_i)$$
(1)
and it can be viewed as a non-isotropic compression mode.
Although some first indications about the energy location of this mode date back to the beginning of the eighties, more recent evidence of the ISGDR in <sup>208</sup>Pb has been reported from the 0<sup>0</sup> measurements of 200 MeV inelastic $`\alpha `$-scattering at the Indiana University cyclotron facility. Further evidence based on extensive angular distributions at near-0<sup>0</sup> angles has since come from 240 MeV inelastic $`\alpha `$-scattering experiments at Texas A&M University. The ISGDR strength has been extracted in a large number of nuclei using a multipole decomposition of the observed inelastic scattering spectra. The results for medium- and heavy-mass nuclei like <sup>90</sup>Zr, <sup>116</sup>Sn, <sup>144</sup>Sm and <sup>208</sup>Pb are reported in . The measured strength is spread over a wide energy range between 15 and 30 MeV and it is claimed to exhaust nearly 100% of the appropriate energy-weighted sum rule (EWSR). The values of the centroid energy E$`{}_{0}{}^{}m_1/m_0`$ are 26.3(4), 24.3(3), 23.0(3) and 20.3(2) MeV respectively in the four nuclei quoted above (the moments $`m_k`$ of the strength distribution are defined as $`m_k_n|<n|\widehat{D}|0>|^2(E_nE_0)^k`$).
Here, we report the ISGDR results calculated with effective Skyrme interactions, within self-consistent Hartree-Fock (HF) plus Random Phase Approximation (RPA) in the case of <sup>208</sup>Pb, and Hartree-Fock-BCS (HF-BCS) plus quasi-particle RPA (QRPA) for the other, non double-magic nuclei. The Skyrme forces used in this work are: SkP , SGII , SKM , SLy4 and SkI2 . They span a range of values of $`K_{\mathrm{}}`$ from 200 to 250 MeV. The HF mean field is first calculated in coordinate space, then the single-particle spectrum of occupied and unoccupied states is built by diagonalizing the mean field on a harmonic oscillator basis. Details of RPA calculations can be found in Ref. . The dimension of the 1particle-1hole (1p-1h) space is fixed by requiring the exhaustion of the RPA $`m_1`$ sum rule. In the HF-BCS calculations constant pairing gaps $`\mathrm{\Delta }`$ are introduced according to the usual 12 MeV/$`\sqrt{A}`$ parametrization. The QRPA matrix equations are solved with a procedure which parallels what has been said for RPA, with the two quasi-particle configurations replacing the 1p-1h ones. The method is the same as that of Ref. .
In the ISGDR problem, one has to face the question of the spurious state associated with the center-of-mass motion which carries the same quantum numbers $`J^\pi =1^{}`$. In a bona fide self-consistent RPA the spurious state would appear as an eigenstate at zero energy, exhausting the whole strength of the operator
$$\widehat{S}=\underset{i=1}{\overset{A}{}}r_iY_{1\mu }(\widehat{r}_i)$$
(2)
and orthogonal to all other physical states. However, in actual calculations the spurious state is at low but not zero energy because of small numerical inaccuracies and therefore, strength associated with the operator $`\widehat{S}`$ will be shared among the physical states. Starting from the actual RPA set of states $`|n^{}`$, we construct a new set of normalized states $`|n`$,
$$|n=𝒩_n(|n^{}\alpha _n|S),$$
(3)
where the state $`|S`$ is defined as
$$|S\widehat{S}|0,$$
(4)
$`|0`$ being the RPA vacuum. According to we associate to $`|S`$ the transition density
$$\alpha _S\frac{d\varrho _0}{dr}$$
(5)
where $`\varrho _0`$ is the HF ground state density. The state $`|n`$ is required to satisfy the condition $`n|\widehat{S}|0=0`$, i.e.,
$$𝑑rr^3(\delta \varrho _n^{}a_n\frac{d\varrho _0}{dr})=0,$$
(6)
where $`\delta \varrho _n^{}`$ is the transition density of the RPA state $`|n^{}`$ defined in the usual way. The problem of the spurious state normalization $`\alpha _S`$ is circumvented by the use of Eq. (6) since $`a_n\alpha _n\alpha _S`$ is well behaved (i.e., not divergent).
The difference between the strength distributions associated with the states $`|n^{}`$ and $`|n`$ is shown in the top-left corner of Fig. 1, for the typical case of <sup>208</sup>Pb with the force SGII. The strengths are essentially the same in the energy range which will be denoted “giant resonance (GR) region” (this range is evident from the plot but it is explicitly indicated in Tables 1 and 2). At lower energies, omitting the projection procedure can lead to a serious overestimation of the ISGDR strength. It is clear nevertheless from Fig. 1 that a non-negligible amount of non-spurious strength is present in the energy range which will be called “low-energy region”. This low-lying strength is due to 1 $`\mathrm{}\omega `$ excitations, which of course can contain strength associated with the $`\widehat{D}`$ operator.
Another way of eliminating the spurious strength is to keep the $`|n^{}`$ states and to replace the operator (1) by
$$\widehat{D}_{modif}=\underset{i=1}{\overset{A}{}}(r_i^3\eta r_i)Y_{1\mu }(\widehat{r}_i),$$
(7)
where $`\eta =\frac{5}{3}r^2`$. This prescription was derived and used in Ref. . Although the derivation is based on hydrodynamical-type arguments, one thus obtains strength distributions which are almost indistinguishable to those calculated with the present projection procedure. One may also note that in Ref. a different prescription was used for the subtraction of spurious strength resulting in an almost disappearance of strength in the low-energy region.
In Fig. 1 we also show center-of-mass corrected strength distributions for the other nuclei calculated with a typical interaction, namely SGII. The general features are: a) a large fraction of the strength lies in the GR region, and b) a non negligible amount of strength is in the low-energy region. The latter region contains about 20% of the ISGDR energy-weighted sum rule. These features are common to the results obtained with the other interactions. A more detailed analysis in terms of the moments $`m_0`$ and $`m_1`$ is reported in Table 1 for <sup>208</sup>Pb and all interactions, and in Table 2 for all 4 nuclei and the SGII interaction.
In comparison with the existing data, there are two main issues to be faced. First, there is a large discrepancy between predicted and measured GR energies, much larger than in all other GR cases. This is the more puzzling that the same model employed here was used successfully to describe the ISGMR in <sup>208</sup>Pb . Second, the calculations predict a sizeable amount of strength at low-energy, which needs to be experimentally confirmed . These features are common to the calculated strength distributions of the operator $`j_1(qr)Y_{1\mu }(\widehat{r})`$, which is a generalization of Eq.(1). In particular, they remain peaked at the same energies as the strength distribution of $`\widehat{D}`$ for values of $`q`$ up to 0.6 fm<sup>-1</sup>.
In Fig. 2 we show the predicted peak and centroid energies of the GR region for various nuclei as a function of $`K_{\mathrm{}}`$. The experimental values of $`E_0`$ for the GR region quoted above would be outside the figure, except for <sup>90</sup>Zr. The discrepancy appears very severe in Pb and Sn. In what follows we concentrate on Pb because it is the nucleus where the HF+RPA model should work better.
Earlier RPA calculations performed with the finite range Gogny interaction already found that the ISGDR energy was in the range of 26 MeV, in qualitative agreement with the present results and with Ref. . One might expect that effects beyond RPA, like the coupling to 2p-2h excitations would somewhat lower the centroid energy. However, the calculations of Ref. find a downward shift of less than 1 MeV. The ISGDR has also been calculated in the relativistic RPA approach in <sup>208</sup>Pb and <sup>144</sup>Sm and it is found that, for effective lagrangian parametrizations corresponding to $`K_{\mathrm{}}`$ in the range 200-270 MeV the energy of the ISGDR is of the order of 25 MeV. Thus, the question of understanding the observed values of $`E_0`$ is still open.
As for the low-energy region, our analysis of the configurations involved for instance in <sup>208</sup>Pb, shows that the strength comes from bound-to-bound neutron transitions like $`h_{9/2}i_{11/2}`$, $`i_{13/2}j_{15/2}`$ and, in some cases, $`f_{5/2}g_{7/2}`$. In the data reported in Ref. no low-lying strength is present. Further analysis of the same data is currently in progress , which may reveal the presence of isoscalar strength around the region of the isovector dipole.
In conclusion, we report in this paper HF+RPA and HF-BCS+QRPA calculations of the ISGDR in <sup>90</sup>Zr, <sup>112</sup>Sn, <sup>144</sup>Sm and <sup>208</sup>Pb nuclei. Two general features appear from the calculated strength distributions: some large resonance-type distribution of strength in the 110$`A^{1/3}`$ MeV energy region and some smaller, but still sizeable fraction of the strength below 20 MeV. These two characteristic features do not seem to agree quantitatively with the observation.
We thank D.H. Youngblood for helpful discussions about the experimental data analysis, and U. Garg for discussions. P.F.B. thanks IPN-Orsay for the warm hospitality during the time when this work was completed. |
warning/0003/nucl-th0003007.html | ar5iv | text | # Cluster states in nuclei as representations of a 𝑈(𝜈+1) group
## Abstract
We propose a description of cluster states in nuclei in terms of representations of unitary algebras $`U(\nu +1)`$, where $`\nu `$ is the number of space degrees of freedom. Within this framework, a variety of situations including both vibrational and rotational spectra, soft and rigid configurations, identical and non-identical constituents can be described. As an example, we show how the method can be used to study $`\alpha `$ clustering configurations in <sup>12</sup>C with point group symmetry $`𝒟_{3h}`$.
The purpose of this letter is to point out that the algebra $`U(\nu +1)`$, which has been suggested to be the spectrum generating algebra for a quantum mechanical problem with $`\nu `$ space degrees of freedom , might provide a framework for a unified description of cluster states in nuclei. We note that the main properties of clustering in nuclei are: (i) the softness of the cluster configuration which makes nuclei appear more like liquid structures rather than rigid molecular structures in which the constituents sit at some definite location in space; (ii) the near equality of vibrational and rotational energies which does not allow a clear-cut distinction between these two types of motion; (iii) the fact that the constituents are not point-like objects but particles with a spatial extent comparable to that of the overall structure and (iv) the fact that the constituents are often identical which implies that permutation symmetry must be imposed. A unified description of clustering in nuclei should be able to accomodate all these properties.
To illustrate the uselfuness of the algebra $`U(\nu +1)`$ in describing the variety of observed situations, we consider the specific case of a cluster composed of three particles (a description of two-body cluster configurations in nuclei in terms of $`U(4)`$ was suggested long ago and has been used to describe resonances in heavy ion scattering ). For a three-body problem, the number of space degrees of freedom (after removal of the center of mass) is $`\nu =3n3=6`$. (We do not consider in this article constituents with an internal structure. For such situation the algebraic structure must be enlarged to $`U(\nu +1)U(\mathrm{\Omega })`$ where $`\mathrm{\Omega }`$ is the number of internal degrees of freedom.) The space degrees of freedom can be taken as the Jacobi coordinates $`\stackrel{}{\rho }=(\stackrel{}{r}_1\stackrel{}{r}_2)/\sqrt{2}`$ and $`\stackrel{}{\lambda }=(\stackrel{}{r}_1+\stackrel{}{r}_22\stackrel{}{r}_3)/\sqrt{6}`$, where $`\stackrel{}{r}_i`$ $`(i=1,2,3)`$ are the coordinates of the three particles. The corresponding algebra is $`U(7)`$. The algebra of $`U(7)`$ is constructed by introducing two vector bosons $`b_\rho `$, $`b_\lambda `$ together with an auxiliary scalar boson $`s`$. It was introduced in where it was used to describe three-quark configurations in baryons. The 49 bilinear products of creation and annihilation operators generate the Lie algebra $`U(7)`$,
$`b_{\rho ,m}^{},b_{\lambda ,m}^{},s^{}`$ $``$ $`c_\alpha ^{}(m=0,\pm 1)(\alpha =1,\mathrm{},7)`$ (1)
$`𝒢`$ $`:`$ $`G_{\alpha \beta }=c_\alpha ^{}c_\beta (\alpha ,\beta =1,\mathrm{},7)`$ (2)
The creation and annihilation operators for vector bosons ($`b_{\rho ,m}^{}`$, $`b_{\lambda ,m}^{}`$ and $`b_{\rho ,m}`$, $`b_{\lambda ,m}`$) represent the second quantized form of the Jacobi coordinates and their canonically conjugate momenta, while the auxiliary scalar boson is introduced in order to construct the spectrum generating algebra. (The method of embedding the problem in a larger dimensional space is similar to that used in Kaluza-Klein theories of particle physics.) The energy levels can be obtained by diagonalizing the Hamiltonian
$`H`$ $`=`$ $`H_0+ϵ_ss^{}\stackrel{~}{s}ϵ_p(b_\rho ^{}\stackrel{~}{b}_\rho +b_\lambda ^{}\stackrel{~}{b}_\lambda )+u_0(s^{}s^{}\stackrel{~}{s}\stackrel{~}{s})u_1s^{}(b_\rho ^{}\stackrel{~}{b}_\rho +b_\lambda ^{}\stackrel{~}{b}_\lambda )\stackrel{~}{s}`$ (6)
$`+v_0\left[(b_\rho ^{}b_\rho ^{}+b_\lambda ^{}b_\lambda ^{})\stackrel{~}{s}\stackrel{~}{s}+s^{}s^{}(\stackrel{~}{b}_\rho \stackrel{~}{b}_\rho +\stackrel{~}{b}_\lambda \stackrel{~}{b}_\lambda )\right]`$
$`+{\displaystyle \underset{l=0,2}{}}c_l\left[(b_\rho ^{}\times b_\rho ^{}b_\lambda ^{}\times b_\lambda ^{})^{(l)}(\stackrel{~}{b}_\rho \times \stackrel{~}{b}_\rho \stackrel{~}{b}_\lambda \times \stackrel{~}{b}_\lambda )^{(l)}+4(b_\rho ^{}\times b_\lambda ^{})^{(l)}(\stackrel{~}{b}_\lambda \times \stackrel{~}{b}_\rho )^{(l)}\right]`$
$`+c_1(b_\rho ^{}\times b_\lambda ^{})^{(1)}(\stackrel{~}{b}_\lambda \times \stackrel{~}{b}_\rho )^{(1)}+{\displaystyle \underset{l=0,2}{}}w_l(b_\rho ^{}\times b_\rho ^{}+b_\lambda ^{}\times b_\lambda ^{})^{(l)}(\stackrel{~}{b}_\rho \times \stackrel{~}{b}_\rho +\stackrel{~}{b}_\lambda \times \stackrel{~}{b}_\lambda )^{(l)},`$
within the space of the totally symmetric representations $`[N]`$ of $`U(7)`$. The coefficients $`ϵ_s`$, $`ϵ_p`$, $`u_0`$, $`u_1`$, $`v_0`$, $`c_0`$, $`c_1`$, $`c_2`$, $`w_0`$ and $`w_2`$ parametrize the interactions. The Hamiltonian $`H`$ is the most general Hamitonian that preserves angular momentum and parity, transforms as a scalar under permutations (we consider here the case of three identical particles) and is at most quadratic (two-body interactions). Associated with the Hamiltonian $`H`$, there are transition operators, $`T`$. Electromagnetic transition rates and form factors can all be calculated by considering the matrix elements of the operator
$`T`$ $`=`$ $`\text{e}^{iq\beta D_{\lambda ,z}/X_D},`$ (7)
$`D_{\lambda ,z}`$ $`=`$ $`(b_\lambda ^{}\times \stackrel{~}{s}s^{}\times \stackrel{~}{b}_\lambda )_z^{(1)},`$ (8)
which is the algebraic image of the operator $`\mathrm{exp}(iqr_{3z})`$ obtained from the full operator $`_{i=1}^3e^{i\stackrel{}{q}\stackrel{}{r}_i}`$ by choosing the momentum transfer $`\stackrel{}{q}`$ in the $`z`$ direction and considering identical particles (the coefficient $`X_D`$ is a normalization factor).
The Hamiltonian of Eq. (6) has two dynamic symmetries corresponding to the breakings of $`U(7)`$ onto $`U(6)`$ and $`SO(7)`$
$$U(7)\{\begin{array}{c}U(6),\\ SO(7).\end{array}$$
(9)
When the Hamiltonian contains only Casimir operators of these chains, the eigenvalue problem can be solved in closed analytic form. The corresponding solutions describe two situations sometimes encountered in the three body problem: (i) six-dimensional vibrational spectra $`U(6)`$, and (ii) an unusual situation which we call $`\omega `$-unstable or $`SO(7)`$ limit. Both situations will be described in a longer publication. Here instead, as an example of application of the algebraic method, we discuss another situation that is appropriate to three particles at the vertices of an equilateral triangle. The spectrum of an equilateral triangle configuration can be obtained from the Hamiltonian of Eq. (6) by setting some coefficients equal to zero and taking specific combinations of others
$`H`$ $`=`$ $`H_0+\xi _1(s^{}s^{}b_\rho ^{}b_\rho ^{}b_\lambda ^{}b_\lambda ^{})(\stackrel{~}{s}\stackrel{~}{s}\stackrel{~}{b}_\rho \stackrel{~}{b}_\rho \stackrel{~}{b}_\lambda \stackrel{~}{b}_\lambda )`$ (13)
$`+\xi _2\left[(b_\rho ^{}b_\rho ^{}b_\lambda ^{}b_\lambda ^{})(\stackrel{~}{b}_\rho \stackrel{~}{b}_\rho \stackrel{~}{b}_\lambda \stackrel{~}{b}_\lambda )+4(b_\rho ^{}b_\lambda ^{})(\stackrel{~}{b}_\lambda \stackrel{~}{b}_\rho )\right]`$
$`+\xi _3(b_\rho ^{}\stackrel{~}{b}_\rho +b_\lambda ^{}\stackrel{~}{b}_\lambda )^{(1)}(b_\rho ^{}\stackrel{~}{b}_\rho +b_\lambda ^{}\stackrel{~}{b}_\lambda )^{(1)}`$
$`+\xi _4(b_\rho ^{}\stackrel{~}{b}_\lambda b_\lambda ^{}\stackrel{~}{b}_\rho )^{(0)}(b_\lambda ^{}\stackrel{~}{b}_\rho b_\rho ^{}\stackrel{~}{b}_\lambda )^{(0)}.`$
This spectrum does not correspond to a dynamic symmetry, since it cannot be written in terms of invariants of a chain of algebras originating from $`U(7)`$. However, an approximate expression for the energy levels can be obtained by making use of the method of intrinsic or coherent states (valid in the limit of large $`N`$). The energy eigenvalues are then given by
$$E(v_1,v_2^l,L,K,M)=E_0+Av_1+Bv_2+CL(L+1)+D(K\pm 2l)^2,$$
(14)
where $`A4N\xi _1`$, $`B2N\xi _2`$, $`C=\xi _3/2`$ and $`D=\xi _4/3`$. The quantum numbers have the following meaning: $`v_1`$, $`v_2`$ are vibrational quantum numbers; for three identical particles one of the vibration ($`v_1`$) is singly degenerate, while the other ($`v_2`$) is doubly degenerate; $`l=v_2,v_22,\mathrm{},1`$ or $`0`$ is the vibrational angular momentum of the doubly degenerate vibration; $`L`$ is the angular momentum, $`M`$ its projection on a laboratory fixed axis and $`K`$ its projection on a body fixed axis. We note the particular angular momentum composition of the rotation-vibration bands. The vibrationless ground state band $`(v_1,v_2^l)=(0,0^0)`$ has $`K=3n`$ ($`n=0,1,2,\mathrm{}`$) with $`L=0,2,4,\mathrm{}`$ for $`K=0`$ and $`L=K,K+1,K+2,\mathrm{}`$ for $`K0`$. The parity is given by $`P=()^K`$. The stretching vibration $`(1,0^0)`$ contains the same angular momenta $`L^P=0^+`$, $`2^+`$, $`3^{}`$, $`4^\pm ,\mathrm{}`$, as the ground state band, while the bending vibration $`(0,1^1)`$ has $`K=3n+1,3n+2`$ ($`n=0,1,2,\mathrm{}`$) with $`L=K,K+1,K+2,\mathrm{}`$. The angular momentum content of the bending vibration is then $`1^{}`$, $`2^\pm `$, $`3^\pm ,\mathrm{}`$. Since we do not consider the excitation of the $`\alpha `$ particles themselves, the wave functions describing the relative motion have to be symmetric. As a consequence, the relative sign in the last term of Eq. (14) is such that $`|K\pm 2l|=3m`$, a multiple of 3 . (The energy formula obtained from the Hamiltonian $`H`$ of Eq. (13) contains a Coriolis term which do not discuss here, since a detailed treatment of this term requires the use of the full Hamiltonian of Eq. (6), rather than the simplified form of Eq. (13)). In Fig. 1 we show the energy spectrum corresponding to Eq. (14). The importance of this figure is the particular nature of the rotation-vibration spectrum of a triangular configuration with $`𝒟_{3h}`$ symmetry. If a physical system is claimed to be composed of three identical structureless particles at the vertices of an equilateral triangle, then its spectrum must be as in Fig. 1. The algebraic framework produces this spectrum automatically by an appropriate choice of parameters.
Another consequence of using the compact algebra $`U(\nu +1)`$ as a spectrum generating algebra is that one can evaluate all observables in exact form. For example, by taking matrix elements of the operator $`T`$ between the eigenstates of $`H`$ obtained by matrix diagonalization, one can evaluate form factors. When the Hamiltonian has a dynamic symmetry these can be derived in closed form. Although the Hamiltonian of Eq. (13) does not correspond to a dynamic symmetry, the form factors can still be obtained in explicit form in the limit of large $`N`$. For transitions among the lowest states they are given by
$`F(0_1^+0_1^+;q)`$ $`=`$ $`j_0(q\beta ),`$ (15)
$`F(0_1^+2_1^+;q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{5}j_2(q\beta ),`$ (16)
$`F(0_1^+3_1^{};q)`$ $`=`$ $`i\sqrt{{\displaystyle \frac{35}{8}}}j_3(q\beta ),`$ (17)
$`F(0_1^+4_1^+;q)`$ $`=`$ $`{\displaystyle \frac{9}{8}}j_4(q\beta ),`$ (18)
$`F(0_1^+0_2^+;q)`$ $`=`$ $`\chi _1q\beta j_1(q\beta ),`$ (19)
$`F(0_1^+1_1^{};q)`$ $`=`$ $`i\chi _2{\displaystyle \frac{1}{2}}\sqrt{3}q\beta j_2(q\beta ).`$ (20)
Here $`q`$ is the momentum transfer and $`\beta `$ is the distance of the particles from the center (the first three form factors were already given in ). The last two form factors correspond to vibrational excitations. The coefficients $`\chi _1`$ and $`\chi _2`$ are proportional to the intrinsic matrix elements for each type of vibration ($`v_1`$ and $`v_2`$). Electromagnetic transition rates can be calculated from the $`B(EL)`$ values, which in turn can be obtained from the long wavelenght limit of the form factors. In the case in which the constituents of the cluster are extended objects (as in nuclei) the form factors and $`B(EL)`$ values can be obtained by folding the point-like distribution with the charge distribution (and eventually magnetic moment distribution) of the constituents. In the case of clusters composed of $`\alpha `$ particles, the folding can be done in a straightforward way, since the charge distribution of the $`\alpha `$ particle can be taken to a very good approximation as $`\mathrm{exp}(\alpha r^2)`$. The form factors for an extended distribution are then obtained from those in Eq. (20) by multiplying by $`\mathrm{exp}(q^2/4\alpha )`$. They are a crucial ingredient in understanding whether a cluster configuration is present or not. When $`N`$ is finite (the situation encountered in nuclei) the energy spectrum and form factors can be evaluated numerically using a computer program written by one of us . In this case, vibrational bands are no longer decoupled, but instead show an appreciable mixing between them and, as a result, the spectrum is considerably distorted from the energy formula of Eq. (14).
The formalism introduced here can be used to study cluster states in <sup>12</sup>C. It was suggested long ago that <sup>12</sup>C in its ground state can be viewed as three $`\alpha `$ particles at the vertices of an equilateral triangle (point group $`𝒟_{3h}`$). The experimental spectrum of <sup>12</sup>C is shown in Fig. 2, where it is compared with that given by Eq. (14). One can see that this spectrum is indeed similar (if not identical) to that of a triangular configuration. The crucial point is the sequence of angular momenta in the ground state rotational band: $`0^+`$, $`2^+`$, $`3^{}`$, $`4^+,\mathrm{}`$. This sequence is typical of a triangular configuration. A linear configuration would not have negative parity states, while a shell-model configuration would not have the $`3^{}`$ state as a member of the rotational band but rather as an octupole vibration, i.e. it would not form a rotational sequence with the $`0^+`$, $`2^+`$, $`4^+`$ states. However, the rotational spectrum does not follow precisely what expected from a triangular configuration (oblate top, $`D<0`$ in Eq. (14)) but it shows rather a spherical or slightly prolate top with $`𝒟_{3h}`$ symmetry. The spectrum also shows an excited $`0^+`$ state at 7.65 MeV and an excited $`1^{}`$ state at 10.84 MeV which could be interpreted as bandheads of the vibrational (stretching and doubly degenerate bending) excitations. Whether or not this is the case or rather those states represent other types of configurations, such as three $`\alpha `$ particles on a line as suggested by several authors, remains an open question. To settle this question uniquely one would have to identify the rotational sequences built on top of them which have a characteristic pattern for triangular configurations and another pattern for linear configurations. In particular the nature of the $`2^+`$ state at 11.16 MeV and $`2^{}`$ state at 11.83 MeV, which could form the rotational excitation of the doubly degenerate vibration, should be further investigated (the role played by the $`2^+`$ state in determining the cluster structure of <sup>12</sup>C has been emphasized before ). We have also calculated form factors and electromagnetic transition rates . All members of the ground rotational band are well described by Eq. (20), as well as the shape of the form factors leading to the $`0^+`$ state at 7.65 MeV and the $`1^{}`$ state at 10.84 MeV. This analysis will be presented in a forthcoming publication . The result of the simultaneous investigation of spectra, transition rates and form factors done within $`U(7)`$ is that an $`\alpha `$ clustering structure (albeit not a rigid one) with $`𝒟_{3h}`$ symmetry is a good description of the ground state configuration of <sup>12</sup>C. However, in order to make this conclusion stronger, we suggest to readdress the problem of $`\alpha `$ clustering in <sup>12</sup>C by a remeasurement of the properties of the high-lying states by $`(\alpha ,\alpha ^{})`$ and $`(e,e^{})`$ inelastic scattering. These experiments were done long ago and can benefit from new and improved techniques. We have predictions for all form factors, transition rates and energies of cluster states in the $`𝒟_{3h}`$ configuration. They can be obtained from us upon request.
In conclusion, we have proposed a description of cluster states in nuclei in terms of the group $`U(\nu +1)`$ and shown that within this algebraic structure one can describe many situations. In particular, for the three-body problem, one can recover the case of three particles at the vertices of a triangle, a configuration of interest in <sup>12</sup>C. We have shown that $`U(7)`$ contains the main properties of clustering in nuclei: the softness of the cluster configuration, the near equality of vibrational and rotational energies, the spatial extension of the constituents and the permutation symmetry. We can also describe the situation of three particles on a line (not discussed here) and of vibrational spectra, in other words the method is flexible enough that it can accomodate many situations encountered in nuclei. We have also constructed the algebra appropriate to four-body problems, $`U(10)`$, where additional geometric arrangements can occur, such as four particles at the vertices of a tetrahedron (point group $`𝒯_d`$) and used it to study cluster configurations in <sup>16</sup>O. In other words, all cluster structures up to four-body clusters can be studied with the algebraic method. The importance of using $`U(\nu +1)`$ for cluster states lies in the possibility of describing the variety of situations encountered in nuclei where clusters are not rigid structures but rather liquid like structures arising from the nature of the nucleon-nucleon force (spin-isospin) and the shell structure. The unitary algebra $`U(\nu +1)`$ can also be of interest in the description of other quantum mechanical systems with non-rigid structure, such as atomic clusters, floppy molecules, and trimers making the method of broad applicability to a large class of problems.
This work was supported in part by DGAPA-UNAM under project IN101997, by CONACyT under project 32416-E, and by D.O.E. Grant DE-FG02-91ER40608. |
warning/0003/math0003194.html | ar5iv | text | # Non-unitary set-theoretical solutions to the Quantum Yang-Baxter Equation
## 1 Introduction
In this paper we study set-theoretical solutions to the Quantum Yang-Baxter equation, i.e. permutations $`R:X\times XX\times X`$ with $`X`$ being a non-empty set such that
$$R^{12}R^{13}R^{23}=R^{23}R^{13}R^{12}inAut(X\times X\times X).$$
In the above $`R^{12}`$, $`R^{13}`$, $`R^{23}Aut(X\times X\times X)`$ stand for $`R`$ acting in 1,2; 1,3; and 2,3 components of $`X\times X\times X`$ correspondingly. The idea to consider set-theoretical solutions first appeared in \[Dr\]. Later on, Etingof, Schedler and the author \[ESS\] studied set-theoretical solutions to the Quantum Yang-Baxter equation which satisfied additional properties of unitarity and nondegeneracy (crossing symmetry). In particular, \[ESS\] contained the classification of nondegenerate unitary set-theoretical solutions to QYBE in group theoretical terms as well as numerous classes of examples of such solutions. Subsequently, Lu, Yan and Zhu in \[LYZ\] showed that many of the constructions from \[ESS\] hold in a more general case of nondegenerate but not necessarily unitary set-theoretical solutions to QYBE.
Following \[ESS\] and \[LYZ\], we develop a theory of nondegenerate set-theoretical solutions to QYBE. Particularly, we show that the unitarity condition that was used in \[ESS\] for group theoretical characterization of unitary nondegenerate set-theoretical solutions to QYBE can be dropped. We give a group theoretical characterization of general set-theoretical nondegenerate solutions to QYBE in Theorem 2.7. We also introduce injective solutions, study their properties, and show that there is a combinatorial criterion (Theorem 2.9) describing the class of injective nondegenerate solutions to QYBE. Injectivity property, which is a generalization of involutivity, is important for studying affine solutions. In particular the classification of unitary affine solutions given in \[ESS\] can be generalized to include injective solutions.
It was shown in \[ESS\] that the structure group of a nondegenerate unitary set-theoretical solution to QYBE on a set with N elements always has an abelian subgroup of finite index and of rank N. We compute (Theorem 2.10) the rank of a finite index abelian subgroup of the structure group for an arbitrary nondegenerate finite solution and show that this rank never exceeds N, with the equality taking place only in the unitary case.
In the second part of the paper we discuss the applications of the developed theory to examples. In particular, we classify affine solutions on an abelian group. It is proved that injective affine solutions are obtained from the representations of the algebra generated by invertible elements $`p`$, $`q`$, $`z`$ subject to $`pq=qp`$ and $`z^2z(p+q)+pq=0`$.
Acknowledgments. The author is thankful to his advisor Pavel Etingof for a valuable exchange of ideas and guidance throughout the work on the paper, and to Lu, Yan, Zhu for making their work available before publication.
## 2 Structure groups
### 2.1 Construction of the structure groups
Let $`X`$ be a nonempty set and $`S:X\times XX\times X`$ a bijective map. We call a pair $`(X,S)`$ a braided set if the following braiding condition holds in $`X\times X\times X`$:
$$S_1S_2S_1=S_2S_1S_2,$$
(2.1)
where $`S_1=S\times id`$, $`S_2=id\times S`$.
Remark. Consider the map $`R:X\times XX\times X`$ given by $`R=\sigma S`$, where $`\sigma (x,y)=(y,x)`$ for $`x,yX`$. Then $`(X,S)`$ is a braided set if and only if $`R`$ satisfies the Quantum Yang-Baxter equation.
It is useful to associate with a braided set $`(X,S)`$ two groups $`G_X`$ and $`A_X`$.
###### Definition 1
Define the group $`G_X`$ as the group generated by the elements of $`X`$ subject to the relations $`xy=y_1x_1`$ if $`S(x,y)=(y_1,x_1)`$, where $`x,yX`$. We call $`G_X`$ the structure group of the braided set $`(X,S)`$.
###### Definition 2
Define the group $`A_X`$ as the group generated by elements of $`X`$ subject to relations $`x_1y=y_2x_1`$, where $`x,yX`$ and $`x_1,y_2X`$ are defined out of relations $`S(x,y)=(y_1,x_1),S(y_1,x_1)=(x_2,y_2)`$. We call $`A_X`$ the derived structure group of the braided set $`(X,S)`$.
We introduce the maps $`g:X\times XX`$ and $`f:X\times XX`$ as components of $`S`$, i.e. for $`x,yX`$
$$S(x,y)=(g_x(y),f_y(x)).$$
###### Definition 3
(i) We call a set $`(X,S)`$ nondegenerate if $`g_x(y)`$ is a bijective function of $`y`$ for fixed $`x`$ and $`f_y(x)`$ is a bijective function of $`x`$ for fixed $`y`$. (ii) We call a set $`(X,S)`$ involutive if $`S^2=id_{X^2}`$. A braided set which is involutive will be called symmetric.
In particular, for a symmetric set $`(X,S)`$ we see that the group $`A_X`$ is the free abelian group generated by elements of $`X`$. Note that the properties of involutivity and nondegeneracy are equivalent to corresponding properties of unitarity and crossing symmetry for the map $`R=\sigma S`$ \[ESS\].
Recall that the braid group $`B_n`$ for $`n2`$ is generated by elements $`b_i,1in1`$, with defining relations
$$b_ib_j=b_jb_i,|ij|>1,b_ib_{i+1}b_i=b_{i+1}b_ib_{i+1},$$
and recall that the symmetric group $`S_n`$ is the quotient of $`B_n`$ by the relations $`b_i^2=1`$. The following obvious proposition explains our terminology. Let $`S_n^{ii+1}:X^nX^n`$ be defined as $`S_n^{ii+1}=id_{X^{i1}}\times S\times id_{X^{ni1}}`$.
###### Theorem 2.1 (\[ESS\])
(i) The assignment $`b_iS_n^{ii+1}`$ extends to an action of $`B_n`$ on $`X^n`$ ($`n3)`$ if and only if $`(X,S)`$ is a braided set.
(ii) The assignment $`b_iS_n^{ii+1}`$ extends to an action of $`S_n`$ on $`X^n`$ ($`n3`$) if and only if $`(X,S)`$ is a symmetric set.
###### Definition 4
The action of Theorem 2.1 is called the twisted action of $`B_n`$ (or $`S_n`$) given by $`S`$.
The following result shows how a nondegenerate braided set $`(X,S)`$ gives rise to two actions of the structure group $`G_X`$ on the set $`X`$.
###### Theorem 2.2
(\[ESS\]) Suppose that $`(X,S)`$ is nondegenerate. Then $`(X,S)`$ is a braided set if and only if the following conditions are simultaneously satisfied:
(i) the assignment $`xf_x`$ is a right action of $`G_X`$ on $`X`$;
(ii) the assignment $`xg_x`$ is a left action of $`G_X`$ on $`X`$;
(iii) the linking relation
$$f_{g_{f_y(x)}(z)}(g_x(y))=g_{f_{g_y(z)}(x)}(f_z(y))$$
holds.
Proof: Conditions (i)-(iii) are exactly the three components of the braid relation (2.1). $`\mathrm{}`$
In light of the above proposition it makes sense to introduce the notations $`xy=g_x(y)`$ and $`yx=f_y^1(x)`$ for $`x,yX`$. Then if $`(X,S)`$ is a nondegenerate braided set, we can extend $``$ and $``$ to left actions of $`G_X`$ on $`X`$. We will denote the actions of an element $`gG_X`$ on an element $`xX`$ by $`gx`$ and $`gx`$ correspondingly.
From now on we always assume $`(X,S)`$ to be a nondegenerate braided set. Sometimes we refer to nondegenerate braided sets as to ”solutions”, meaning that $`S`$ is a solution to braid equation (2.1).
Define $`\varphi :X\times XX`$ by
$$\varphi (y,x)=x^1((yx)y)$$
(2.2)
and $`S^{}:X\times XX\times X`$ as $`S^{}(x,y)=(\varphi (y,x),x)`$.
###### Theorem 2.3
(i) $`\varphi (y,z)`$ is $`G_X`$-invariant w.r.t. \*-action, i.e. for $`gG_X`$, $`g\varphi (y,z)=\varphi (gy,gz)`$.
(ii) $`(X,S^{})`$ is a nondegenerate braided set. We call this set the derived braided set or the derived solution.
(iii) The structure group of the derived solution is the derived structure group.
(iv) For each integer $`n2`$ there exists a bijection $`J_n:X^nX^n`$ such that $`J_nS_n^{ii+1}J_n^1=(S^{})_n^{ii+1}`$, where $`S_n^{ii+1}`$ is the same as in Theorem 2.1. In this way, twisted actions of $`B_n`$ given by $`S`$ and $`S^{}`$ are conjugate.
We remark that the statement (iv) of Theorem 2.3 was proved in \[ESS\] (cf. Prop. 1.7) for unitary solutions and in \[LYZ\] (cf. Th. 3) for injective solutions that are defined in Section 2.3.
Proof:
It is easy to see that (iv) implies (ii). Statement (iii) follows from definitions of structure group, derived structure group and $`\varphi `$. Let us show that (i) implies (iv). Define $`J_n`$ inductively as $`J_1=id_X`$, $`J_n=Q_n(J_{n1}\times id_X)`$, where $`Q_n:X^nX^n`$ is defined as $`Q_n(x_1,\mathrm{},x_n)=(x_n^1x_1,\mathrm{},x_n^1x_{n1},x_n)`$. We prove formula $`J_nS_n^{ii+1}J_n^1=(S^{})_n^{ii+1}`$ by induction on $`n`$. For $`n=2`$ (induction base) the relation follows directly from definition of $`\varphi `$. Suppose the relation holds for $`n=k`$, let us prove it holds for $`n=k+1`$. Since $`J_{k+1}=Q_{k+1}(J_k\times id_X)`$ and $`Q_{k+1}`$ commutes with $`(S^{})_{k+1}^{ii+1}`$ for $`i<k`$ (by (i)) the relation is true for $`i<k`$. So, it remains to prove it for $`i=k`$ when it becomes identical to the relation of the induction base.
Now we have to show that (i) is true. It is enough to check that for every $`t,y,zX`$, $`t^1\varphi (y,z)=\varphi (t^1y,t^1z)`$ or, equivalently by (2.2):
$$t^1(z^1((yz)y))=(t^1z)^1(((t^1y)(t^1z))(t^1y)).$$
(2.3)
The linking relation of Theorem 2.2 states that for $`x,y,tX`$
$$((y^1x)t)^1(xy)=((yt)^1x)(t^1y)$$
(2.4)
holds. If we substitute $`x=yz`$ in relation (2.4) we can rewrite its right hand side as follows:
$$((yt)^1(yz))(t^1y)=((yt)^1(yt(t^1z)))(t^1y)$$
where the product $`yt`$ is the product of elements $`y,t`$ in group $`G_X`$, thus $`yt=(yt)(t^1y)`$ and
$$((yt)^1(yt(t^1z)))(t^1y)=((t^1y)(t^1z))(t^1y).$$
In this way, relation (2.4) is equivalent to relation (2.5) below:
$$(zt)^1((yz)y)=((t^1y)(t^1z))(t^1y).$$
(2.5)
Substituting, (2.5) into the right side of (2.3) we get
$$(t^1z)^1(((t^1y)(t^1z))(t^1y))=((t^1z)^1(zt)^1)((yz)y).$$
Using the relation $`t^1z^1=(t^1z)^1(zt)^1`$ in group $`G_X`$ we verify that invariance relation (2.3) is true. The theorem is proved. $`\mathrm{}`$
Note that by construction we have two natural maps $`\psi _G:XG_X`$ and $`\psi _A:XA_X`$ that are not necessarily embeddings (see Examples below). Theorem 2.4 shows that the latter map is $`G_X`$-invariant with respect to a suitable $`G_X`$-module structure on $`A_X`$. Define the group $`Permut(X)`$ as the group of all permutations of the set $`X`$. Then both $``$ and $``$ define homomorphisms from $`G_X`$ to $`Permut(X)`$. Denote by $`Aut_X(A_X)`$ the group of automorphisms of $`A_X`$ that map the generating set $`\psi _A(X)`$ onto itself.
###### Theorem 2.4
For a nondegenerate braided set $`(X,S)`$ the group homomorphism $`:G_XPermut(X)`$ can be uniquely lifted to the homomorphism $`\widehat{}:G_XAut_X(A_X)`$ such that for any $`gG_X,xX`$ $`\psi _A(gx)=g\widehat{}\psi _A(x)`$.
Proof:
It is clear that the map $`:G_XPermut(X)`$ can be uniquely extended to the homomorphism from the group $`G_X`$ to the group of automorphisms of the free group generated by $`X`$. So, it is enough to show that such an extension respects the generating relations in group $`A_X`$. This immediately follows from statements (i), (iii) of Theorem 2.3. $`\mathrm{}`$
Example\[Dr\] Let $`X`$ be a union of conjugacy classes in a group $`G`$, i.e. $`gXg^1=X`$ for any $`gG`$. Define $`S:X\times XX\times X`$ as $`S(x,y)=(xyx^1,x)`$ for $`x,yX`$. It is easy to see that $`(X,S)`$ is a braided nondegenerate set. It was called the conjugate solution in \[LYZ\]. Clearly, this solution coincides with its own derived solution - therefore $`G_X=A_X`$. The action of $`G_X`$ on $`A_X`$ is easily seen to be trivial, i.e. $`g\widehat{}h=h`$ for any $`gG_X,hA_X`$. The embedding $`i:XG`$ can be extended to a homomorphism $`I:G_XG`$ that maps elements of $`\psi _G(X)`$ onto $`X`$, thus the map $`\psi _G:XG_X`$ is injective. So $`(X,S)`$ is an injective solution (cf. Definition 8).
### 2.2 Bijective 1-cocycle
Recall that we started with a nondegenerate braided set $`(X,S)`$ and constructed two groups $`G_X`$, $`A_X`$ and the action $`\widehat{}`$ of $`G_X`$ on $`A_X`$. In the future we drop $`\widehat{}`$ in $`\widehat{}`$ and denote the action of $`gG_X`$ on $`hA_X`$ from Theorem 2.4 simply by $`gh`$. Theorem 2.5 is the main step towards classification of nondegenerate braided sets in group theoretical terms.
###### Definition 5
Suppose a group $`G`$ acts on a group $`A`$ by automorphisms meaning that there is a homomorphism from $`G`$ to $`Aut(A)`$. Denote the product of elements $`g_1,g_2`$ in $`G`$ by $`g_1g_2`$ and the product of elements $`a_1,a_2`$ in $`A`$ by $`a_1a_2`$. We call a map $`\pi :GA`$ a 1-cocycle if for any $`g_1,g_2G`$
$$\pi (g_1g_2)=(g_2^1\pi (g_1))\pi (g_2)$$
(2.6)
where $``$ stands for the action of $`G`$ on $`A`$.
###### Theorem 2.5
(cf. \[LYZ\], Th. 2 and \[ESS\], Prop. 2.5) For a nondegenerate braided set $`(X,S)`$ there exists a unique bijective 1-cocycle $`\pi :G_XA_X`$ such that $`\pi \psi _G=\psi _A`$ on $`X`$.
Remark. In \[LYZ\] the authors introduced on $`G_X`$ another group structure. It follows from our and their results that $`G_X`$ with the new group structure is isomorphic to $`A_X`$ via $`\pi :G_XA_X`$.
Proof:
Construction of the 1-cocycle.
Let us construct the 1-cocycle $`\pi :G_XA_X`$. Consider the semidirect product $`G_XA_X`$. The group $`G_XA_X`$ consists of pairs $`(g,h)`$, $`gG_X,hA_X`$ with the group operation given by the formula
$$(g,h)(g^{},h^{})=(gg^{},((g^{})^1h)h^{}).$$
Define the map $`s:XG_XA_X`$ via the formula $`s(x)=(\psi _G(x),\psi _A(x))`$ for any $`xX`$. We claim that there exists a group homomorphism $`\overline{s}:G_XG_XA_X`$ such that $`\overline{s}\psi _G=s`$. Indeed, since $`\psi _G(X)`$ generates the whole $`G_X`$ to show that $`\overline{s}`$ exists it is necessary and sufficient to check that $`s`$ respects the relations in $`G_X`$, i.e. that for $`x,yX`$ $`s(x)s(y)`$ is equal to $`s(xy)s(y^1x)`$ in group $`G_XA_X`$. When we formally multiply terms out, the above condition transforms into the relation
$$(xy,(y^1x)y)=((xy)(y^1x),((y^1x)^1(xy))(y^1x)),$$
which coincides componentwise with the defining relations in groups $`G_X`$, $`A_X`$. Let us define the projection $`p:G_XA_XA_X`$ by the formula $`p(g,h)=h`$. Introduce $`\pi :G_XA_X`$ as $`\pi =p\overline{s}`$. Clearly, $`\pi `$ is a 1-cocycle. The tricky part is to show that $`\pi `$ is bijective.
Proof that $`\pi `$ is bijective.
###### Lemma 1
(i) For $`x_1,x_2X`$, $`gG_X`$ if $`\psi _G(x_1)=\psi _G(x_2)`$ then $`\psi _G(gx_1)=\psi _G(gx_2)`$ and $`\psi _G(gx_1)=\psi _G(gx_2)`$.
(ii) For $`xX`$ $`\psi _G((xx)x)=\psi _G(xx)`$ and $`\psi _G((x^1x)^1x)=\psi _G(x^1x)`$.
(iii) For $`xX`$ $`\psi _G(x)=\psi _G((xx)^1(xx))=\psi _G((x^1x)(x^1x))`$.
To prove the lemma we notice that $`S(xx,x)=((xx)x,x)`$ therefore $`\psi _G((xx)x)\psi _G(x)=\psi _G(xx)\psi _G(x)`$ and $`\psi _G((xx)x)=\psi _G(xx)`$. Similarly $`S(x,x^1x)=(x,(x^1x)^1x)`$ implies that $`\psi _G((x^1x)^1x)=\psi _G(x^1x)`$. So statement (ii) of the lemma is proved. It is clear that statement (iii) follows from (i) and (ii) thus it remains to prove (i).
Let us show that $`\psi _G(x_1)=\psi _G(x_2)`$ implies that $`\psi _G(gx_1)=\psi _G(gx_2)`$. For a fixed $`zX`$ it is enough to reason that $`\psi _G(x_1)=\psi _G(x_2)`$ if and only if $`\psi _G(z^1x_1)=\psi _G(z^1x_2)`$. Since $`S(x_1,z)=(x_1z,z^1x_1)`$ and $`S(x_2,z)=(x_2z,z^1x_2)`$ we see that
$$\psi _G(x_1)\psi _G(z)=\psi _G(x_1z)\psi _G(z^1x_1)$$
(2.7)
and
$$\psi _G(x_2)\psi _G(z)=\psi _G(x_2z)\psi _G(z^1x_2).$$
(2.8)
Suppose $`\psi _G(x_1)=\psi _G(x_2)`$ then since $``$ is an action of $`G_X`$ on $`X`$ we immediately see that $`\psi _G(x_1z)=\psi _G(x_2z)`$ therefore equations (2.7)-(2.8) imply that $`\psi _G(z^1x_1)=\psi _G(z^1x_2)`$. Conversely, assume that $`\psi _G(z^1x_1)=\psi _G(z^1x_2)`$. Our task is to prove that $`\psi _G(x_1)=\psi _G(x_2)`$. We use the statement (iv) of Theorem 2.3 in the following form: $`S=J_2^1S^{}J_2`$, where $`J_2(x,y)=(y^1x,y)`$. We notice that $`S^{}J_2(x_1,z)=(\varphi (z,z^1x_1),z^1x_1)`$ and $`S^{}J_2(x_2,z)=(\varphi (z,z^1x_2),z^1x_2)`$. Since $`\varphi (z,z^1x_1)`$ is the action of $`\psi _A(z^1x_1)A_X`$ on $`zX`$ and $`\psi _A=\pi \psi _G`$ we see that $`\varphi (z,z^1x_1)=\varphi (z,z^1x_2)`$. In this way, first components of $`S(x_1,z)=J_2^1S^{}J_2(x_1,z)`$ and $`S(x_2,z)=J_2^1S^{}J_2(x_2,z)`$ coincide, i.e. $`x_1z=x_2z`$. So the equations (2.7)-(2.8) imply that $`\psi _G(x_1)=\psi _G(x_2)`$. In a similar fashion one can show that $`\psi _G(x_1)=\psi _G(x_2)`$ if and only if $`\psi _G(zx_1)=\psi _G(zx_2)`$. Lemma is proved. $`\mathrm{}`$.
We aim to construct the map $`h:A_XG_X`$ inverse to $`\pi `$. At first we define $`h`$ on $`F_X`$ \- the free group generated by $`X`$ and then show that it descends to $`A_X`$. We note that $`G_X`$ acts on $`F_X`$ via $`fgf`$ for $`fF_X,gG_X`$ \- the action induced from \*-action of $`G_X`$ on $`X`$. For $`xXF_X`$ define
$$h(x)=\psi _G(x),h(x^1)=(\psi _G(x^1x))^1$$
Notice that $`x^1x`$ is the same as $`\psi _G(x)^1x`$. For an element $`f=x_1\mathrm{}x_kF_X`$ of length $`k`$ define inductively $`h(x_1\mathrm{}x_k)=h(h(x_2\mathrm{}x_k)x_1)h(x_2\mathrm{}x_k)`$, where for each $`i`$, $`x_iXF_X`$ or $`x_i^1XF_X`$. In the above $`f=x_1\mathrm{}x_k`$ was the minimum decomposition of $`fF_X`$ and $`k`$, the length of f, is the number of elements in such a decomposition. The only element of length 0 in $`F_X`$ is identity $`e`$ and we put $`h(e)=1G_X`$. We claim that for $`a,bF_X`$
$$h(ab)=h(h(b)a)h(b).$$
(2.9)
Indeed, we can verify (2.9) by induction on the length of $`a`$. Suppose we know that (2.9) holds for for all elements $`a`$ of length $`k`$ and want to check that $`h(ayb)=h(h(b)(ay))h(b)`$ for $`yXF_X`$ or $`y^1XF_X`$. We simplify $`h(h(b)(ay))h(b)`$ as follows:
$`h(h(b)(ay))h(b)=h((h(b)a)(h(b)y))h(b)`$
$`=h(h(h(b)y)(h(b)a))h(h(b)y)h(b).`$
On the other hand, we know that
$`h(ayb)=h(h(yb)a)h(yb)=h(h(h(b)y)h(b)a)h(h(b)y)h(b)`$
$`=h(h(h(b)y)(h(b)a))h(h(b)y)h(b).`$
We see that $`h(ayb)=h(h(b)(ay))h(b)`$. In this way, we just need to check the induction base, namely that $`h(yb)=h(h(b)y)h(b)`$. There are two cases to consider:
1. $`Length(yb)=1+Length(b)`$
2. $`Length(yb)=Length(b)1`$, i.e. $`b=y^1b^{}`$ and $`Length(b)=Length(b^{})+1`$
In the first case the formula $`h(yb)=h(h(b)y)h(b)`$ follows from definition of $`h`$. In the second case without loss of generality assume that $`yX`$. Then, $`h(b)=h(y^1b^{})=h((h(b^{})y)^1)h(b^{})`$ and letting $`h(b^{})y=z`$ we get
$`h(h(b)y)h(b)=h(h((h(b^{})y)^1)h(b^{})y)h((h(b^{})y)^1)h(b^{})`$
$`=h(h(z^1)z))h(z^1)h(b^{})=\psi _G(\psi _G(z^1z)^1z)\psi _G(z^1z)^1h(b^{})`$
$`=(byLemma)\psi _G((z^1z)\psi _G(z^1z)^1h(b^{})=h(b^{})=h(yb).`$
Now we are going to check that $`h`$ descends to $`A_X`$. Recall that $`A_X`$ is given by generators - elements of $`X`$ and relations $`(z^1((yz)y))z=zy`$ for each $`y,zX`$. We see that
$$h((z^1((yz)y))z)=\psi _G((yz)y)h(z)=h(zy).$$
(2.10)
If we define $`p(y,z)=(z^1((yz)y))zF_X`$ and $`q(y,z)=zyF_X`$ then for any $`gG_X`$, $`h(g(p(y,z)^1q(y,z)))=h(g(q(y,z)^1p(y,z)))=1`$.
Indeed, we checked in the proof of Theorem 2.4 that the defining relations of group $`A_X`$ are $`G_X`$ -invariant with respect to * action, i.e. $`gp(y,z)=p(gy,gz)`$ and $`gq(y,z)=q(gy,gz)`$. Formula (2.10) states that $`h(p(y,z))=h(q(y,z))`$, therefore
$`h(p(y,z)^1q(y,z))=h(h(q(y,z))p(y,z)^1)h(q(y,z))`$
$`=h(h(p(y,z))p(y,z)^1)h(p(y,z))=h(p(y,z)^1p(y,z))=1.`$
Thus $`h(g(p(y,z)^1q(y,z)))=h(p(gy,gz)^1q(gy,gz))=1`$. Similarly $`h(g(q(y,z)^1p(y,z)))=1`$.
In order to check that $`h:F_XG_X`$ descends to the map $`h:A_XG_X`$ it is enough to check that for any $`a,bF_X,y,zX`$ $`h(ap(y,z)^1q(y,z)b)=h(aq^1(y,z)p(y,z)b)=h(ab)`$ holds. By formula (2.9)
$`h(ap^1(y,z)q(y,z)b)=h(h(p(y,z)^1q(y,z)b)a)`$
$`h(p(y,z)^1q(y,z)b),`$
so it is enough to show that $`h(p^1(y,z)q(y,z)b)=h(b)`$. But $`h(p^1(y,z)q(y,z)b)=h(h(b)(p^1(y,z)q(y,z)))h(b)=h(b)`$ holds. Similar reasoning allows to check that $`h(aq^1(y,z)p(y,z)b)=h(ab)`$. In this way we constructed the map $`h:A_XG_X`$ such that
1. $`h(ab)=h(h(b)a)h(b)`$
2. $`h\psi _A=\psi _G`$
Condition 1 above implies that both $`\pi h:A_XA_X`$ and $`h\pi :G_XG_X`$ are group homomorphisms while Condition 2 implies that these homomorphisms are identities if restricted to corresponding generating sets $`\psi _A(X),\psi _G(X)`$, thus $`\pi h=id_{A_X},h\pi =id_{G_X}`$. Theorem 2.5 is proved. $`\mathrm{}`$
It turns out then that the groups $`G_X`$ and $`A_X`$ both have abelian subgroups of finite index when $`X`$ is finite. Namely, define $`\mathrm{\Gamma }=\{gG_X|gx=x,gx=xforallxX\}`$. In other words, $`\mathrm{\Gamma }`$ is the intersection of the kernels of left and right actions from Theorem 2.2.
###### Theorem 2.6
(cf. Section 2.5 in \[ESS\] and Prop. 6 in \[LYZ\])
(i) $`\pi (\mathrm{\Gamma })`$ is a normal $`G_X`$-invariant (w.r.t. to $``$-action) subgroup lying in the center of $`A_X`$. $`\mathrm{\Gamma }`$ is a normal abelian subgroup in $`G_X`$, and $`\pi :\mathrm{\Gamma }\pi (\mathrm{\Gamma })`$ is an isomorphism.
(ii) The 1-cocycle $`\pi :G_XA_X`$ can be factored out through $`\mathrm{\Gamma }`$ giving rise to the bijective 1-cocycle $`\overline{\pi }:G_X/\mathrm{\Gamma }A_X/\pi (\mathrm{\Gamma })`$.
(iii) If X is finite then both $`G_X/\mathrm{\Gamma }`$ and $`A_X/\pi (\mathrm{\Gamma })`$ are finite groups.
Proof:
(i) Let us show that $`\pi (\mathrm{\Gamma })`$ is $`G_X`$-invariant and central subgroup in $`A_X`$. From the defining relations in $`G_X`$ we see that $`\psi _G(xy)=\psi _G(x)\psi _G(y)\psi _G(y^1x)^1`$ for all $`x,yX`$. In fact, one can easily check that the above formula can be generalized for $`xG_X`$, $`yX`$.
###### Lemma 2
For $`gG_X,yX`$
$$\psi _G(gy)=g\psi _G(y)\pi ^1(\psi _G(y)^1\pi (g))^1.$$
Proof of Lemma 2. Let us show that if the statement of the Lemma holds for some $`gG_X`$ then it holds for $`g^1`$. We need to check that
$$g^1\psi _G(gy)\pi ^1(\psi _G(gy)^1\pi (g^1))^1=y.$$
(2.11)
Notice that $`\pi (g^1)=(g\pi (g))^1`$ therefore
$$\psi _G(gy)^1\pi (g^1)=((\psi _G(gy)^1g)\pi (g))^1.$$
Since the statement of the Lemma holds for $`g`$, we have
$`\pi ^1((\psi _G(gy)^1g)\pi (g))^1=\pi ^1((\pi ^1(\psi _G(y)^1\pi (g))\psi _G(y)^1)\pi (g))^1`$
$`=\pi ^1(\pi ^1(\psi _G(y)^1\pi (g))(\psi _G(y)^1\pi (g)))^1`$
$`=\pi ^1(\psi _G(y)^1\pi (g)).`$
This implies the validity of equation (2.11). Similar, a simple computation shows that if the statement of the lemma is true for $`g=g_1`$ and $`g=g_2`$ then it is true for $`g=g_1g_2`$. In this way since by definition of $`G_X`$ Lemma is true for $`g\psi _G(X)`$ we prove the Lemma. $`\mathrm{}`$
In particular, for $`g\mathrm{\Gamma },y\psi _G(X)`$ one has $`y=gy(\pi ^1(y^1\pi (g)))^1`$. Thus,
$$y^1gy=\pi ^1(y^1\pi (g)).$$
(2.12)
The condition (2.12) implies that $`\pi \mathrm{\Gamma }`$ belongs to the center of $`A_X`$. Indeed, applying $`\pi `$ to relation $`y^1g=\pi ^1(y^1\pi (g))y^1`$ we get that $`\pi (y^1)`$ commutes with $`\pi (g)`$. Moreover, for $`g\mathrm{\Gamma },hG_X`$ 1-cocycle condition implies that
$$\pi ^1(\pi (h)\pi (g))=hg.$$
(2.13)
So, in particular the product of elements in $`\pi (\mathrm{\Gamma })`$ is in $`\pi (\mathrm{\Gamma })`$ and $`\pi `$ restricted to $`\mathrm{\Gamma }`$ becomes an isomorphism between $`\mathrm{\Gamma }`$ and $`\pi (\mathrm{\Gamma })`$.
(ii) The relation (2.13) shows that $`\pi `$ can be lifted to $`\overline{\pi }:G_X/\mathrm{\Gamma }A_X/\pi (\mathrm{\Gamma })`$.
(iii) The kernels of each of the actions $``$, $``$ are of finite indexes since the corresponding quotients are isomorphic to subgroups in Permut(X). So, the intersection of kernels, i.e. subgroup $`\mathrm{\Gamma }`$ has finite index as well. $`\mathrm{}`$
###### Definition 6
We call a 7-tuple
$$(G,A,X,\rho _{GA},\rho _{GAX},\overline{\pi },\overline{\psi _A})$$
a bijective cocycle 7-tuple if $`G,A`$ are groups, $`\rho _{GA}:GAut(A)`$ is an action of $`G`$ on $`A`$, $`\rho _{GAX}:GAPermut(X)`$ is an action of $`GA`$ on $`X`$, $`\overline{\pi }:GA`$ is a bijective 1-cocycle, $`\overline{\psi _A}:XA`$ is $`GA`$-equivariant, where $`GA`$ acts on $`A`$ by conjugation.
Note that the action of $`GA`$ gives rise to two actions $`\rho _{GX}:GPermut(X)`$ and $`\rho _{AX}:APermut(X)`$. Starting with a bijective cocycle 7-tuple, let $`:X\times XX`$ and $`\varphi :X\times XX`$ be defined as $`yx=\rho _{GX}(\overline{\psi _G}(y))(x)`$ and $`\varphi (x,y)=\rho _{AX}(\overline{\psi _A}(y))(x)`$, where $`\overline{\psi _G}=\overline{\pi }^1\overline{\psi _A}`$. Define $`S:X\times XX\times X`$ and $`S^{}:X\times XX\times X`$ via $`S(x,y)=(xy,y^1x)`$ and $`S^{}(x,y)=(\varphi (y,x),x)`$, where $`xy`$ is defined in such a way that $`\varphi (y,x)=x^1((yx)y)`$ holds.
###### Lemma 3
(i) For an arbitrary bijective cocycle 7-tuple, $`(X,S)`$ constructed above is a braided nondegenerate set.
(ii) $`(X,S^{})`$ is a derived nondegenerate braided set corresponding to $`(X,S)`$.
Proof:
It is obvious from the definition of a bijective cocycle 7-tuple that $`(X,S^{})`$ is a braided nondegenerate set. We introduced $`xy`$ in such a way that $`(X,S^{})`$ becomes a derived solution corresponding to $`(X,S)`$ as long as we are able to show that $`(X,S)`$ is braided nondegenerate itself. Since $`\varphi `$ is $`G`$-invariant the argument identical to the one made in the proof of Theorem 2.3 mitigates that for each integer $`n2`$ there exists a bijection $`J_n:X^nX^n`$ such that $`J_nS_n^{ii+1}J_n^1=(S^{})_n^{ii+1}`$. This implies that $`S`$ is bijective and satisfies the braid relation (2.1).
It remains to prove that $`(X,S)`$ is nondegenerate, i.e. that $`xy`$ depends bijectively on $`yX`$ for a fixed $`xX`$. For that purpose we show that we can define $`gyX`$ for $`yX,gG`$ such that for any $`x,yX`$, $`g,hG`$ $`xy=\overline{\psi _G}(x)y`$ and $`(gh)y=g(hy)`$, i.e. $``$ is an action of $`G`$ on $`X`$.
Let a group homorphism $`P:GGA`$ be defined by the formula $`P(g)=(g,\overline{\pi }(g))`$ for $`gG`$. Then the map $`\rho _{GX}^{}=\rho _{GAX}P`$ is an action of $`G`$ on $`X`$. For $`gG,yX`$ define
$$gy=\rho _{GX}^{}(\overline{\pi }^1(\rho _{GA}(\overline{\psi _G}(y)^1)(\overline{\pi }(g))))(y).$$
We want to check $`(gh)y=g(hy)`$. For notational convience let $`\rho _{GA}(\overline{\psi _G}(y)^1)(\overline{\pi }(g))=\overline{\psi _G}(y)^1\overline{\pi }(g)`$, i.e. by $``$ we will mean both actions $`\rho _{GA}`$ and $`\rho _{GX}`$. Then due to the fact that $`P`$ is a homomorphism it is enough to verify that
$$\overline{\pi }^1(\overline{\psi _G}(hy)^1\overline{\pi }(g))=\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(gh))(\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h))^1.$$
Since $`\overline{\pi }(gh)=(h^1\overline{\pi }(g))\overline{\pi }(h)`$,
$`\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(gh))=\overline{\pi }^1((\overline{\psi _G}(y)^1h^1\overline{\pi }(g))(\overline{\psi _G}(y)^1\overline{\pi }(h)))`$
$`=\overline{\pi }^1(\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h))\overline{\psi _G}(y)^1h^1\overline{\pi }(g))\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h)).`$
In this way, if we check that
$$\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h))\overline{\psi _G}(y)^1h^1=\overline{\psi _G}(hy)^1$$
(2.14)
we conclude that $`g(hy)=(gh)y`$. Let us rewrite (2.14) as (cf. Lemma 2) $`\overline{\psi _G}(hy)\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h))=h\overline{\psi _G}(y)`$ and apply $`\overline{\pi }`$ to it. We get that (2.14) is equivalent to
$$\overline{\psi _A}(\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h))^1(hy))(\overline{\psi _G}(y)^1\overline{\pi }(h))=(\overline{\psi _G}(y)^1\overline{\pi }(h))(\overline{\psi _A}(y)).$$
The last equality follows from $`A`$-equivariance of $`\overline{\psi _A}`$ since $`\overline{\pi }^1(\overline{\psi _G}(y)^1\overline{\pi }(h))^1(hy)=\rho _{AX}(\overline{\psi _G}(y)^1\overline{\pi }(h))(y)`$ by definition of $`hy`$. Lemma is proved.
$`\mathrm{}`$
One can combine two actions of $`G`$ on $`X`$ \- $`\rho _{GX}`$ and $`\rho _{GX}^{}`$ into the action $`\rho =\rho _{GX}\times \rho _{GX}^{}:GPermut(X)\times Permut(X)`$ of $`G`$ on $`X^2`$.
###### Definition 7
We call a bijective cocycle 7-tuple $`(G,A,X,\rho _{GA},\rho _{GAX},\overline{\pi },\overline{\psi _A})`$ faithful if $`\overline{\psi _G}(X)=\overline{\pi }^1\overline{\psi _A}(X)`$ generates $`G`$ and the action $`\rho :GPermut(X)\times Permut(X)`$ is faithful.
The following theorem is a characterization of finite braided nondegenerate sets in group theoretical terms. We notice that this result is a generalization of Proposition 2.11 in \[ESS\].
###### Theorem 2.7
The construction of Lemma 3 establishes a 1-1 correspondence between nondegenerate braided sets $`(X,S)`$ and faithful bijective cocycle 7-tuples.
Proof:
Having a nondegenerate braided set $`(X,S)`$ it is straightforward to construct the faithful bijective cocycle 7-tuple. Indeed, following the notations of Theorem 2.6 we put $`G=G_X/\mathrm{\Gamma }`$, $`A=A_X/\pi (\mathrm{\Gamma })`$ then $`G`$ acts faithfully on $`X^2`$ and we have a faithful cocycle 7-tuple. Conversely, if we have a faithful 7-tuple we can construct a nondegenerate braided set by Lemma 3. Let $`G_X`$ be the structure group of so constructed nondegenerate braided set. It follows from (2.14) that there is a group homomorphism $`Q:G_XG`$ such that $`Q\psi _G=\overline{\psi _G}`$. Since $`\overline{\psi _G}(X)`$ generates $`G`$, $`Q`$ is surjective, and since $`G`$ acts faithfully on $`X^2`$, $`Ker(Q)=\mathrm{\Gamma }`$, i.e. $`G_X/\mathrm{\Gamma }`$ is isomorphic to $`G`$. $`\mathrm{}`$
### 2.3 Injective solutions
In this section we talk about most tractable braided nondegenerate sets - injective solutions.
###### Definition 8
We call a braided nondegenerate set $`(X,S)`$ an injective solution if the map $`\psi _G:XG_X`$ is injective.
In \[LYZ\] the authors noticed that in the absence of involutivity the natural map $`\psi _G:XG_X`$ is not obviously injective which creates difficulties in characterization of solutions. It turned out that injectivity may indeed fail (see examples below). This motivates Definition 8.
###### Lemma 4
A nondegenerate braided set $`(X,S)`$ is injective if and only if its derived solution $`(X,S^{})`$ is injective.
Proof:
The statement of Theorem 2.5 implies that injectivity of the map $`\psi _A:XA_X`$ is equivalent to injectivity of $`\psi _G:XG_X`$. Since $`A_X`$ is the structure group of the derived solution Lemma 4 is proved. $`\mathrm{}`$
The importance of injective solutions is in the fact that their properties and group-theoretical characterization are very similar to that of involutive solutions \[ESS\].
###### Theorem 2.8
(i) Let a group $`G`$ act on a group $`A`$ by $`\rho _{GA}:GAut(A)`$ such that the bijective map $`\pi :GA`$ is a 1-cocycle. Then any $`GA`$-invariant subset $`XA`$ has a natural structure of a nondegenerate braided injective set given by
$$S(x,y)=(\pi (\pi ^1(x)\pi ^1(y)(\pi ^1(\rho (\pi ^1(y)^1)(x)))^1),\rho (\pi ^1(y)^1)(x)),$$
(2.15)
for $`x,yX`$.
(ii) Any nondegenerate braided injective set can be obtained by the method just described.
Proof:
To prove (i) we can use Lemma 3. Indeed, we have the following bijective cocycle 7-tuple: $`G`$, $`A`$, $`X`$, $`\overline{\pi }=\pi `$, $`\rho _{GA}`$ as given; $`\rho _{GAX}`$ is induced from the adjoint action of $`GA`$ on its subgroup $`A`$, $`\overline{\psi _A}=id_X`$. It is straightforward to check (cf. (2.14)) that the map $`S`$ constructed in Lemma 3 coincides with the map $`S`$ given by formula (2.15). So it remains to prove that the set $`(X,S)`$ is injective. Let $`G_X`$ be its structure group. Arguing as in the proof of Theorem 2.7 we see that there is a group homomorphism $`Q:G_XG`$ such that $`\pi Q\psi _G=id_X`$. This implies that $`\psi _G`$ is injective.
Conversely, if $`(X,S)`$ is an injective nondegenerate braided set then $`X=\psi _A(X)`$ is a $`G_X`$-invariant (w.r.t. $``$ \- action) subset in $`A_X`$. Note that this subset is automatically $`G_XA_X`$-invariant. Recall that $`S`$ is given by the formula $`S(x,y)=(xy,y^1x)`$ for $`x,yX`$ according to notations after Theorem 2.2. Let us make sure that the construction of Theorem 2.8 yields the same map $`S`$ we already have. Indeed, for $`x,y\psi _A(X)=X`$, $`\pi ^1(y)^1x=\psi _G(y)^1x=y^1x`$ and $`\pi (\pi ^1(x)\pi ^1(y)(\pi ^1(\pi ^1(y)^1x))^1)=\pi (\psi _G(x)\psi _G(y)\psi _G(y^1x)^1)=\pi (\psi _G(xy))=\psi _A(xy)=xy`$. Theorem is proved. $`\mathrm{}`$
Lemma 4 implies that injectivity of a given solution is determined by the properties of the function $`\varphi :X\times XX`$. In particular, for a symmetric set $`(X,S)`$, $`\varphi (y,x)=y`$ and $`A_X`$ is the free abelian group generated by $`X`$. Hence, symmetric sets are injective. We don’t know any easy way to check that a given function $`\varphi (y,x)`$ corresponds to an injective solution. While an injectivity criterion is provided by Theorem 2.9, we give two simple necessary conditions below that are in many cases sufficient to check that a given solution is not injective.
###### Lemma 5
If $`(X,S)`$ is an injective braided set then
(a) $`\varphi (x,x)=x`$ for all $`xX`$,
(b) a pair $`(y,x)X\times X`$ satisfies $`\varphi (y,x)=y`$ if and only if $`\varphi (x,y)=x`$.
Proof:
Suppose $`(X,S)`$ is injective. The group $`A_X`$ is generated by the elements of $`X`$ subject to relations $`\varphi (y,x)y=yx`$ for $`x,yX`$. Consequently $`\psi _A(\varphi (x,x))\psi _A(x)=\psi _A(x)\psi _A(x)`$ and $`\psi _A(\varphi (x,x))=\psi _A(x)`$ in $`A_X`$. Since $`\psi _A:XA_X`$ is injective $`\varphi (x,x)=x`$ on $`X`$. Now, assume that $`\varphi (y,x)=y`$. Then $`\psi _A(y)\psi _A(x)=\psi _A(x)\psi _A(y)`$. On the other hand, $`\psi _A(\varphi (x,y))\psi _A(y)=\psi _A(y)\psi _A(x)`$, therefore $`\psi _A(\varphi (x,y))=\psi _A(x)`$ and $`\varphi (x,y)=x`$. Lemma is proved. $`\mathrm{}`$
Example. Let $`c,bPermut(X)`$. Define $`S:X\times XX\times X`$ by the formula $`S(x,y)=(by,cx)`$. It is easy to see that $`(X,S)`$ is a nondegenerate braided set if $`bc=cb`$. We claim that this solution is injective if and only if $`cb=id_X`$. Indeed, suppose the solution is injective, then by Lemma 5 $`\varphi (x,x)=x`$, i.e. $`x^1((xx)x)=cbx=x`$. Therefore $`cb=id_X`$. Conversely, if $`cb=id_X`$ then $`(X,S)`$ is symmetric and hence injective.
We remark that with each nondegenerate braided set we associated two actions of the group $`G_X`$ ($`,`$) and an action of $`A_X`$ (via $`\varphi (x,y)`$, which is a $``$ action of the derived solution) on $`X`$. In particular, the latter action allows us to construct a finite group $`A_X^0Permut(X)`$, as the image of $`A_X`$ under that action, and a surjective homomorphism $`p:A_XA_X^0`$. Define $`M_X`$ \- a module over $`A_X^0`$ generated by $`v_x`$, $`xX`$ subject to relations
$$p(y)^1v_x+v_y=p(x)^1v_{\varphi (y,x)}+v_x.$$
By construction we have a natural map $`\psi _M:XM_X`$ given by $`xv_x`$. It turns out that injectivity of $`\psi _M`$ is equivalent to injectivity of $`\psi _G`$.
###### Theorem 2.9
(i) There exists a unique 1-cocycle $`\theta :A_XM_X`$ such that $`\theta \psi _A=\psi _M`$, where $`A_X`$ acts on $`M_X`$ via $`p`$.
(ii) $`\theta `$ is injective on $`\psi _A(X)`$.
Proof:
Statement (i) is clear from definitions of $`A_X`$ and $`M_X`$. Let us show that (ii) holds. Let $`Ker(p)=\mathrm{\Gamma }_A`$, i.e. $`A_X^0=A_X/\mathrm{\Gamma }_A`$. Let $`x_1,x_2\psi _A(X)`$, $`x_1x_2`$. We want to show that $`\theta (x_1)\theta (x_2)`$ in $`M_X`$. Fix a character $`\xi :\mathrm{\Gamma }_A^{}`$. Define a vector space
$$V_\xi =\{f:A_X|f(a\gamma )=f(a)\xi (\gamma ),\gamma \mathrm{\Gamma }_A\}.$$
$`V_\xi `$ clearly has an $`A_X`$-module structure defined as $`(bf)(a)=f(b^1a)`$, where $`fV_\xi ,b,aA_X`$. Choose a lifting $`g:A_X^0A_X`$ (as a set only). Then $`V_\xi `$ is identified with $`Fun(A_X^0,)`$, the space of functions on $`A_X^0`$ via $`ff|_{g(A_X^0)}`$. This space has a basis $`\delta _a`$, $`aA_X^0`$ such that $`\delta _a(b)=0`$ if $`ab`$ and $`\delta _a(a)=1`$. Let us define $`\epsilon :\psi _A(X)Fun(A_X^0,)`$ by the formula $`x\delta _a=\delta _{p(x)a}\epsilon (x)(a)`$ for $`x\psi _A(X)`$, $`aA_X^0`$. We can always choose $`\xi `$ and $`g`$ such that $`\epsilon (x_1)\epsilon (x_2)`$. Indeed, if $`p(x_1)=p(x_2)`$ it suffices to choose $`\xi `$ such that $`\xi (x_1x_2^1)1`$ for any lifting $`g`$. If $`p(x_1)p(x_2)`$ then for any character $`\xi 1`$ there is a lifting $`g`$ such that $`\epsilon (x_1)\epsilon (x_2)`$. Moreover, since $`xy=\varphi (y,x)x`$ for $`x,y\psi _A(X)`$ and
$$xy\delta _a=x\delta _{p(y)a}\epsilon (y)(a)=\delta _{p(x)p(y)a}\epsilon (x)(p(y)a)\epsilon (y)(a),$$
we get that
$$\epsilon (x)(p(y)a)\epsilon (y)(a)=\epsilon (\varphi (y,x))(p(x)a)\epsilon (x)(a).$$
In this way, we have the equality $`(p(y)^1\epsilon (x))\epsilon (y)=(p(x)^1\epsilon (\varphi (y,x)))\epsilon (x)`$ in $`Fun(A_X^0,)`$. Hence we can construct an $`A_X^0`$-homomorphism from $`M_X`$ to $`Fun(A_X^0,)`$ given by $`v_z\epsilon (\psi _A(z))`$, $`zX`$. But $`\epsilon (x_1)\epsilon (x_2)`$ thus $`\theta (x_1)\theta (x_2)`$.
Corollary 1 The map $`\psi _M:XM_X`$ is injective if and only if $`(X,S)`$ is an injective solution.
### 2.4 Rank of the structure group
In this section we show how to compute the rank of the structure group $`G_X`$ for a finite nondegenerate braided set $`(X,S)`$.
###### Definition 9
(i) The rank of a group G having an abelian subgroup of finite index $`\mathrm{\Gamma }`$ is defined as the rank of $`\mathrm{\Gamma }`$.
(ii) We define the rank of a finite nondegenerate braided set $`(X,S)`$ to be the rank of its structure group $`G_X`$.
Clearly the above definition doesn’t depend on the choice of $`\mathrm{\Gamma }`$, for any two abelian subgroups of finite index $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ have the same rank that is equal to the rank of their intersection, which has finite index in each of them.
###### Lemma 6
The rank of a solution $`(X,S)`$ is equal to the rank of the derived solution $`(X,S^{})`$.
Proof:
According to statement (i) of Theorem 2.6 abelian subgroups $`\mathrm{\Gamma }`$ and $`\pi (\mathrm{\Gamma })`$ of finite indexes in $`G_X`$ and $`A_X`$ are isomorphic, thus $`G_X`$ and $`A_X`$ have the same rank. Lemma is proved.
We aim to compute the rank of $`A_X`$. Note that defining relations in group $`A_X`$ can be rewritten as
$$\varphi (y,x)=xyx^1.$$
(2.16)
Note that since $`S(x,y)=(\varphi (y,x),x)`$ gives rise to a nondegenerate braided set, $`\varphi (,x)`$ can be extended to the action of $`A_X`$ on $`X`$. Introduce an equivalence relation $``$ on $`X`$ such that the orbits of the above action become equivalence classes. Namely, consider the minimum equivalence relation $``$ such that $`y\varphi (y,x)`$ for $`x,yX`$.
###### Theorem 2.10
The rank of group $`A_X`$ is equal to the number of equivalence classes with respect to equivalence relation $``$ on $`X`$.
The proof relies on the following lemma. Let $`H_1`$, $`H_2`$, $`H_3`$ be three groups, such that there is an exact sequence
$$1H_1H_2H_31$$
and $`H_1`$ is cenral in $`H_2`$.
###### Lemma 7 (\[CR\])
There exists an exact sequence (The Hochschild - Serre sequence) (2.17) for any abelian group B, where $`H^2(H_3,B)`$ stands for the second cohomology group of $`H_3`$ with coefficients in B.
$$1Hom(H_3,B)Hom(H_2,B)Hom(H_1,B)H^2(H_3,B)$$
(2.17)
Proof of Theorem 2.10.
Let us use the sequence 2.17 in the following situation: $`H_2=A_X`$, $`H_1\pi (\mathrm{\Gamma })A_X`$ is the free abelian group of rank $`r`$ of finite index in $`\pi (\mathrm{\Gamma })`$, $`H_3=H_2/H_1`$ \- a finite group, $`B=^{}`$. Notice that $`\pi (\mathrm{\Gamma })`$ is central in $`A_X`$ therefore $`H_1`$ is central in $`H_2`$. Then, since both $`Hom(H_3,B)`$ and $`H^2(H_3,B)`$ are finite groups and $`Hom(H_1,B)=(^{})^r`$, the dimension of $`Hom(H_2,B)=Hom(A_X,^{})`$ is equal to $`r`$, the rank of $`A_X`$. On the other hand, it is clear from formula (2.16) that $`Hom(A_X,^{})=(^{})^k`$, where $`k`$ is the number of equivalence classes in $`X`$. Theorem is proved.
Corollary 1 The rank of any solution $`(X,S)`$ is less or equal than $`n`$, the number of elements in $`X`$, with the equality taking place if and only if $`(X,S)`$ is symmetric.
Proof:
It is clear that if rank of $`A_X`$ is equal to $`n`$ then equivalence relation $``$ on $`X`$ is trivial, $`\varphi (y,x)=y`$ and the map $`\psi _A:XA_X`$ is injective. This immediately implies that the pair $`(X,S)`$ is symmetric.
Example. Consider a permutation solution $`S(x,y)=(bx,cy)`$, $`b,cPermut(X)`$, $`bc=cb`$. It is easy to check that $`\varphi (y,x)=bcy`$, hence $`y(bc)^my`$ for any integer m. The rank of the permutation solution is equal to the number of equivalence classes with respect to this relation, which, in turn, is equal to the number of independent cyclic permutations in canonical decomposition of $`bc`$ (counting cyclic permutations of length 1).
## 3 Linear And Affine Solutions
### 3.1 Linear Braided Sets
In this section we will look for nondegenerate braided sets of the following form: $`X`$ is an abelian group, and $`S`$ is an affine linear transformation of $`X\times X`$. Such braided sets will be called affine solutions. Considering affine solutions was motivated by the results in \[ESS\].
We will start with considering a special case, when $`S`$ is an automorphism of $`X\times X`$. In this case, an affine solution will be called a linear solution. For a linear solution, $`S`$ has the form
$$S(x,y)=(ax+by,cx+dy),a,b,c,dEnd(X).$$
(3.1)
It is easy to check that for $`S`$ of the form (3.1) the nondegeneracy is equivalent to invertibility of both $`b`$ and $`c`$ while braid relation is equivalent to the equations \[Hi\]
$`a(1a)=bac,`$ $`d(1d)=cdb,`$ (3.2)
$`ab=ba(1d),`$ $`ca=(1d)ac,`$ $`dc=cd(1a),`$ (3.3)
$`bd=(1a)db,`$ $`cbbc=adadad.`$ (3.4)
###### Lemma 8
Braided nondegenerate linear sets $`(X,S)`$ are in 1-1 correspondence with the quadruples $`(a,b,d,s)End(X)^4`$ such that:
(i) $`1a,1d,b,1+s`$ are invertible,
(ii) s commutes with a,b,d and $`sa=sd=0`$,
(iii) $`bdb^1=(1a)d,b^1ab=a(1d).`$ The 1-1 correspondence is given via the formula
$$bc=(1d+ad)(1a)+s.$$
(3.5)
Proof:
Suppose $`(a,b,c,d)`$ solves (3.2)-(3.4). Note that the first of equations (3.4) implies that $`bdb^1=(1a)d`$, therefore $`b(1d)b^1=1d+ad`$. Moreover, if we multiply the first equation of (3.3) by $`b^1`$ on the right we get that $`a=bab^1b(1d)b^1`$ thus relations (3.6) hold. Similarly from equations two and three of (3.3) we obtain relations (3.7).
$`bab^1(1d+ad)=a,bdb^1=(1a)d,`$ (3.6)
$`cac^1=(1d)a,cdc^1(1a+da)=d.`$ (3.7)
Above formulas show how to conjugate the elements of subalgebra generated by $`a`$ and $`d`$ by elements $`b`$, $`c`$ and their products. We define $`s`$ from the relation (3.5). Notice that $`sa=sd=0`$. Indeed, multiplying (3.5) by $`a`$ on the right and using first of relations (3.2) and second of relations (3.3) we get that
$$sa=bca(1d+ad)(1a)a=b(1d)ac(1d+ad)bac.$$
Since according to (3.6) $`b(1d)b^1=1d+ad`$ we see that $`b(1d)ac(1d+ad)bac=0`$. Similarly $`as=0`$. The last of equations (3.4) imply that
$$cb=(1a+da)(1d)+s.$$
(3.8)
Multiplying the relation (3.8) just obtained by $`d`$ we get that $`sd=ds=0`$. Now we have everything to show that $`1a`$, $`1+s`$ and $`1d`$ are invertible. Indeed, $`b,c`$ are invertible because of nondegeneracy of $`(X,S)`$. Since
$$bc=(1d+s+ad)(1a)=(1a)(1d+s+da),$$
$$bc=(1d+ad)(1a+s)=(1a+s)(1d+ad),$$
and
$$cb=(1a+s+da)(1d)=(1d)(1a+s+ad)$$
we conclude that $`1a`$ , $`1a+s`$ and $`1d`$ have right and left inverses and hence invertible. But $`1a+s=(1a)(1+s)`$, thus $`1+s`$ is invertible. Let us show that $`s`$ commutes with $`c`$. We conjugate relation (3.5) by $`c`$ and use relations (3.7) to conclude that
$$cb=(1a+da)(1d)+csc^1.$$
(3.9)
Comparing (3.8) to (3.9) we get that $`s=csc^1`$, i.e. $`s`$ commutes with $`c`$. Now, relation (3.5) implies that $`s`$ commutes with $`b`$ as well. In this way, starting from $`(a,b,c,d)`$, a solution to (3.2)-(3.4) we constructed a quadruple $`(a,b,d,s)`$ satisfying conditions of the lemma. Conversely, let us assume we start with $`(a,b,d,s)`$. Define $`c`$ from the formula (3.5). Since $`1d+ad=b(1d)b^1`$ is invertible, $`c`$ is invertible as well. It is straightforward to check that so defined $`(a,b,c,d)`$ satisfy the braid relations (3.2)-(3.4), and that $`S:X\times XX\times X`$ is bijective. The lemma is proved.
### 3.2 Injective Linear Solutions
It turns out that injective linear solutions are easy to characterize.
###### Theorem 3.1
A linear nondegenerate braided set of the form (3.1) is injective iff $`bc=(1d+ad)(1a)`$ or, in the language of Lemma 8, $`s=0`$.
Proof:
Assume that $`(X,S)`$ is injective. Then, by Lemma 5 $`\varphi (x,x)=x`$ on $`X`$. In the linear case it is easy to compute $`\varphi (y,z)`$ explicitly:
$$\varphi (y,z)=c((yz)y)+dz=c(a(yz)+by)+dz=cac^1(zdy)+cby+dz.$$
According to relation (3.8), $`cb=(1a+da)(1d)+s`$. Plugging it into formula for $`\varphi (y,z)`$ and using (3.7) we conclude that
$$\varphi (y,z)=(1(1d)(1a))z+((1d)(1a)+s)y.$$
Condition $`\varphi (x,x)=x`$ immediately implies that $`s=0`$.
Conversely, assume that $`s=bc(1d+ad)(1a)=0`$. Let us show that $`\psi _A:XA_X`$ is injective. The group $`A_X`$ is generated by elements of $`X`$ subject to relations
$$\varphi (y,z)z=zy,wherey,zXand\varphi (y,z)=z^1((yz)y).$$
(3.10)
Denoting $`K=(1d)(1a)`$ we see that the group $`A_X`$ is given by relations $`((1K)z+Ky)z=zy`$. By Lemma 8, $`K`$ is invertible. Let us now define the action of $`A_X`$ on $`X`$. We let the elements of generating set $`\psi _A(X)`$ act with $`K`$ on $`X`$ and then extend this action to arbitrary elements of $`A_X`$. Consider a semidirect product $`A_XX`$ with respect to this action, and define an embedding $`J:XA_XX`$ given by the formula $`J(x)=(\psi _A(x),x)`$. We notice that $`J((1K)z+Ky)J(z)=J(z)J(y)`$ in $`A_XX`$. Indeed,
$`J((1K)z+Ky)J(z)=(\psi _A((1K)z+Ky),(1K)z+Ky)(\psi _A(z),z)`$
$`=(\psi _A((1K)z+Ky)\psi _A(z),K^1((1K)z+Ky)+z)`$
$`=(\psi _A(z)\psi _A(y),K^1z+y)=J(z)J(y).`$
So, $`J`$ can be extended to a homomorphism $`\widehat{J}:A_XA_XX`$ such that $`Pr\widehat{J}\psi _A=id`$, where $`Pr:A_XXX`$ stands for the projection to the second component. Therefore, the map $`\psi _A:XA_X`$ is injective. Theorem is proved.
Corollary 1 (i) Let $`(X,S)`$ be a nondegenerate braided linear set of the form (3.1). Then the pair $`(X,\widehat{S})`$ with $`\widehat{S}:X\times XX\times X`$ given by $`\widehat{S}(x,y)=(ax+by,cx+(ds)y)`$, where $`s`$ is defined in 3.5, is an injective solution.
(ii) Suppose $`(X,S)`$ is an injective linear solution and $`s:XX`$ satisfies $`sa=as=0`$, $`sb=bs`$, $`sd=ds=s^2`$, $`sc=cs`$. Then, $`(X,\stackrel{˘}{S})`$ with $`\stackrel{˘}{S}:X\times XX\times X`$ given by $`\stackrel{˘}{S}(x,y)=(ax+by,cx+(d+s)y)`$ is a nondegenerate braided set that corresponds under the correspondence of Lemma 8 to the quadruple $`(a,b,d,s)`$.
Proof:
(i) It is easy to check $`(X,\widehat{S})`$ is a braided nondegenerate set by directly checking relations (3.2) - (3.4). Since $`s`$ commutes with everything and $`sa=sd=0`$ the above task is pretty simple. Also, it is obvious that $`bc=(1(ds)+a(ds))(1a)`$, thus by Theorem 3.1 $`(X,\widehat{S})`$ is injective. Proof of (ii) is similar to (i) and is left to the reader.
Examples.
1. Consider the linear solution $`S(x,y)=(cy,bx)`$ with $`c,b`$ being linear automorphisms of $`X`$ subject to $`cb=bc`$. We see that s=bc-1, therefore $`(X,\widehat{S})`$ with $`\widehat{S}(x,y)=(cy,bx+(1bc)y)`$ is an injective solution. We obtain the same solution by a different method in Example 2 at the end of this section.
2. It was shown in \[ESS\] that symmetric nondegenerate linear solutions of the form 3.1 are given as the solutions to the following equations: $`bab^1=\frac{a}{a+1}`$, $`c=b^1(1a^2)`$, $`d=\frac{a}{a1}`$ with $`b,c`$ being invertible. In particular a large class of solutions of this kind considered corresponded to nilpotent $`a`$, i.e. there was $`n`$ such that $`a^n=0`$. Define $`s=ma^{n1}`$, $`m`$ being any integer. It is easy to see that so defined $`s`$ satisfies conditions of Corollary 1 (ii) hence $`S^{\prime \prime }(x,y)=(ax+by,cx+(d+ma^{n1})y)`$ is a nondegenerate braided set, which is not symmetric and not injective unless $`ma^{n1}=0`$.
###### Theorem 3.2
Linear braided nondegenerate injective sets $`(X,S)`$ on an abelian group $`X`$ are in 1-1 correspondence with triples $`(a,b,d)`$ of endomorphisms of $`X`$ such that $`b,1a,1d`$ are invertible and $`bdb^1=(1a)d,b^1ab=a(1d)`$.
Proof: Straightforward application of Lemma 8 and Theorem 3.1.
Corollary 1 Linear braided nondegenerate injective sets $`(X,S)`$ on an abelian group $`X`$ are in 1-1 correspondence with triples $`(p,q,z)`$ of automorphisms of $`X`$ such that $`pq=qp`$ and $`z^2z(p+q)+pq=0`$.
The statement of Corollary 1 follows from Theorem 3.2 via change of variables $`p=b^1,q=(1a)(1d)b^1,z=(1a)b^1`$.
Examples.
1. If $`(X,S)`$ is unitary then $`(1a)(1d)=1`$ and therefore $`p=q`$. In this way, nondegenerate unitary linear braided sets are characterized as representations of algebra generated by invertible $`p,z`$ subject to $`z^22zp+p^2=0`$.
2. Put $`z=p`$, then $`q`$ can be anything as long as it is invertible and $`pq=qp`$. Correspondingly, $`a=0,b=p^1,1d=p^1q,c=q`$. The map $`S`$ is defined in the following way: $`S(x,y)=(p^1y,qx+(1p^1q)y)`$. Similarly if we let $`z=q`$ we obtain the solution given by the formula: $`S(x,y)=(py+(1q^1p)x,q^1x)`$.
3. Let $`ϵ_1`$ and $`ϵ_2`$ be two nilpotent operators on $`X`$ of order two, i.e. $`ϵ_1^2=ϵ_2^2=0`$. Then we let $`p=1+ϵ_1`$, $`q=1ϵ_1`$, $`z=1+ϵ_2`$. It is easy to check that so defined $`p`$ and $`q`$ commute and that $`z^2z(p+q)+pq=0`$. We recover $`a`$, $`b`$, $`c`$, $`d`$ from $`ϵ_1`$ and $`ϵ_2`$ via the formulas $`b=1ϵ_1`$, $`a=ϵ_1ϵ_2+ϵ_2ϵ_1`$, $`d=1(1+ϵ_1)(1ϵ_2)(12ϵ_1)`$, $`c=(1+ϵ_1)(1ϵ_2)(1ϵ_1)(1+ϵ_2)(1ϵ_1)`$. In particular, if $`ϵ_1ϵ_2=ϵ_2ϵ_1`$ the corresponding linear solution has the form:
$$S(x,y)=((ϵ_1ϵ_2+ϵ_2ϵ_1)x+(1ϵ_1)y,(1ϵ_1)x+(ϵ_1+ϵ_2ϵ_2ϵ_1)y)$$
4.Suppose $`X`$ is a n-dimensional vector space and $`p,q`$ are invertible operators on it having $`2n`$ distinct eigenvalues. Let $`(v_1\mathrm{}v_n)`$ be the basis of $`X`$ in which both $`p`$ and $`q`$ are diagonalized, i.e. $`pv_i=p_iv_i`$, $`qv_i=q_iv_i`$. Then $`zv_i=z_iv_i`$, where for each $`i`$ either $`z_i=p_i`$ or $`z_i=q_i`$. Indeed, $`0=(z^2z(p+q)+pq)v_i=(zp_i)(zq_i)v_i`$. The vector subspace of $`X`$ generated by applying $`z`$ to $`v_i`$ ($`i`$ is fixed) is annihilated by $`(zp_i)(zq_i)`$ therefore it has a basis of $`p_i`$ and $`q_i`$ eigenvectors of $`z`$. Therefore, it has to be one dimensional, i.e. $`zv_i=p_iv_i`$ or $`zv_i=q_iv_i`$. It is also easy to see that $`z`$ given by diagonal matrix in $`(v_1\mathrm{}v_n)`$ basis with $`p_i`$ or $`q_i`$ entries on the $`i`$-th place does satisfy $`z^2z(p+q)+pq=0`$ on $`X`$.
### 3.3 Affine solutions
In this section we talk about general affine solutions on an abelian group $`X`$.
###### Definition 10
The solution (braided, nondegenerate) $`(X,S)`$ of the form (3.11) is called affine.
$$S(x,y)=(ax+by+z,cx+dy+t),a,b,c,dEnd(X),t,zX.$$
(3.11)
###### Lemma 9
The pair $`(X,S)`$ of the form (3.11) is a solution if and only if (3.2)-(3.4) and (3.12) hold. Therefore, any affine solution gives rise to a linear solution $`(X,S^{})`$ with $`S^{}(x,y)=(ax+by,cx+dy)`$.
$$cdz+dt=0,az+bat=0,(c+dad1)z+(da+1ab)t=0.$$
(3.12)
Proof: Straightforward.
###### Definition 11
We call $`(X,S^{})`$ from Lemma 9 the linear part of an affine solution $`(X,S)`$. Conversely, we call $`(X,S)`$ an affine solution associated with $`(X,S^{})`$.
###### Theorem 3.3
(i) Let $`(X,S^{})`$ be a linear braided nondegenerate set, $`S^{}(x,y)=(ax+by,cx+dy)`$. Then, $`(X,S)`$ given by (3.11) is an affine solution associated with $`(X,S^{})`$ if and only if $`t=c(1a)^1z+k`$ and $`ak=dk=0`$, $`(b1)k=sz`$, where s is defined in (3.5).
(ii) An affine solution $`(X,S)`$ is injective if and only if its linear part is injective and $`k=0`$ in the above characterization. In this way, injective affine solutions associated with a given injective linear solution $`(X,S^{})`$ are in 1-1 correspondence with elements $`zX`$, $`t`$ being given by $`t=c(1a)^1z`$.
Proof:
Let us proof part (i). Note that since $`cdz+dt=d(c(1a)^1z+t)`$ and az+bat=$`a(z+(1a)c^1t)`$ by equations (3.2)-(3.3) we can rewrite two of the relations (3.12) as
$$d(c(1a)^1z+t)=0,a(z+(1a)c^1t)=0.$$
Therefore, if we define $`k`$ from the relation
$$t=c(1a)^1z+k,$$
(3.13)
we see using (3.2) that $`ak=dk=0`$. Now, (3.7) implies $`c(1a)^1c^1=(1a+da)^1`$ thus we can transform (3.13) into $`(1a+da)t=cz+(1a+da)k=cz+k`$. In this way, we rewrite last of the relations (3.12) as
$$(dad1)zbt+k=0$$
(3.14)
If we substitute $`t`$ from (3.13) into (3.14) and use (3.5) we get that $`(b1)k=sz`$. Part (i) is proved. In order to prove part (ii) we compute the function $`\varphi :X\times XX`$. Define by $`\varphi ^{}`$ the corresponding function $`\varphi `$ of the linear part $`(X,S^{})`$. Then, it is easy to check that $`\varphi (y,x)=\varphi ^{}(y,x)+k`$. Now, assume that $`(X,S)`$ is injective. Then by Lemma 5 $`\varphi (x,x)=x`$, hence $`\varphi ^{}(x,x)=x`$ and $`k=0`$. As we saw in the proof of Theorem 3.1 $`\varphi ^{}(x,x)=x`$ implies that $`s=0`$ and thus $`(X,S^{})`$ is injective. Conversely, if $`(X,S^{})`$ is injective and $`k=0`$ then $`\varphi (y,x)=\varphi ^{}(y,x)`$ and hence the derived solution of $`(X,S)`$ coincides with the derived solution of $`(X,S^{})`$ and hence is injective. Theorem is proved. |
warning/0003/math0003229.html | ar5iv | text | # There are enough Azumaya algebras on surfaces
## Introduction
Generalizing the classical theory of central simple algebras over fields, Grothendieck introduced the Brauer group $`\mathrm{Br}(X)`$ and the cohomological Brauer group $`\mathrm{Br}^{}(X)`$ for schemes.
Let me recall the definitions. The *Brauer group* $`\mathrm{Br}(X)`$ comprises equivalence classes of Azumaya algebras. Two Azumaya algebras $`𝒜,`$ are called equivalent if there are everywhere nonzero vector bundles $`,`$ with $`𝒜\text{nd}()\text{nd}()`$. Let us define the *cohomological Brauer group* $`\mathrm{Br}^{}(X)`$ as the torsion part of the étale cohomology group $`H^2(X,𝔾_m)`$. Nonabelian cohomology gives an inclusion $`\mathrm{Br}(X)\mathrm{Br}^{}(X)`$, and Grothendieck asked whether this is bijective.
It would be nice to know this for the following reason: The cohomological Brauer group is related to various other cohomology groups via exact sequences, and this is useful for computations. In contrast, it is almost impossible to calculate the Brauer group of a scheme directly from the definition. Here is a list of schemes with $`\mathrm{Br}(X)=\mathrm{Br}^{}(X)`$:
1. Schemes of dimension $`1`$ and regular surfaces (Grothendieck ).
2. Abelian varieties (Hoobler ).
3. The union of two affine schemes with affine intersection (Gabber ).
4. Smooth toric varieties (DeMeyer and Ford ).
On the other hand, a nonseparated normal surface with $`\mathrm{Br}(X)\mathrm{Br}^{}(X)`$ recently appeared in . I wonder how the final answer to this puzzle will look like. The goal of this paper is to prove the following Theorem.
###### Theorem.
For separated geometrically normal algebraic surfaces, the inclusion $`\mathrm{Br}(X)\mathrm{Br}^{}(X)`$ is a bijection.
This adds some singular and nonprojective schemes to the preceding list. For quasiprojective surfaces, Hoobler ( Cor. 9) deduced the result directly from Gabber’s Theorem on affine schemes. Without ample line bundles, a different approach is required. Indeed, my initial motivation was to disprove the Theorem, rather than to prove it. The new idea is to use Maruyama’s theory of elementary transformations.
Here is an application of the preceding result:
###### Theorem.
Each proper normal algebraic surface admits a nonfree vector bundle.
It might easily happen that all line bundles are free . The existence of nonfree vector bundles can be viewed as a generalization, in dimension two, of Winkelmann’s Theorem , which asserts that each compact complex manifold has nonfree holomorphic vector bundles.
This paper has four sections. In the first section, I relate Azumaya algebras that are trivial on large open subsets to certain reflexive sheaves. In Section 2, we turn to normal surfaces and construct Azumaya algebras that are generically trivial by constructing the corresponding reflexive sheaves. This prepares the proof of the main Theorem, which appears in Section 3. The idea in the proof is to apply elementary transformations to Brauer–Severi schemes. The last section contains the existence result for nonfree vector bundles.
###### Acknowledgments.
This research was done in Bologna, and I am grateful to the Mathematical Department for its hospitality. I wish to thank Angelo Vistoli for suggestions, encouragement, and many stimulating discussions. Furthermore, I wish to thank Ofer Gabber, the referee, for his precise report. He found and corrected several mistakes. Several crucial steps are entirely due to Gabber, and the paper would be impossible without his contribution. The revision was done at M.I.T., and I wish to thank the Mathematical Department for its hospitality. Finally, I thank the DFG for financial support.
## 1. Azumaya algebras via reflexive sheaves
In this section, we shall describe Azumaya algebras that have trivial Brauer class on certain large open subsets. Throughout, $`X`$ will be a noetherian scheme. Let us call an open subset $`UX`$ *thick* if it contains all points $`xX`$ with $`0pt(𝒪_{X,x})1`$. In other words, $`0pt_{XU}(𝒪_X)2`$. A coherent $`𝒪_X`$-module $``$ is called *almost locally free* if it is locally free on some thick open subset $`UX`$, and has $`0pt_{XU}()2`$. Such sheaves behave well under suitable restriction and extension functors:
###### Lemma 1.1.
Let $`i:YX`$ be a thick open subset. Then the restriction map $`i^{}()`$ and the direct image map $`𝒢i_{}(𝒢)`$ induce inverse equivalences between the categories of almost locally free sheaves on $`X`$ and $`Y`$, respectively.
###### Proof.
This is similar to the proof of Theorem 1.12. Fix an almost locally free $`𝒪_X`$-module $``$. First, we check that $`\mathrm{\Gamma }(V,)\mathrm{\Gamma }(VY,)`$ is bijective for all affine open subsets $`VX`$. Setting $`A=VVY`$, we have an exact sequence of local cohomology groups
(1)
$$0H_A^0(V,)H^0(V,)H^0(VY,)H_A^1(V,)0.$$
Since $`0pt_A()2`$, the cohomology groups with supports vanish by Theorem 3.8. Therefore, the map in the middle is bijective. As a consequence, the adjunction map $`i_{}i^{}()`$ is bijective, so that the restriction functor $`i^{}()`$ is fully faithful.
Second, we check that the functor $`i^{}()`$ is essentially surjective. Fix an almost locally free $`𝒪_Y`$-module $`𝒢`$. By Corollary 6.9.8, the sheaf $`𝒢`$ extends to a coherent $`𝒪_X`$-module $``$. I claim that $`=^{}`$ is almost locally free. This is a local problem, so we may assume that there is a partial resolution $`_1_0^{}0`$ with coherent locally free sheaves, hence an exact sequence
(2)
$$0_0^{}_1^{}.$$
Let $`A=XU`$, where $`UY`$ is a thick open subset on which $`𝒢`$ is locally free. The exact sequence (2) gives an inclusion $`H_A^0(X,_0^{}/)H_A^0(X,_1^{})`$ and an exact sequence of local cohomology groups
$$H_A^0(X,_0^{}/)H_A^1(X,)H_A^1(X,_0^{}).$$
Since $`0pt_A(𝒪_X)2`$, the outer groups vanish, and we conclude $`0pt_A()2`$. Consequently, $``$ is almost locally free. By the same argument, we see that $`𝒢^{}`$ is almost locally free. Since the canonical map $`𝒢𝒢^{}`$ is bijective on some thick open subset, we conclude that it is bijective on $`Y`$, hence $`𝒢i^{}()`$.
It remains to check that $`𝒢i_{}(𝒢)`$ is the desired inverse equivalence. Extend $`𝒢`$ to an almost locally free $`𝒪_X`$-module $``$. Since the adjunction map $`i_{}i^{}()`$ is bijective, we are done. ∎
###### Remark 1.2.
Given two almost locally free sheaves $`_1`$ and $`_2`$, the sheaf $`=\text{om}(_1,_2)`$ is almost locally free as well. This is because, by Lemma 1.1, the middle map in (1) is bijective, so that the cohomology groups with support vanish. As a consequence, almost locally free sheaves are reflexive.
An Azumaya algebra $`𝒜`$ is called *almost trivial* if its Brauer class $`\mathrm{cl}(𝒜)\mathrm{Br}(X)`$ vanishes on some thick open subset. This easily implies that $`𝒜\text{nd}()`$ for some almost locally free sheaf $``$. However, the condition that the $`𝒪_X`$-algebra $`\text{nd}()`$ is an Azumaya algebra implies more.
###### Definition 1.3.
A coherent $`𝒪_X`$-module $``$ is called *balanced* if for each geometric point $`\overline{x}X`$, there is a decomposition $`_{\overline{x}}_{i=1}^r_{\overline{x}}`$ with $`r>0`$ for some almost invertible $`𝒪_{X,\overline{x}}`$-module $`_{\overline{x}}`$.
Here $`𝒪_{X,\overline{x}}`$ is the strict henselization of the local ring $`𝒪_{X,x}`$. Perhaps it goes without saying that *almost invertible sheaves* are invertible on a thick open subset and have depth $`2`$ outside. By fpqc-descent, balanced sheaves are almost locally free. They are closely related to Azumaya algebras, and the following result reduces the existence of certain Azumaya algebras to the existence of balanced sheaves.
###### Proposition 1.4.
Let $``$ be a balanced $`𝒪_X`$-module. Then $`\text{nd}()`$ is an almost trivial Azumaya algebra, and each almost trivial Azumaya algebra has this form.
###### Proof.
Obviously, $`𝒜=\text{nd}()`$ is a trivial Azumaya algebra on some thick open subset. To check that it is an Azumaya algebra on $`X`$, we may assume that $`X`$ is strictly local, and that $`=_{i=1}^r`$ for some almost invertible sheaf $``$. Being bijective on some thick open subset, the map $`𝒪_X\text{nd}()`$ is everywhere bijective by Lemma 1.1. Consequently, $`𝒜\text{at}_r(𝒪_X)`$ is an Azumaya algebra.
Conversely, let $`𝒜`$ be an almost trivial Azumaya algebra. Choose a thick open subset $`i:UX`$ on which the Brauer class is trivial. Then there is an isomorphism $`𝒜_U\text{nd}(𝒢)`$ for some locally free $`𝒪_U`$-module $`𝒢`$. By Lemma 1.1, this induces an isomorphism of algebras $`𝒜\text{nd}()`$, where $`=i_{}(𝒢)`$.
If remains to check that the almost locally free sheaf $``$ is balanced. To do so, we may assume that $`X`$ is strictly local, so $`𝒜=\text{at}_r(𝒪_X)`$. Now $``$ is a module over $`\text{at}_r(𝒪_X)`$. By Morita equivalence (see e.g. p. 53), $`=_{i=1}^r`$ for some coherent $`𝒪_X`$-module $``$. Clearly, $``$ is invertible on a thick open subset and has depth $`2`$ outside. In other words, $``$ is almost invertible. ∎
Next, we shall generalize some notions from Hartshorne’s paper on generalized divisors . For simplicity, we assume that $`X`$ satisfies Serre’s condition $`(S_2)`$, such that the points of codimension one are precisely the points of depth one. Set
$$𝒜𝒟\text{iv}_X=\underset{xX^{(1)}}{}(i_x)_{}(\mathrm{Div}(𝒪_{X,x}))$$
where the sum runs over all points of codimension one. The elements of the group $`\mathrm{ADiv}(X)=\mathrm{\Gamma }(X,𝒜𝒟\text{iv}_X)`$ are called *almost Cartier divisors*. For normal schemes, $`𝒜𝒟\text{iv}_X`$ is just the sheaf of Weil divisors. As in the normal case, an almost Cartier divisor $`D\mathrm{ADiv}(X)`$ defines an almost invertible sheaf $`𝒪_X(D)`$, which is invertible in codimension one and satisfies Serre’s condition $`(S_2)`$. According to Proposition 21.1.8, the canonical map $`𝒟\text{iv}_X𝒜𝒟\text{iv}_X`$ is injective. The exact sequence
$$0𝒟\text{iv}_X𝒜𝒟\text{iv}_X𝒫_X0$$
defines an abelian sheaf $`𝒫_X`$ on the étale site $`X_{\text{ét}}`$. For a geometric point $`\overline{x}X`$ with corresponding strict localization $`𝒪_{X,\overline{x}}=𝒪_{X,x}^{\mathrm{sh}}`$, the stalk is
$$𝒫_{X,\overline{x}}=\mathrm{ADiv}(𝒪_{X,\overline{x}})/\mathrm{Div}(𝒪_{X,\overline{x}}).$$
For normal schemes, this reduces to the class group $`\mathrm{Cl}(𝒪_{X,\overline{x}})`$. The preceding short exact sequence gives a long exact sequence in étale cohomology
$$\mathrm{ADiv}(X)H^0(X,𝒫)H^1(X,𝒟\text{iv}_X)H^1(X,𝒜𝒟\text{iv}_X).$$
We also have an exact sequence
$$0H^1(X,𝒟\text{iv}_X)H^2(X,𝔾_m)H^2(X^{(0)},𝔾_m),$$
where $`X^{(0)}=\mathrm{Spec}(𝒪_{X,\eta })`$ is the scheme of generic points, and you easily infer that the classes of almost trivial Azumaya algebra lie in the image of the iterated coboundary map
(3)
$$H^0(X,𝒫)H^1(X,𝒟\text{iv}_X)H^2(X,𝔾_m).$$
If $``$ is an almost locally free sheaf with decompositions $`_{\overline{x}}=_{i=1}^r_{\overline{x}}`$, the function
$$\overline{x}\mathrm{cl}(_{\overline{x}})\mathrm{ADiv}(𝒪_{X,\overline{x}})/\mathrm{Div}(𝒪_{X,\overline{x}})=𝒫_{X,\overline{x}}$$
depends only on $``$ and yields a section $`s_{}\mathrm{\Gamma }(X,𝒫_X)`$. The following result is due to Gabber:
###### Proposition 1.5 (Gabber).
Let $``$ be an balanced $`𝒪_X`$-module. Then the class of the Azumaya algebra $`\text{nd}()`$ in $`H^2(X,𝔾_m)`$ is the inverse of the image of the section $`s_{}H^0(X,𝒫_X)`$ under the iterated coboundary map in (3).
###### Proof.
Set $`𝒜=\text{nd}()`$. According to p. 341, its cohomology class in $`H^2(X,𝔾_m)`$ is given by the gerbe $`d(𝒜`$) of trivializations for $`𝒜`$, which associates to each étale $`UX`$ the groupoid of pairs $`(,\phi )`$, where $``$ is a locally free $`𝒪_U`$-module, and $`\phi :\text{nd}()𝒜_U`$ is an isomorphism. The action $`𝔾_m𝒜\text{ut}(,\phi )`$ is given by homotheties on $``$.
Set $`_X^\times =g_{}g^{}(𝔾_m)`$, where $`g:X^{(0)}X`$ is the inclusion of generic points, and let $`\pi :_X^\times 𝒟\text{iv}_X`$ and $`p:𝒜𝒟\text{iv}_X𝒫_X`$ be the natural surjections. Then the image of $`s_{}H^0(X,𝒫_X)`$ under the iterated coboundary map is given by the gerbe of $`_X^\times `$-liftings of the $`𝒟\text{iv}_X`$-torsor $`p^1(s_{})`$. This gerbe associates to each étale $`UX`$ the groupoid of pairs $`(𝒯,\psi )`$, where $`𝒯`$ is a $`_U^\times `$-torsor, and $`\psi :𝒯p^1(s_{})`$ is a $`\pi `$-morphism of torsors. Moreover, the action $`𝔾_m𝒜\text{ut}(𝒯,\psi )`$ is given by translation on $`𝒯`$.
We have to construct an equivalence between the preceding two stacks that is equivariant for the sign change map $`1:𝔾_m𝔾_m`$. First note that the stack of pairs $`(,\phi )`$ is $`𝔾_m`$-equivalent to the stack of triples $`(,,\psi )`$, where $``$ is a locally free $`𝒪_U`$-module, and $``$ is an almost invertible $`𝒪_U`$-module, and $`\psi :`$ is an isomorphism. The action $`𝔾_m𝒜\text{ut}(,,\psi )`$ is given by homotheties on $``$ and inverse homotheties on $``$.
Obviously, the local section $`s_{}\mathrm{\Gamma }(U,𝒫_X)`$ is nothing but the restriction of the global section $`s_{}\mathrm{\Gamma }(X,𝒫_X)`$. Now let $`_U^\times ()`$ be the sheaf of invertible meromorphic sections in $`|_U`$, which is a $`_U^\times `$-torsor. Dividing out the induced $`𝔾_m`$-action, we obtain a $`\pi `$-morphism
$$\phi :_U^\times ()p^1(s_{})=p^1(s_{}|_U).$$
Consequently, the functor $`(,,\psi )(_U^\times (),\phi )`$ gives the desired antiequivariant equivalence of $`𝔾_m`$-gerbes. ∎
For the rest of the section, we assume that $`X`$ satisfies Serre’s condition $`(S_1)`$, that is, $`X`$ has no embedded components. Let $`g:X^{(0)}X`$ be the inclusion of the generic points. We can relate generically trivial Brauer classes to torsors of Cartier divisors as follows. The exact sequence
$$0𝔾_mg_{}g^{}(𝔾_m)𝒟\text{iv}_X0,$$
together with $`R^1g_{}(𝔾_{m,X^{(0)}})=0`$ and $`\mathrm{Pic}(X^{(0)})=0`$, gives an exact sequence
(4)
$$0H^1(X,𝒟\text{iv}_X)H^2(X,𝔾_m)H^2(X^{(0)},𝔾_m).$$
Hence each $`\alpha \mathrm{Br}(X)`$ with $`g^{}(\alpha )=0`$ comes from a $`𝒟\text{iv}_X`$-torsor. To make this explicit, choose an Azumaya algebra $`𝒜`$, say of rank $`r`$, representing the class $`\alpha `$, and let $`f:BX`$ be the associated *Brauer–Severi* scheme. This is the $`^{r1}`$-bundle
$$B=\mathrm{Isom}(\text{at}_r(𝒪_X),𝒜)\times _{\mathrm{PGL}_r}^{r1}$$
on the étale site $`X_{\text{ét}}`$. Here we use the left $`\mathrm{PGL}_r`$-action coming from the canonical representation $`\mathrm{PGL}_r\mathrm{Aut}(^{r1})`$ described in Chap. 0 §5. Set $`P=B\times _XX^{(0)}`$, and pick an invertible $`𝒪_P`$-module $`𝒪_P(1)`$ of fiber degree one. This leads to a $`𝒟\text{iv}_X`$-torsor $`𝒯`$ as follows. Define
$$\mathrm{\Gamma }(X,𝒯)=\left\{\mathrm{cl}(,t)\right\},$$
where $``$ is an invertible $`𝒪_B`$-module, and $`t:𝒪_P(1)|_P`$ is an isomorphism. Here $`\mathrm{cl}(,t)`$ denotes isomorphism class, and two pairs $`(,t)`$ and $`(^{},t^{})`$ are called isomorphic if there is an isomorphism $`\varphi :^{}`$ with $`t\varphi =t^{}`$. Define $`𝒯`$ in the same way on the étale site.
By Proposition 21.2.11, the sections of $`𝒟\text{iv}_X`$ correspond to $`\mathrm{cl}(,s)`$, where $``$ is an invertible $`𝒪_X`$-module, and $`s:𝒪_{X^{(0)}}|_{X^{(0)}}`$ is a trivialization. The map
$$(,s),(,t)(f^{}(),f^{}(s)t)$$
turns $`𝒯`$ into a $`𝒟\text{iv}_X`$-torsor.
###### Proposition 1.6.
The Brauer class $`\mathrm{cl}(𝒜)H^2(X,𝔾_m)`$ is the opposite for the image of the torsor class $`\mathrm{cl}(𝒯)H^1(X,𝒟\text{iv}_X)`$ under the coboundary map in (4).
###### Proof.
As in the proof of Proposition 1.5, the Brauer class is given by the gerbe of trivializations $`\phi :\text{nd}()𝒜`$. This gerbe is equivalent to the gerbe of trivializations $`(,u)`$, where $`u:B(^{})`$ is an isomorphism. According to Chap. V Lemma 4.8.1, the latter gerbe is antiequivalent to the $`𝔾_m`$-gerbe of invertible $`𝒪_P`$-modules $``$ of fiber degree one.
Set $`_X^\times =g_{}g^{}𝔾_m`$. The image of the torsor class $`\mathrm{cl}(𝒯)H^1(X,𝒟\text{iv}_X)`$ is the gerbe of $`_X^\times `$-liftings $`\psi :𝒮𝒯`$ for the $`𝒟\text{iv}_X`$-torsor $`𝒯`$. Given an invertible $`𝒪_P`$-module $``$ of fiber degree one, we obtain an $`_X^\times `$-torsor
$$𝒮=\{(,t)t:𝒪_P(1)\stackrel{}{}_P\},$$
where the action is by multiplication on $`t`$. Clearly, the quotient $`𝒮/𝔾_m=\left\{\mathrm{cl}(,t)\right\}`$ is canonically isomorphic to $`𝒯`$, so we obtain a morphism of torsors $`\psi :𝒮𝒯`$. To see that $`(𝒮,\psi )`$ is a $`𝔾_m`$-equivalence, note that a $`_X^\times `$-lifting of the torsor $`𝒯`$ exists if and only if $`𝒯`$ is trivial, because $`H^1(X,_X^\times )=0`$. ∎
## 2. Generically trivial Brauer classes
In this section, we turn to normal surfaces. The task is to prove the following result, which is a major step towards showing $`\mathrm{Br}(X)=\mathrm{Br}^{}(X)`$.
###### Proposition 2.1.
Let $`X`$ be a separated normal algebraic surface. Then $`\mathrm{Br}(X)`$ contains each class $`\alpha \mathrm{Br}^{}(X)`$ that is generically trivial.
###### Proof.
We start with some preliminary reductions. By Chap. II Lemma 4, we may assume that the ground field $`k`$ is separably closed. Since $`X`$ is separated and of finite type, the Nagata Compactification Theorem gives a compactification $`X\overline{X}`$. By resolution of singularities, we may assume that $`\mathrm{Sing}(X)=\mathrm{Sing}(\overline{X})`$. As in Chap. II Theorem 2.1, each generically trivial Brauer class $`\alpha \mathrm{Br}^{}(X)`$ extends to $`\mathrm{Br}^{}(\overline{X})`$, so we may begin the proof with the additional assumption that $`X`$ is proper.
As discussed in Section 1, the exact sequence
$$0H^1(X,𝒟\text{iv}_X)H^2(X,𝔾_m)H^2(X^{(0)},𝔾_m)$$
shows that our cohomology class $`\alpha \mathrm{Br}^{}(X)`$ lies in $`H^1(X,𝒟\text{iv}_X)`$. The exact sequence
$$0𝒟\text{iv}_X𝒵_X^1𝒫_X0$$
yields an exact sequence
$$Z^1(X)H^0(X,𝒫_X)H^1(X,𝒟\text{iv}_X)0.$$
Choose a global section $`sH^0(X,𝒫_X)`$ mapping to $`\alpha `$. Since $`r\alpha =0`$ for some integer $`r>0`$, there is a global Weil divisor $`EZ^1(X)`$ mapping to the section $`rsH^0(X,𝒫_X)`$.
The sheaf $`𝒫_X`$ is supported by the singular locus $`\mathrm{Sing}(X)`$. For each singular point $`xX`$, the stalk is $`𝒫_{X,x}=\mathrm{Cl}(𝒪_{X,x}^h)`$, where $`𝒪_{X,x}𝒪_{X,x}^h`$ is the henselization (which is the strict localization, because $`k`$ is separably closed). Suppose $`x_1,\mathrm{},x_mX`$ are the singularities, and set
$$X^h=\mathrm{Spec}(\underset{i=1}{\overset{m}{}}𝒪_{X,x_i}^h)=\underset{i=1}{\overset{m}{}}\mathrm{Spec}(𝒪_{X,x_i}^h).$$
Then $`H^0(X,𝒫_X)=\mathrm{Cl}(X^h)`$. Choose a Weil divisor $`DZ^1(X^h)`$ representing the section $`sH^0(X,𝒫_X)`$, such that $`E|_{X^h}rD`$. According to Proposition 1.5, it suffices to construct a reflexive $`𝒪_X`$-module $``$ with
$$𝒪_{X^h}=\underset{i=1}{\overset{r}{}}𝒪_{X^h}(D),$$
for then $`𝒜=\text{nd}()`$ would be the desired Azumaya algebra.
Let $`f:YX`$ be a resolution of singularities, and $`Y_0Y`$ be the reduced exceptional curve. Set $`Y^h=Y\times _XX^h`$. The following crucial result, which is due to Gabber, tells us that certain vector bundles on $`Y^h`$ are already determined on suitable infinitesimal neighborhoods of $`Y_0`$.
###### Lemma 2.2 (Gabber).
Let $``$ be a family of locally free $`𝒪_{Y_0}`$-modules of fixed rank $`n0`$. Suppose that $``$ is, up to tensoring with line bundles, a bounded family. Then there is an exceptional curve $`RY`$ so that the map $`H^1(Y^h,\mathrm{GL}_n)H^1(R,\mathrm{GL}_n)`$ is bijective on the subsets of vector bundles whose restriction to $`Y_0`$ lies in $``$.
###### Proof.
Let $`𝒪_Y`$ be the ideal of $`Y_0Y`$. Since the intersection matrix for the irreducible components of $`Y_0`$ is negative definite, there is an exceptional curve $`AY`$ with support $`Y_0`$ so that $`𝒪_{Y_0}(A)`$ is ample. Then $`𝒪_A(A)`$ is ample as well. Since the family $`\{\text{nd}()\}`$ is bounded, there is an integer $`m_0>0`$ so that
(5)
$$H^1(A,\text{nd}()𝒪_A(mA))=0$$
for all $`mm_0`$ and all $``$.
Let $`𝔜Y`$ be the formal completion along $`Y_0Y`$. We first check the statement of the Lemma for the formal scheme $`𝔜`$ instead of $`Y^h`$. Note that the canonical map
$$H^1(𝔜,\mathrm{GL}_n)\underset{}{\mathrm{lim}}H^1(mA,\mathrm{GL}_n)$$
is bijective, as explained in , proof of Theorem 3.5. Let $`𝒥𝒪_Y`$ be the ideal of $`AY`$. The obstruction to lifting a vector bundle $``$ on $`mA`$ to $`(m+1)A`$ lies in
$$H^2(R,\text{nd}()𝒥^m/𝒥^{m+1})=0,$$
so the restriction maps $`H^1(𝔜,\mathrm{GL}_n)H^1(mA,\mathrm{GL}_n)`$ are surjective for all $`m0`$.
Fix an integer $`mm_0`$, and let $`,^{}`$ be two vector bundles on $`(m+1)A`$ that are isomorphic on $`mA`$ and whose restrictions to $`Y_0`$ belong to the family $``$. Choose an isomorphism $`\psi :|_{mA}^{}|_{mA}`$. Locally on $`(m+1)A`$, we can lift this isomorphism to an isomorphism $`^{}`$. The sheaf of such liftings is a torsor under
$$\text{om}(,^{})𝒥^m/𝒥^{m+1}\text{nd}()𝒥^m/𝒥^{m+1}\text{nd}()𝒪_A(mA).$$
This sheaf has no first cohomology by (5). The upshot is that a global lifting of $`\psi :|_{mA}^{}|_{mA}`$ exists. Consequently, for $`R=m_0A`$, the mapping $`H^1(𝔜,\mathrm{GL}_n)H^1(R,\mathrm{GL}_n)`$ is bijective on the subsets of vector bundles whose restriction belongs to the family $``$.
Finally, we pass to $`Y^h`$. By the Artin Approximation Theorem ( Thm. 3.5), the map $`H^1(Y^h,\mathrm{GL}_n)H^1(𝔜,\mathrm{GL}_n)`$ is injective and has dense image. So given a formal vector bundle $``$ with $`|_{Y_0}`$, we find a vector bundle $`^h`$ on $`Y^h`$ with $`^h|_R|_R`$. By the choice of $`R`$, this implies $`^h|_𝔜`$. ∎
We proceed with the proof of Proposition 2.1. Let $``$ be the family of vector bundles on $`Y_0`$ of rank $`r`$ that are free up to tensoring with line bundles, and choose an exceptional divisor $`RY`$ as in the preceding Lemma. Let $`D^{}Z^1(Y^h)`$ be the strict transform of $`DZ^1(X^h)`$. Let $`E^{}Z^1(Y)`$ be the unique Weil divisor that is the strict transform of $`EZ^1(X)`$ on $`YY_0`$, and satisfies $`E^{}rD^{}`$ on $`Y^h`$. Set $`=𝒪_R(D^{})`$. According to Lemma 2.2, we have to construct a locally free $`𝒪_Y`$-module $``$ with $`_R_{i=1}^r^{}`$. For then the double dual $`=f_{}()^{}`$ would be the desired reflexive $`𝒪_X`$-module. We shall construct such a vector bundle as an elementary transformation of the trivial bundle $`𝒪_Y^r`$.
Choose an ample divisor $`AY`$. Replacing the divisors $`D`$ and $`E`$ by $`D+f_{}(tA)`$ and $`E+f_{}(rtA)`$, respectively, does not change the class $`\alpha \mathrm{Br}^{}(X)`$. Choosing $`t0`$, we may assume that $`=𝒪_R(D^{})`$ is very ample, and that $`H^1(Y,𝒪_Y(E^{}R))=0`$. Next, choose pairwise disjoint effective Cartier divisors $`D_1,\mathrm{},D_rR`$, each one representing $``$. Let $`s_i\mathrm{\Gamma }(R,)`$ be the corresponding sections. Regard their product $`s_1\mathrm{}s_r`$ as a section of $`𝒪_R(E^{})`$. By construction, the group on the right in the exact sequence
$$H^0(Y,𝒪_Y(E^{}))H^0(R,𝒪_R(E^{}))H^1(Y,𝒪_Y(E^{}R))$$
is zero. Consequently, $`E^{}Z^1(Y)`$ is linearly equivalent to an effective divisor $`HY`$ with $`HR=D_1\mathrm{}D_r`$. Now choose a closed subset $`SHR`$ so that each Cartier divisor $`D_iH`$ is principal on $`HS`$. For each $`1ir`$, choose an exact sequence
(6)
$$0𝒪_H(C_i)\stackrel{t_i}{}𝒪_H\underset{ji}{}𝒪_{D_j}0$$
for certain Cartier divisors $`C_i\mathrm{Div}(H)`$ supported by $`S`$. As explained in , p. 152, it suffices to construct the desired Azumaya $`𝒪_X`$-algebra on $`Xf(S)`$. Hence it suffices to construct the desired locally free $`𝒪_Y`$-module $``$ on $`YS`$, and we may replace $`X`$, $`Y`$ by the complements $`Xf(S)`$, $`YS`$, respectively. Now the exact sequence (6) induces an exact sequence
$$𝒪_Y\stackrel{t_i}{}𝒪_H\underset{ji}{}𝒪_{D_j}0.$$
The map $`t=(t_1,\mathrm{},t_r):_{i=1}^r𝒪_Y𝒪_H`$ is surjective, because the $`D_i`$ are pairwise disjoint. The exact sequence
$$0\underset{i=1}{\overset{r}{}}𝒪_Y\stackrel{t}{}𝒪_H0$$
defines a locally free $`𝒪_Y`$-module $``$, because the cokernel $`𝒪_H`$ has homological dimension $`\mathrm{hd}(𝒪_H)=1`$. Restricting to the curve $`RY`$, we obtain an exact sequence
$$𝒯\text{or}_{𝒪_Y}^1(𝒪_H,𝒪_R)_R\underset{i=1}{\overset{r}{}}𝒪_R\stackrel{t_R}{}𝒪_{HR}0.$$
The term on the left is zero, because the curves $`H,RY`$ have no common components. By construction, $`t_i|_{D_j}=0`$ for $`ij`$, so the induced surjection $`t_R`$ is a diagonal matrix of the form
$$t_R=\left(\begin{array}{ccc}t_1|_{D_1}& & 0\\ & \mathrm{}& \\ 0& & t_r|_{D_r}\end{array}\right):\underset{i=1}{\overset{r}{}}𝒪_R\underset{i=1}{\overset{r}{}}𝒪_{D_i}=𝒪_{HR}.$$
Consequently $`_R_{i=1}^r^{}`$. Hence $``$ is the desired locally free $`𝒪_Y`$-module. ∎
###### Remark 2.3.
In my first proof of Proposition 2.1, I used a result of Treger ( Prop. 3.5), which states that the map $`H^1(Y^h,\mathrm{GL}_n)H^1(R,\mathrm{GL}_n)`$ is bijective for a suitable exceptional curve $`RY`$. As Gabber pointed out, this statement is wrong for $`n2`$. His counterexample goes as follows.
Let $`X=\mathrm{Spec}(A)`$ be a complete normal local surface singularity, and $`f:YX`$ a resolution of singularity, and $`Y_mY`$ the infinitesimal neighborhoods of the exceptional curve $`Y_0`$. Suppose $`n=2`$ for simplicity, and assume there is an exceptional curve $`R`$ as above. Fix an ample invertible $`𝒪_Y`$-module $``$, set $`_m=|_{Y_m}`$, and choose $`m0`$ with $`RY_m`$. The exact sequence
$$0_0^t(Y_m)_{m+1}^t_m^t0$$
yields an exact sequence
$$H^0(Y_m,_m^t)H^1(Y_0,_0^t(Y_m))H^1(Y_{m+1},_{m+1}^t)H^1(Y_m,_m^t)0.$$
For $`t0`$, the group on the left is zero, and $`H^1(Y_0,_0^t(Y_m))`$ is nonzero. It follows that there is a nonzero class $`\zeta H^1(Y,^t)`$ restricting to zero in $`H^1(Y_m,_m^t)`$. This defines a nonsplit extension
(7)
$$0^t\stackrel{\psi }{}𝒪_Y0,$$
which splits on $`R`$. By the defining property of the curve $`R`$, there is a bijection $`\varphi :𝒪_Y^t`$. The composition
$$𝒪_Y\stackrel{\varphi }{}\stackrel{\psi }{}𝒪_Y$$
is surjective on $`Y_m`$, because $`H^0(Y_m,_m^t)=0`$. By the Nakayama Lemma, the composition is surjective on the formal completion, and hence on $`Y`$ as well. So the extension (7) splits, contradicting $`\zeta 0`$.
## 3. Elementary transformations of Brauer–Severi schemes
We come to the main result of this paper.
###### Theorem 3.1.
Let $`X`$ be a separated geometrically normal algebraic surface. Then we have $`\mathrm{Br}(X)=\mathrm{Br}^{}(X)`$.
###### Proof.
According to Chap. II Lemma 4, we may assume that the ground field is algebraically closed. Fix a class $`\alpha \mathrm{Br}^{}(X)`$. In light of Proposition 2.1, it suffices to construct an Azumaya algebra representing $`\alpha \mathrm{Br}^{}(X)`$ generically. Choose a resolution of singularities $`f:YX`$, and let $`𝔜Y`$ be the formal completion along the reduced exceptional curve $`Y_0Y`$. According to Theorem 2.1, there is an Azumaya $`𝒪_Y`$-algebra $`𝒜`$ representing $`f^{}(\alpha )`$. The task now is to choose such an Azumaya algebra so that the formal vector bundle $`𝒜|_𝔜`$ is trivial. For then, as explained in p. 152, the $`𝒪_X`$-algebra $`f_{}(𝒜)`$ is an Azumaya algebra, which represents $`\alpha \mathrm{Br}^{}(X)`$ generically. Note that we may remove finitely many closed smooth points from $`X`$.
First, we check that the formal Azumaya algebra $`𝒜|_𝔜`$ is trivial. Since the ground field is separably closed, there is a locally free $`𝒪_{Y_0}`$-module $`_0`$ and an isomorphism $`\phi _0:\text{nd}(_0)𝒜|_{Y_0}`$. The following argument due to Gabber shows that the pair $`(_0,\phi _0)`$ extends over all infinitesimal neighborhoods $`Y_0Y_n`$. Let $`d(𝒜|_{Y_n})`$ be the $`𝔾_m`$-gerbe of trivializations of $`𝒜|_{Y_n}`$. The restriction map gives a cartesian functor $`d(𝒜|_{Y_{n+1}})d(𝒜|_{Y_n})`$. Consider the corresponding stack of liftings of trivializations. This is a gerbe for the abelian sheaf $`^{n+1}/^{n+2}`$ on the étale site of $`Y_0`$, where $`=𝒪_Y(Y_0)`$. Since $`H^2(Y_0,^{n+1}/^{n+2})=0`$, the gerbe of liftings is trivial. Consequently, we have $`𝒜|_𝔜\text{nd}()`$ for some locally free $`𝒪_𝔜`$-module $``$.
The idea now is to make an elementary transformation along a curve $`HY`$, so that $``$ becomes free on certain infinitesimal neighborhood $`Y_0R`$. Furthermore, we shall choose the curve $`R`$ so that the freeness of $`_R`$ implies the freeness of $``$. This requires some preparation. Let me introduce three numbers $`m,k,q`$ depending on $`Y`$ and $``$. First, set $`m=\mathrm{rank}()`$. Second, let $`k1`$ be the order of the cokernel for the map $`\mathrm{Pic}(Y)\mathrm{NS}(Y_0)`$ onto the Néron–Severi group. Third, define $`q=1`$ in characteristic zero. In characteristic $`p>0`$, let $`q>0`$ be a $`p`$-th power so that the unipotent part of $`\mathrm{Pic}^0(𝔜)`$ is $`q`$-torsion. This works, because $`\mathrm{Pic}_𝔜^0`$ is an algebraic group scheme.
Now set $`r=mkq`$, and let $``$ be the family of locally free $`𝒪_{Y_0}`$-modules of rank $`r`$ which are free up to tensoring with line bundles. Choose a curve $`RY`$ as in Lemma 2.2. Finally, let $`X\overline{X}`$ be a compactification with $`\mathrm{Sing}(X)=\mathrm{Sing}(\overline{X})`$, and let $`Y\overline{Y}`$ be the corresponding compactification.
###### Claim.
We can modify the Azumaya algebra $`𝒜`$ and the vector bundle $``$ so that $``$ is a globally generated formal vector bundle of rank $`r`$, and that there is a very ample invertible $`𝒪_{\overline{Y}}`$-module $``$ with $`det()=_𝔜`$ and $`H^1(\overline{Y},(R))=0`$.
###### Proof.
Tensoring $``$ with an ample line bundle, we archive that $``$ is globally generated and that $`det()`$ is ample. Next, we replace the Azumaya algebra $`𝒜`$ by $`𝒜\text{nd}(𝒪_Y^k)`$. This replaces the vector bundle $``$ by $`^k`$, and $`det()`$ by $`det()^k`$. Consequently, we may assume that $`det()`$ is numerically equivalent to some ample invertible $`𝒪_𝔜`$-module $`_𝔜`$. Twisting $``$ by a suitable power of $`_𝔜`$, we may assume that $`_𝔜`$ extends to an ample invertible $`𝒪_{\overline{Y}}`$-module $``$ with $`H^1(\overline{Y},^s(R))=0`$ for all integers $`s>0`$.
In characteristic zero, $`\mathrm{Pic}^0(𝔜)`$ is a divisible group. Hence we find an invertible $`𝒪_𝔜`$-module $``$ with $`^{mk}det()_𝔜`$. Replacing $``$ by $``$, we have $`det()=_𝔜`$. Now assume that we are in characteristic $`p>0`$. Set $`G=\mathrm{Pic}_{𝔜/k}^0`$, and let $`G^{}G`$ be the unipotent part. Then the quotient $`G^{\prime \prime }=G/G^{}`$ is semiabelian, and $`G^{\prime \prime }(k)`$ is a divisible group. As in characteristic zero, we may assume that $`det()_𝔜^{}`$ lies in the unipotent part of $`\mathrm{Pic}^0(𝔜)`$, which is a $`q`$-torsion group. Passing to $`𝒜\text{nd}(𝒪_X^q)`$ and $`^q`$, we are done. ∎
We continue proving Theorem 3.1. Set $`r=\mathrm{rank}()`$ and $`\mathrm{\Gamma }=\mathrm{\Gamma }(R,_R)`$. The canonical surjection $`\mathrm{\Gamma }𝒪_R_R`$ yields a morphism $`\phi :R\mathrm{Grass}_r(\mathrm{\Gamma })`$ into the Grassmannian of r-dimensional quotients. Choose a generic $`r`$-dimensional subvector space $`\mathrm{\Gamma }^{}\mathrm{\Gamma }`$. For each integer $`k0`$, let $`G_k\mathrm{Grass}_r(\mathrm{\Gamma })`$ be the subscheme of surjections $`\mathrm{\Gamma }\mathrm{\Gamma }^{\prime \prime }`$ such that the composition $`\mathrm{\Gamma }^{}\mathrm{\Gamma }^{\prime \prime }`$ has rank $`k`$. Note that $`G_{r1}`$ is a reduced Cartier divisor, and that $`G_{r2}`$ has codimension four (see , Sec. II.2).
By the dimensional part of Kleiman’s Transversality Theorem ( Thm. 2), which is valid in all characteristics, the map $`\phi :R\mathrm{Grass}_r(\mathrm{\Gamma })`$ is disjoint to $`G_{r2}`$ and passes through $`G_{r1}`$ in finitely many points. The upshot of this is that the quotient of the canonical map $`\mathrm{\Gamma }^{}𝒪_X`$ is an invertible sheaf on some Cartier divisor $`DR`$. Consequently, we have constructed an exact sequence
$$0𝒪_R^r_R𝒪_D0.$$
In other words, the trivial vector bundle is the elementary transformation of $``$ with respect to the surjection $`_R𝒪_D`$. In geometric terms: Blowing up $`(_R)`$ along the section $`(𝒪_D)(_D)`$ and contracting the strict transform of $`(_D)`$ yields $`_R^{r1}`$ (see Thm. 1.4).
We seek to extend this elementary transformation from the curve to the surface. Note that $`_R=det(_R)=𝒪_R(D)`$. The exact sequence
$$H^0(\overline{Y},)H^0(R,_R)H^1(\overline{Y},(R))=0$$
implies that $`D=HR`$ for some ample curve $`HY`$. Since the ground field $`k`$ is algebraically closed, Tsen’s Theorem gives $`H^2(H,𝔾_m)=0`$. Removing finitely many smooth points from the open subset $`X\overline{X}`$, we may assume that $`𝒜_H=\text{nd}(𝒪_H^r)`$. In other words, if $`PY`$ is the Brauer–Severi scheme corresponding to $`𝒜`$, we have $`P_H=_H^{r1}`$. Let $`A`$ be the semilocal ring of the curve $`H`$ corresponding to the closed points $`DH`$. The section $`(𝒪_D)P_D`$ is given by a surjection $`A^rA/I`$, where $`IA`$ is the ideal of $`D`$. By Nakayama’s Lemma, this lifts to a surjection $`A^rA`$. So, if we shrink $`X`$ further, we can extend the section $`(𝒪_D)P_D`$ to a section $`SP_H`$.
Let $`h:\widehat{P}P`$ be the blowing-up with center $`SP`$, and let $`E\widehat{P}`$ be the strict transform of the Cartier divisor $`P_HP`$. I claim that there is a birational contraction $`\widehat{P}P^{}`$ contracting precisely the fibers of $`EH`$ to points such that $`P^{}`$ is a Brauer–Severi scheme. Over suitable étale neighborhoods, this follows from Theorem 1.4. You easily check that these contractions glue together and define a contraction in the category of schemes.
By construction, the new Brauer–Severi scheme $`P^{}X`$ has a trivial restriction $`P_R^{}=_R^{r1}`$. If $`𝒜^{}`$ is the Azumaya algebra corresponding to the Brauer–Severi scheme $`P^{}`$, this implies $`𝒜_R^{}=\text{nd}(𝒪_R^r)`$. By the choice of the curve $`RY`$, this forces the formal vector bundle $`𝒜_𝔜^{}`$ to be free. Consequently, the direct image $`f_{}(𝒜^{})`$ is an Azumaya $`𝒪_X`$-algebra representing the class $`\alpha \mathrm{Br}^{}(X)`$ generically. ∎
###### Remark 3.2.
The proof works for separated geometrically normal 2-dimensional algebraic spaces as well. This is because their resolutions are schemes.
###### Question 3.3.
The hypothesis of *geometric* normality annoys me. What happens for separated normal 2-dimensional noetherian schemes that are of finite type over nonperfect fields, or over Dedekind rings, or have no base ring at all?
## 4. Existence of vector bundles
Given a scheme $`X`$, one might ask whether $`X`$ admits a nonfree vector bundle. In dimension two, we can use Brauer groups to obtain a positive answer:
###### Theorem 4.1.
Let $`X`$ be a proper normal surface over a field $`k`$. Then there is a locally free $`𝒪_X`$-module of finite rank that is not free.
###### Proof.
Seeking a contradiction, we assume that each vector bundle is free. We may assume that $`k=\mathrm{\Gamma }(X,𝒪_X)`$. First, I reduce to the case that the ground field $`k`$ is separably closed. Let $`kL`$ be a separable closure, and let $`_L`$ be a vector bundle on $`X_L=XL`$, say of rank $`r0`$. Then there is a finite separable field extension $`kK`$, say of degree $`d1`$, such that $`_L`$ comes from a vector bundle $`_K`$ on $`X_K`$. Let $`p:X_KX`$ be the canonical projection. Then $`=p_{}(_K)`$ is a vector bundle of rank $`dr`$, hence free by assumption. This gives
$$\mathrm{\Gamma }(X_K,_K)=\mathrm{\Gamma }(X,)k^{rd}K^r.$$
Now you easily choose $`r`$ sections of $``$ that are linearly independent over $`K`$, which gives the desired trivialization of $`_K`$.
From now on, assume that $`k`$ is separably closed. Note that $`\mathrm{Pic}(X)=0`$, so $`X`$ is nonprojective, hence it must contain some singularities. Let $`f:YX`$ be a resolution of singularities. Choose an exceptional divisor $`RX`$ so that $`\mathrm{Pic}(R)=\mathrm{Pic}(𝔜)`$, where $`𝔜Y`$ is the formal completion along the exceptional curve. The spectral sequence for $`𝔾_{m,X}=f_{}(𝔾_{m,Y})`$ gives an exact sequence
$$0\mathrm{Pic}(Y)\mathrm{Pic}(R)H^2(X,𝔾_m)H^2(Y,𝔾_m).$$
Set $`G=\mathrm{Pic}^0(R)/\mathrm{Pic}^0(Y)`$, and let $`H\mathrm{NS}(Y)`$ be the kernel of the restriction map $`\mathrm{NS}(Y)\mathrm{NS}(R)`$ for the Néron–Severi groups. The snake lemma gives an inclusion $`G/HH^2(X,𝔾_m)`$. To proceed, we need a well-known fact:
###### Lemma 4.2.
For each $`l>0`$ prime to $`\mathrm{char}(k)`$, the group $`\mathrm{Pic}^0(R)`$ is $`l`$-divisible.
###### Proof.
We have an exact sequence
$$0\mathrm{Pic}^0(R)\mathrm{Pic}_{R/k}^0(k)\mathrm{Br}(k).$$
Since $`\mathrm{Br}(k)=0`$, we have to see that the multiplication morphism $`l:PP`$ is surjective on the smooth algebraic group scheme $`P=\mathrm{Pic}_{R/k}^0`$. Since $`P`$ is connected, it suffices to check that $`l:PP`$ is open. The completion at the origin $`0P`$ is a formal group, given by a formal power series ring $`k[[X_1,\mathrm{},X_n]]`$ together with $`n`$ formal power series
$$F_i(X_1,\mathrm{},X_n,Y_1,\mathrm{},Y_n)=X_i+Y_i+\text{terms of higher order.}$$
Multiplication by $`l`$ is given by $`l1`$ substitutions
$$[l]^{}X_i=F_i([l1]^{}X_1,\mathrm{},[l1]^{}X_n,X_1,\mathrm{},X_n)lX_i,$$
modulo terms of higher order. Since $`l`$ is prime to the characteristic of the ground field, this is bijective. Consequently, $`l`$ is étale on $`𝒪_{P,0}^{}`$, hence étale on $`P`$, and therefore open. ∎
We continue with the proof of the Proposition. If $`G=0`$, then $`X`$ would have nontrivial line bundles (, Prop. 4.2), which is absurd. So $`G`$ is a nonzero $`l`$-divisible group. On the other hand, $`H`$ is finitely generated. We conclude that $`G/HH^2(X,𝔾_m)`$ contains many torsion points. Consequently, $`\mathrm{Br}^{}(X)`$ contains nonzero generically trivial classes. By Proposition 2.1 there is a nontrivial Azumaya $`𝒪_X`$-algebra $`𝒜`$.
Setting $`r^2=\mathrm{rank}(𝒜)`$, we have $`𝒜𝒪_X^{r^2}`$ as $`𝒪_X`$-module. The multiplication map $`𝒜𝒜𝒜`$ and the unit $`𝒪_X𝒜`$ induce a $`k`$-algebra structure on $`A=\mathrm{\Gamma }(X,𝒜)`$. You easily check that $`A\kappa (x)𝒜(x)`$ for each point $`xX`$. So $`A`$ is a central simple $`k`$-algebra, which is trivial because $`k`$ is separably closed. Consequently $`𝒜=A𝒪_X`$ is also trivial, contradiction. ∎
###### Remark 4.3.
It might easily happen that $`X`$ has trivial Picard group . However, the preceding results ensures the existence of vector bundles of higher rank. |
warning/0003/hep-th0003023.html | ar5iv | text | # Embedding Branes in Flat Two-time Spaces
## 1 Embedding: The geometry
In this section we describe the embedding of a $`SO(n)`$ invariant $`p`$-brane in a $`(D+2)=(n+p+3)`$-dimensional spacetime. We will obtain the embedding by requiring that the known metric of the brane is obtained from a flat $`(D,2)`$ metric, i.e. we demand that the embedding is isometric . The $`D`$-dimensional $`p`$-brane geometry can generally be described by a metric of the form
$`ds^2`$ $`=`$ $`A(r)^2\left[dt^2+dx_pdx_p\right]+`$ (1)
$`+B(r)^2dr^2+C(r)^2d\mathrm{\Omega }_{n1}^2,`$
where $`dx_pdx_p`$ is the $`p`$-dimensional spacelike part of the worldvolume, and $`d\mathrm{\Omega }_{n1}^2`$ is the metric for the $`n`$-sphere.
On the embedding space $`\mathrm{IE}^{(D,2)}`$ we take cartesian coordinates $`X^M`$, with $`M=0,\mathrm{},D+1`$, which we divide as $`X^M=(X^\mu ,X^{p+1},X^{p+2},X^\alpha )`$ (with $`\mu =0,\mathrm{},p`$, $`\alpha =p+3,\mathrm{},D+1`$). Using these coordinates, the metric looks as follows
$`ds^2`$ $`=`$ $`(dX^0)^2+(dX^1)^2+\mathrm{}+(dX^{p+1})^2+`$ (2)
$`(dX^{p+2})^2+\mathrm{}+(dX^{D+1})^2.`$
To obtain the two embedding constraints, we start by making a change of the $`(D+2)`$-dimensional coordinates, such that we make a subgroup $`SO(p,1)\times SO(n)SO(p+n+1,2)`$ manifest. This is achieved by using a mixture of hyperspherical and horospherical coordinates $`\{\rho ,z,x^\mu ,\beta ,n^\alpha \}`$
$`X^{}X^{p+2}X^{p+1}={\displaystyle \frac{\rho }{z}},`$
$`X^+X^{p+2}+X^{p+1}=\rho z+{\displaystyle \frac{\rho }{z}}x^\mu x_\mu ,`$
$`X^\mu =\rho {\displaystyle \frac{x^\mu }{z}},X^\alpha =\beta n^\alpha ,`$ (3)
where $`n^\alpha `$ ($`_\alpha (n^\alpha )^2=1`$) parametrise the sphere $`S^{n1}`$. In these new coordinates, the metric reads
$$ds^2=\frac{\rho ^2}{z^2}\left[dx^\mu dx_\mu +dz^2\right]d\rho ^2+d\beta ^2+\beta ^2dn^\alpha dn^\alpha .$$
(4)
Comparing (4) and (1) we identify $`dx^\mu dx_\mu `$ with $`dt^2+dx_p.dx_p`$ and $`dn^\alpha dn^\alpha `$ with $`d\mathrm{\Omega }_{n1}^2`$. Then $`\beta `$, $`\rho `$ and $`z`$ are functions of $`r`$ and are still to be determined. The comparison gives
$`\beta =C(r),{\displaystyle \frac{\rho }{z}}=A(r),`$
$`d\rho ^2+{\displaystyle \frac{\rho ^2}{z^2}}dz^2+d\beta ^2=B(r)^2dr^2.`$ (5)
The differential equation can be rewritten to give
$$\frac{C^2B^2}{A^{}}=(\rho z)^{}F^{}.$$
(6)
From all this we can derive the following embedding functions
$`X^{}`$ $`=`$ $`A(r),X^+=F(r)+A(r)x^\mu x_\mu `$
$`X^\mu `$ $`=`$ $`A(r)x^\mu ,X^\alpha =C(r)n^\alpha .`$ (7)
We can, furthermore, express the constraints in terms of the $`X^M`$ coordinates only. Denoting the inverse function with an overbar, i.e., $`\overline{f}(f)=f(\overline{f})=\text{ }\text{}`$, we can write $`r=\overline{A}(X^{})`$. Thus, our two constraints are
$`\varphi _1`$ $`=`$ $`X^{}X^+X^\mu X_\mu X^{}F(\overline{A}(X^{}))=0`$
$`\varphi _2`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}(X^\alpha )^2\left[C(\overline{A}(X^{}))\right]^2=0.`$ (8)
These constraints are therefore determined by the functions $`A`$, $`C`$ and $`F`$. The latter is determined up to a constant by (6) in terms of $`A`$, $`B`$ and $`C`$.
Note that so far there is no definition of the radial variable $`r`$. We can use different parametrizations, e.g. it will turn out that in some cases it is useful to take $`A`$ or $`C`$ itself as the radial variable. In the standard brane cases, the functions $`A`$, $`B`$ and $`C`$ will take the form of some harmonic function to some power in the transverse space of the brane. We will further adopt the name $`r`$ for that transverse coordinate, use just the name $`A`$ for the parameter in the first mentioned parametrization, and use $`R`$ for the radial coordinate such that $`C(R)=R`$.
From now on we will assume that the functions $`A`$, $`B`$ and $`C`$ are indeed harmonic functions in $`n`$ dimensions with a flat limit at $`r\mathrm{}`$. For non-dilatonic D- and M-branes, they are of the following form
$`H=\left(1+{\displaystyle \frac{1}{r^\kappa }}\right),A(r)=H^{\frac{1}{p+1}}`$
$`B(r)=H^{\frac{1}{\kappa }},C(r)=rH^{\frac{1}{\kappa }}`$ (9)
where $`\kappa n2=Dp3`$. Here a priori $`r>0`$ and $`r=0`$ corresponds to the horizon, but we will come back to this later.
With this explicit form for the functions $`A`$, $`B`$ and $`C`$ we can evaluate the function $`F`$. Using (6) we get
$$F^{}(r)=wr^{1\kappa }(1+r^\kappa )^{\frac{2}{\kappa }+\frac{1}{p+1}1}(1+2r^\kappa ),$$
(10)
(where $`w=\frac{p+1}{\kappa }`$), which can be integrated to give (up to a constant)
$`F(r)`$ $`=`$ $`\frac{w}{\kappa }[B_{\frac{r^\kappa }{r^\kappa +1}}(\frac{1}{p+1},1\frac{2}{\kappa })+`$ (11)
$`+2B_{\frac{r^\kappa }{r^\kappa +1}}(\frac{p}{p+1},\frac{2}{\kappa })].`$
Here we used the incomplete Beta function
$`B_x(a,b)`$ $`=`$ $`{\displaystyle _0^x}t^{a1}(1t)^{b1}𝑑t=`$
$`=`$ $`a^1x^a{}_{2}{}^{}F_{1}^{}(a,1b;a+1;x),`$
which is defined for $`0<x1`$. This means that $`F(r)`$ is well defined in the region $`r>0`$, which is what we were looking for.
Note that near the horizon ($`r0`$), where the brane geometry is well described by $`AdS_{p+2}\times S_{Dp2}`$ , we get $`Fw^2r^{\frac{1}{w}}`$ and the embedding functions (7) reduce to those used in .
Using the embedding (7) we can now study the global properties of the brane geometries. Before considering the higher dimensional D- and M-branes, let us first look at the simpler example of the extreme Reissner–Nordstrøm (RN) black hole. (A large list of embedding functions for other solutions of General Relativity is given in ). The RN black hole fits our general embedding scheme with $`D=4`$ and $`p=0`$, $`\kappa =1`$, $`w=1`$. Here, rather than working with the radial variable $`r`$ as in (9), we use the variable $`Rr+1`$, which has the property $`C(R)=R`$. Then the functions $`A`$ and $`B`$ are given by $`A(R)=B(R)^1=11/R`$. The horizon is now at $`R=1`$ and $`R=0`$ corresponds to the singularity. Using (6), we then find
$$F_{RN}(R)=\frac{1}{R1}3RR^24\mathrm{log}|R1|.$$
The entire Reissner–Nordstrøm black hole geometry can be drawn using parametrization (7) as is shown in figure 1. We only draw the relevant directions ($`X^{}`$, $`X^+`$, and $`X^0`$), which basically means we only draw the $`R`$ and $`t`$ coordinates of the black hole (every point in the graph should be thought of as a 2-sphere). The lines in the graph are therefore constant $`t`$ and constant $`R`$ lines.
We can read off the following global features from the picture. The geometry consists of 2 distinct regions: region I, the asymptotically flat region for $`R>1`$ which corresponds to $`X^{}>0`$. For big $`R`$ the surface flattens and $`X^{}1`$, which is the flat limit. Region II is the region inside the horizon ($`X^{}<0`$), the singularity ($`R=0`$) corresponds to $`X^{}\mathrm{}`$. This is the global picture we recognize from the familiar Penrose diagram for extreme RN black hole as can be found in .
The two regions are connected in an $`AdS`$-throat. It seems that these two regions are disconnected, the constant time lines all diverge near $`X^{}=0`$ and never cross the horizon, but this is just an artifact of the parametrization. Actually, we know that the near-horizon geometry is equivalent to $`AdS_2`$, which is known to have no problems at its ’horizon’. Indeed, a different parametrization exists (the advanced or retarded Finkelstein coordinates) in which lightlike geodesics pass smoothly through the horizon into the interior region, as is depicted in figure 2.
One of the features of $`AdS`$ spaces is that they admit closed timelike curves. The usual remedy for this is to consider the covering space $`CAdS`$ instead of $`AdS`$ itself. Looking at figures 1 and 2 we see that the RN black hole geometry suffers from the same problem, it admits closed timelike curves. Again this is remedied by considering the covering space. The result of this of course is that the space then consists of multiple universes.
Let us now move to the non-dilatonic branes. As discussed in , the general brane solution case (9) can be divided in two classes: $`p`$ odd or $`p`$ even, with quite different global properties.
Let us first consider the $`p`$ *odd* case. In the exterior region ($`r>0`$), the function $`A(r)`$ is analytic and positive and vanishes as $`r0`$. If we take $`A`$ to be our new radial variable instead of $`r`$, we see that $`A`$ can be continued through the horizon to negative $`A`$ . The range of $`A`$ is from -1 to 1. The analytic extension of the metric is
$`ds^2=A^2dx_\mu dx^\mu +(1A^{p+1})^{\frac{2}{\kappa }}\times `$ (12)
$`\times `$ $`\left[w^2\left(1A^{p+1}\right)^2A^2dA^2+d\mathrm{\Omega }^2\right],`$
which is even in $`A`$. This leads to
$`F(A)`$ $`=`$ $`\frac{w}{\kappa }(\text{sign }A)[B_{A^{p+1}}(\frac{1}{p+1},1\frac{2}{\kappa })+`$ (13)
$`+B_{A^{p+1}}(\frac{p}{p+1},\frac{2}{\kappa })].`$
The embedding functions (7) are then odd in $`A`$. This means that the embedded space is symmetric around the horizon and completely nonsingular. For the non-dilatonic branes, the D3 and M5 fit this picture. The embedding, depicted in figure 3 for the D3-brane case, nicely shows these features (an analogous picture for the M5-brane can be found in ). It is clearly visible there is no interior region, just two symmetric ’exterior’ regions connected in the AdS-throat as was expected from the Penrose diagram .
In the $`p`$ *even* case, the metric and embedding functions are neither even nor odd. It is useful in this case to adopt so-called Schwarzschild coordinates, defined by $`R^\kappa =r^\kappa +1`$. In these coordinates the horizon (which is still a coordinate singularity) is at $`R=1`$. At $`R=0`$ there is a true curvature singularity. Expressed in this coordinate, $`A(R)`$ can be continued through the horizon into negative $`A`$ and its range is $`\{\mathrm{},1\}`$. As already stated in , the Penrose diagram for these spaces is equivalent to the extreme Reissner–Nordstrøm diagram.
The embedding of the M2-brane metric illustrates these features. The expression (11) of $`F`$ is only well defined in the region $`R>1`$. It is not possible to find a continuous expression for $`F`$ valid in both regions ($`0<R1`$ and $`R>1`$). But, nevertheless, a continuous embedding is obtained using in the interior region
$`F(R<1)`$ $`=`$ $`\frac{w}{\kappa }\left[B_{R^\kappa }\right(\frac{1}{p+1}+\frac{2}{\kappa },\frac{1}{p+1})+`$
$`2B_{R^\kappa }(\frac{1}{p+1}+\frac{3}{\kappa },\frac{1}{p+1})].`$
The global properties of the M2-brane are qualitatively the same as those of the RN black hole depicted in figure 1. We refer to for the M2-picture.
## 2 The brane action
We would like to write the action of a brane placed in the background of other branes using the embedding of the previous section. A typical (schematic) form of the action is
$`S_{p+1}`$ $`=`$ $`{\displaystyle _W}d^{p+1}\xi \sqrt{det𝒢_{\mu \nu }}+{\displaystyle _B}\mathrm{\Omega }_{(p+2)}+`$ (14)
$`+{\displaystyle _W}d^{p+1}\xi [\lambda _1\varphi _1+\lambda _2\varphi _2],`$
where $`W=B`$ is the $`(p+1)`$-dimensional world volume of the brane. The expression for $`𝒢_{\mu \nu }`$ differs in each case. For example, for Dp-branes $`𝒢_{\mu \nu }_\mu X^M_\nu X^N\eta _{MN}+_{\mu \nu }`$, with $`_{\mu \nu }`$ the field strength of the gauge field living on the world volume of the brane. The fields $`\lambda _1,\lambda _2`$ are two Lagrange multipliers implementing the constraints (8). $`\mathrm{\Omega }_{(p+2)}(X^M)`$ is a function of the forms coupling to the brane, such that it reduces to the appropriate Wess–Zumino term when projected onto the physical hypersurface. Its explicit form will be determined for the D3-brane case in section 2.2. An analogous treatment for M2 and M5 is given in .
### 2.1 Embedding the field strength
Let us now try to embed the field strengths appearing in the Wess–Zumino term. We will assume that a brane (extended in $`p`$ spatial directions) fluctuating in a spacetime with two times should evolve in both time directions, and therefore couple to a (p+3)-form field strength. We assume therefore that the $`(D+2)`$-dimensional theory can be coupled to a rank $`p+3`$ electric field strength $`K_e`$, and to a rank $`n`$ magnetic field strength $`K_m`$.
This ansatz is the most natural one for the D3-brane, because in this case the $`10`$-dimensional self-dual field strength is extended to a self-dual field strength in $`12`$ dimensions. If there would be a supergravity theory in $`D=12`$, the bosonic configuration with flat $`(10,2)`$ space and a constant self-dual field strength would solve the equations of motion. This is obvious for the Maxwell equation (there can be no Chern–Simons terms built from a 5 form potential in 12 dimensions and so the Maxwell equation would take the standard form), but for the Einstein equations it is only true because the field strength is self-dual. In a $`D`$-dimensional spacetime with zero or two times, a self-dual field strength has a vanishing energy momentum tensor for $`D=4`$ mod 4. (For Lorentzian signature it is $`D=2`$ mod 4). What this would mean is that the ten-dimensional D3-brane solution would just be the projection to a complicated hypersurface of an almost trivial 12 dimensional supergravity solution.
Let us start by analysing how an electric $`(p+2)`$-form field strength $`F^{(p+2)}`$ gets embedded in the $`(D+2)`$-dimensional space. Our aim is to obtain $`F`$ as a restriction of a $`p+3`$-form $`K^{(p+3)}`$ to the $`D`$-dimensional hypersurface $`\mathrm{\Sigma }`$. A general non-dilatonic brane is described in $`D`$ dimensions by the fields
$`ds^2`$ $`=`$ $`H^{\frac{2}{p+1}}\left[dt^2+dx_1^2+\mathrm{}+dx_p^2\right]+`$
$`+H^{\frac{2}{\kappa }}\left[dr^2+r^2d\mathrm{\Omega }_{Dp2}^2\right],`$
$`G_{01\mathrm{}p}`$ $`=`$ $`H^1=A^{p+1},`$
$`\mathrm{\Phi }`$ $`=`$ $`0,`$ (15)
using the notation of section 1. We can write the (electric) field strength as ($`FdG`$)
$$F=(p+1)A^pA^{}drdtdx_1\mathrm{}dx_p,$$
(16)
where the prime denotes differentiation with respect to $`r`$.
To find the embedding, we start by considering a constant $`(p+3)`$-form in $`D+2`$ dimensions
$$K_e=\frac{p+1}{(p+3)!}ϵ_{\mu _0^{}\mathrm{}\mu _{p+2}^{}}dX^{\mu _0^{}}dX^{\mu _1^{}}\mathrm{}dX^{\mu _{p+2}^{}}$$
(17)
($`\mu ^{}=0,\mathrm{},p+2`$). In order to get a rank $`(p+2)`$ field strength, we contract $`K_e`$ with a vector field $`V`$, with components $`V=V^M(\frac{}{X^M}`$), which so far remains arbitrary. (There is a sign ambiguity in this contraction; we chose to make it on the left, i.e. $`K(V)_{\mu _1^{}\mathrm{}\mu _{p+2}^{}}V^{\mu _0^{}}K_{\mu _0^{}\mathrm{}\mu _{p+2}^{}}`$). Such a contraction yields
$$K_e(V)=\frac{p+1}{(p+2)!}ϵ_{\mu _0^{}\mathrm{}\mu _{p+2}^{}}V^{\mu _0^{}}dX^{\mu _1^{}}\mathrm{}dX^{\mu _{p+2}^{}}.$$
(18)
Then we reduce the resulting $`(p+2)`$-form to the $`D`$ dimensional hypersurface by using the embedding functions (7),
$`K_e(V)|_\mathrm{\Sigma }={\displaystyle \frac{p+1}{2}}A^{}A^{p+1}drdtdx_1\mathrm{}dx_p\times `$
$`\times \left[2V^\mu x_\mu +V^+({\displaystyle \frac{F^{}}{A^{}}}x^\mu x_\mu )V^{}\right],`$ (19)
where we defined $`V^\pm V^{p+2}\pm V^{p+1}`$. Next we impose that $`K_e(V)|_\mathrm{\Sigma }=F`$. From this we can determine $`V^M`$, using the ansatz $`V^\mu ^{}=\alpha (r)X^\mu ^{}`$. Because $`K_e`$ only has components in the longitudinal directions, $`V^\alpha `$ stays undetermined. When the field strength also includes a magnetic part, this $`V^\alpha `$ comes into play, as we will see in the next subsection. It follows that, in order for (19) to match with (16), $`\alpha (r)`$ has to obey
$$\alpha (r)(\frac{AF^{}}{A^{}}F)=\frac{2}{A}$$
(20)
which gives, using (6)
$$\alpha (r)=\frac{2}{AF+w^2C^2(2C^\kappa 1)}.$$
(21)
or, in terms of the embedding coordinates ($`\alpha (r)\alpha (r(X))\alpha (X)`$)
$$\alpha (X)=\frac{2}{w^2(X_\alpha )^2\left[2(X_\alpha )^\kappa 1+w^2\right](X_M)^2}$$
(22)
and $`V^\mu ^{}(X)=\alpha (X)X^\mu ^{}`$.
### 2.2 D3-brane embedding
Let us now discuss, as an example, how this construction works for the D3-brane in the background produced by other D3-branes. We refer to for a discussion of the embedding of the M2- and M5-brane.
The 10-dimensional Wess–Zumino term is the integral of the self-dual field strength $`F`$ that couples to the D3-branes solution of the type IIB supergravity theory. For the 12-dimensional theory we construct a self-dual 6 form $`K`$, i.e.
$$KK=\eta _{12}|K|^2,$$
(23)
where $`\eta _{12}`$ is the volume form on $`\mathrm{IE}^{(D,2)}`$. Our aim is to obtain $`F`$ as a restriction of $`K`$ to the 10 dimensional surface $`\mathrm{\Sigma }`$. The D3-brane is described by the fields (15) with $`p+1=\kappa =4`$, $`D=10`$. We can therefore write the self-dual field strength in terms of the embedding functions as
$$F=4A^{}A^3dtdxdydzdr+4\omega _{(5)},$$
(24)
where $`\omega _{(5)}\mathrm{sin}(\theta )^4\mathrm{sin}(\varphi _1)^3\mathrm{sin}(\varphi _2)^2\mathrm{sin}(\varphi _3)d\theta d\varphi _1\mathrm{}d\varphi _4`$ is the volume form on the unit 5-sphere.
To find the embedding, we again start by considering a constant form in the embedding space, which in this case we take to be a self-dual six-form
$`K`$ $`=`$ $`{\displaystyle \frac{4}{6!}}(ϵ_{\mu _0^{}\mathrm{}\mu _5^{}}dX^{\mu _0^{}}dX^{\mu _1^{}}\mathrm{}dX^{\mu _5^{}}+`$ (25)
$`+`$ $`ϵ_{\alpha _1\mathrm{}\alpha _6}dX^{\alpha _1}dX^{\alpha _1}\mathrm{}dX^{\alpha _6}),`$
In order to get a rank 5 field strength, we contract, as we have done for the general electric case, $`K`$ with a vector field $`V`$. Such a contraction yields
$`K(V)`$ $`=`$ $`{\displaystyle \frac{4}{5!}}(ϵ_{\mu _0^{}\mathrm{}\mu _5^{}}V^{\mu _0^{}}dX^{\mu _1^{}}\mathrm{}dX^{\mu _5^{}}+`$ (26)
$`+`$ $`ϵ_{\alpha _1\mathrm{}\alpha _6}V^{\alpha _1}dX^{\alpha _2}\mathrm{}dX^{\alpha _6}).`$
Again we reduce the resulting 5 form to the 10-dimensional hypersurface by using the embedding functions (7). By requiring the matching $`K(V)|_\mathrm{\Sigma }=F`$, we get the constraints on our vector field $`V`$. The resulting 5 form $`K(V)|_\mathrm{\Sigma }`$ is precisely the Wess–Zumino term $`\mathrm{\Omega }_5`$ we were looking for.
Let us analyse separately the two terms in the right hand side of (24), (25) and (26). The electric part has already been studied in the general case in the previous subsection. In this case it gives $`V^\mu ^{}=\alpha (r)X^\mu ^{}`$ with
$$\alpha (X)=\frac{2}{2(X^\alpha X_\alpha )^3X^MX_M}.$$
(27)
The magnetic part in (25) can be rewritten in terms of the radial coordinate $`r`$ and the angular coordinates $`\theta `$,$`\varphi _i`$ ($`i=1,\mathrm{},4`$)
$`{\displaystyle \frac{1}{6!}}ϵ_{\alpha _1\mathrm{}\alpha _6}dX^{\alpha _1}\mathrm{}dX^{\alpha _6}=`$
$`=C^{}C^5dr\omega _5.`$ (28)
For the second term in (26) to match with the second term in (24), we have to require that the vector $`V^\alpha `$ points in the radial direction when decomposed in the $`r,\theta ,\varphi _i`$ basis, that is
$$V^\alpha \frac{}{X^\alpha }V^\alpha \frac{r}{X^\alpha }\frac{}{r}.$$
(29)
This gives
$`\frac{1}{5!}ϵ_{\alpha _1\mathrm{}\alpha _6}V^{\alpha _1}dX^{\alpha _2}\mathrm{}dX^{\alpha _6}=`$
$`=C^{}C^5V^\alpha {\displaystyle \frac{r}{X^\alpha }}\omega _5.`$ (30)
Matching this with (24) requires
$$V^\alpha \frac{r}{X^\alpha }=(C^{}C^5)^1,$$
(31)
which, using the ansatz $`V^\alpha =ϵ(r)X^\alpha `$, is solved by<sup>3</sup><sup>3</sup>3 We used the relation $`C^2(r)=X^\alpha X_\alpha `$, from which $`\frac{r}{X^\alpha }=\frac{r}{C^2(r)}\frac{C^2}{X^\alpha }=(CC^{})^1X^\alpha `$.
$$V^\alpha =C^6X^\alpha =(1+r^4)^{\frac{3}{2}}X^\alpha .$$
(32)
We notice that $`ϵ(r0)=\alpha (r0)1`$, so that in the near-horizon approximation we have $`V^M=X^M`$ as was already found in .
The general form of the vector field in terms of the 12-dimensional coordinates is
$$V^\mu ^{}=\frac{2X^\mu ^{}}{2(X^\alpha X_\alpha )^3X^MX_M},V^\alpha =\frac{X^\alpha }{(X^\beta X_\beta )^3}.$$
(33)
## 3 Discussion
The aim of this talk has been to report on a global description of non-dilatonic branes by isometrically embedding them in flat space with two extra dimensions and two times, thus extending the ideas of . We have gained a rather clear global picture of the geometry, giving insight in the structure around coordinate singularities and in the symmetries. In particular, the differences between $`p`$-branes with $`p`$ even and $`p`$ odd, previously pointed out in , are clearly apparent. Like the familiar embedding of anti-de Sitter spacetime as a quadric, our embeddings are periodic in time. This is consistent with some suggestions in , but one may of course always pass to the covering space.
In the context of supergravity and string theory, $`p`$-branes are coupled to $`(p+2)`$-form field strengths. An embedding of the brane thus has to include, besides the embedding of the geometry, a prescription for the forms in the higher dimensional space. This is obtained by defining constant $`(p+3)`$-forms in $`D+2`$ dimensions, and contracting them using a vector $`V`$. The form of $`V`$ is determined by matching the projection on the surface with the known forms for the branes field strengths.
Unfortunately, the geometric significance of the vector field $`V`$ remains unclear. In the case of the M2-brane it is not even unique, since the $`V^\alpha `$ components are arbitrary. A co-dimension 2 surface has a 2-dimensional normal plane. In the D3 and M5 cases, the vector $`V`$ does not lie in this 2-plane, except in the near-horizon limit. Specifically, the normal 2-plane is spanned by $`_\mu \varphi _1`$ and $`_\mu \varphi _2`$. One may check that $`V`$ is not a linear combination of $`_\mu \varphi _1`$ and $`_\mu \varphi _2`$. The bosonic action for probe branes in the embedded background (14) is completely determined after the construction of $`V`$. It could be interesting to investigate if the vector $`V`$ can have some role in the context of F-theory .
Finally it is possible that the methods developed in this paper may be applicable to scenarios in which one regards the universe as a brane embedded in a higher-dimensional spacetime.
## Acknowledgments
This work is supported by the European Commission TMR programme ERBFMRX - CT96 - 0045. C.H is funded by FCT (Portugal) through grant no. PRAXIS XXI/BD/13384/97. |
warning/0003/nlin0003006.html | ar5iv | text | # Dynamics of lattice kinks
## 1 Introduction
Coherent structures, e.g. kinks, solitary waves, vortices, play a central role, as carriers of energy, in many physical systems. An understanding of their dynamical properties, e.g. stability, instability, metastability, is an important problem. While for many years Hamiltonian partial differential equations and their coherent structures, defined on a spatial continuum, have received a great deal of attention, there has been increasing interest in spatially discrete systems. Two important reasons are that (a) certain phenomenoma are intrinsically associated with discreteness and (b) numerical approximations of continuum systems involve the introduction of discreteness, which may lead to spurious numerical phenomena which require recognition. Some examples of the use of discrete systems in the modeling of physical phenomena are: the problem of dislocations propagating on a lattice (for which the Frenkel - Kontorova or discrete sine-Gordon model was originally proposed) , arrays of coupled Josephson junctions, , the problem of the local denaturation of the DNA double helix and coupled optical waveguide arrays .
A result of the many investigations of discrete systems over the last ten to fifteen years has been a recognition, mainly through numerical experiments and heuristic arguments, of the often sharp contrast in behavior between the dynamics of discrete systems and their continuum analogues. Some of these contrasts are easy to anticipate. For example, continuum systems modeling phenomena in a homogeneous environment are translation invariant in space, and may have further symmetry, e.g. Galilean or Lorentz invariance that enable one to construct traveling solutions from static solutions. The analogous discrete system is expected to lose these symmetries and therefore the existence of traveling wave solutions is now brought into question. A natural question concerns the propagation of energy in the lattice; if one initializes the system with data which, for the continuum model would result in a coherent structure propagating through the continuum, what is the corresponding behavior for the energy distribution on intermediate and long times scales on the lattice? Does the system “find” a coherent structure to carry the energy? Does the energy get trapped or pinned? In this paper we seek to obtain insight into these questions for a class of discrete nonlinear wave equations. The special cases we consider in detail are the discrete sine-Gordon (SG) and discrete $`\varphi ^4`$ equations. The methods however are rather general and apply to systems which can be viewed as the interaction between a finite dimensional system of “oscillators” with an infinite-dimensional system governing a continuous spectrum of waves; see the further discussion below and in section 7.
The basic characteristics of the dynamics of coherent structures in discrete systems were systematically explored in a pioneering paper by Peyrard and Kruskal ; see also the contemporaneous papers . Of particular relevance to our work are the more recent articles of Boesch, Willis and coworkers . In this work the discrete sine-Gordon equation is solved with kink-like initial data on the lattice. Observed is a rapid initial velocity adjustment of the kink, a quasi-steady state phase, resonances of the kink oscillations with phonons (continuous spectral modes) and the emission of radiation, which results in the kink’s deceleration and eventual pinning in the Peierls-Nabarro (PN) potential. The final stage involves the relaxation to an asymptotic state, which is either a static or time-periodically “dressed” kink. In a rough sense, this is a result of the absence of a smooth family of traveling kink solutions due to broken Lorentz invariance. Insight into this process can be gleaned by studying the linearized spectrum about the kink.
An important feature of the discrete systems that makes them very different from their continuum analogues is the presence of additional neutral oscillatory modes in the spectrum of the linearization about the coherent structures. The sources of these internal modes are principally of two types; see for example . (1) The continuum system is translation invariant, a symmetry which leads to the continuum linearization having zero modes, eigenmodes and generalized eigenmodes corresponding to zero frequency. Viewed as a perturbation of the continuum system, the discrete, but “nearly” continuum problem is expected to have nearby modes to which these zero modes have been perturbed. If the coherent state is stable, these modes should be neutrally stable, i.e. that is, they must correspond to purely imaginary eigenvalues. (2) It is possible that discrete neutral modes may emerge from the continuous (phonon) spectrum.
Although some of these phenomena have been identified in the early numerical investigations , there has not been a systematic dynamical systems study relating particular resonances to the rate with which the kink is trapped by a “valley” in the PN potential or, once inside such a valley, the rate with which the kink relaxes to its asymptotic state. The main results in this direction date from the work of Ishimori and Munakata , in which the McLaughlin-Scott direct perturbation scheme is used, valid in only in the nearly continuum regime, and the work of Boesch, Willis and El-Batanouny , which is based on numerical simulations and heuristic arguments.
In this paper, we introduce a systematic approach to the study of these phenomena. The work is based on the recent work of one of us (MIW) with A. Soffer on a time dependent theory of metastable states in the context of: quantum resonances , ionization type problems (parametrically excited Hamiltonians) and resonance and radiation damping of bound states in nonlinear wave equations ; see also A fruitful point of view adopted in these works and in the present work is that the dynamics can be understood as the interaction between a finite dimensional dynamical system, governing the bound state (internal modes plus kink) part of the solution and an infinite dimensional dynamical system and governing radiative behavior. Using the tools of scattering theory and the idea of normal forms, we derive a dispersive normal form, which is a closed (up to controllable error terms) finite dimensional system governing the internal mode components of the perturbation. From this normal form many aspects of the the large time asymptotic state are deducible.
The discrete systems we study (e.g. discrete sine-Gordon and discrete $`\varphi ^4`$) depend on a discreteness parameter, $`d`$; as $`d`$ increases the continuum limit is approached. The character of the normal form (topological character of its phase portrait) changes with $`d`$ because the kinds of resonances which occur among discrete mode oscillations and continuum radiation (classifiable in terms of integer linear combinations of internal mode frequencies) change with $`d`$.
Our analytical and numerical studies lead us to the following picture concerning the dynamics in a neighborhood of the kink:
(1) The ground state kink, $`K_{gs}`$, is always Lyapunov stable <sup>1</sup><sup>1</sup>1 If the initial data is in a small neighborhood of $`K_{gs}`$, then the solution remains in a small neighborhood of $`K_{gs}`$ for all time. This notion of stability does not however imply convergence to $`K_{gs}`$.
(2) If $`d`$, the discretization parameter, is sufficiently large (approaching the spatial continuum case), the ground state kink is an attractor, i.e. is asymptotically stable. <sup>2</sup><sup>2</sup>2 By asymptotic stability we mean that a small perturbation of the kink gives rise to a solution which converges as $`t\pm \mathrm{}`$, in some physically relevant norm, to a kink. Asymptotic stability is a notion of stability commonly associated with dissipative dynamical systems and is not commonly associated with general energy conserving systems. In this work we are studying infinite dimensional Hamiltonian systems on infinite spatial domains. These systems have the possibility of radiation of energy to infinity, while keeping the total energy of the system preserved. Thus dissipative behavior is realized through dispersion and eventual radiation of energy out of any compact set; see, for example, the short overview in and references therein. In this case the kink is approached at different algebraic rates depending on the range of $`d`$ values. We infer this from the normal form analysis of section 4.
(3) For $`d`$ sufficiently small, our discrete nonlinear wave models have one or more branches of finite energy time-periodic solutions which bifurcate from $`K_{gs}`$; see Theorem 5.1 and Corollary 5.1. For the specific cases of discrete SG and discrete $`\varphi ^4`$ our results imply, for various regimes of $`d`$, that there exist wobbling kinks, $`W`$. These are time periodic solutions with the same spatial symmetry as the kink and have been previously observed in numerical simulations. Additionally, our results imply in certain regimes of $`d`$, the existence of time periodic solutions (such as $`gW`$, $`eW`$ in the $`\varphi ^4`$ and $`gW`$ in the SG), $`u(t)`$, with the property that $`u(t)K_{gs}`$ is, to leading order in the direction of an even internal mode. Tables 3 and 7 in section 4 indicate the regimes in which these various solutions denoted occur. The wobbling solutions ($`W,gW,eW`$) may be viewed as a class of discrete breather / topological defect states. They are spatially localized and time-periodic oscillations mounted on a static kink background. The usual discrete breathers are mounted on a zero background. This is analogous to the situation with solitons. Dark solitons of the defocusing nonlinear Schrödinger (NLS) equation are mounted on a non-zero continuous wave background, while standard solitons of focusing NLS sit on a zero background.
(4) Numerical simulations and the normal form analysis indicate that in various $`d`$ regimes these time periodic states are local attractors for the dynamics.
(5) The analysis of section 5 implies some nonexistence results concerning solutions of discrete nonlinear wave equations in a neighborhood of the kink. We have that in certain regimes of the discreteness parameter, $`d`$, no nontrivial time-periodic solution exists, and that time quasiperiodic solutions do not exist in any regime of $`d`$. On the other hand, the normal form analysis and numerical simulations indicate, in some regimes of $`d`$, that periodic or quasiperiodic oscillations can be very long lived; see section 6, and in particular, the discussion of Regime VI in section 6.1. Thus, we may think of the system as possessing metastable periodic and quasi-periodic solutions. In regimes of $`d`$ where we show that the wobbling kink, $`W`$, is unstable on long time scales due to a resonance of the “shape mode” with continuum radiation modes, we have the analogue of the the wobbling solution of continuum $`\varphi ^4`$, proved to be stable on large but finite time scales by Segur . On an infinite time scale this wobbling solution behaves as those of the discrete system for large enough $`d`$; eventually the oscillations damp at an algebraic rate leaving a kink in the limit . For the continuum system, the limiting kink may be traveling, while for the discrete system it is a static ground state (centered between lattice sites) kink.
In (2) the dispersive normal form has a dissipative character; the effect of coupling the internal mode oscillations to radiation is modeled by an appropriate nonlinear friction; see section 4 and the discussion in the introduction to on another related model. The damping or friction coefficients are given by formulae which can be understood as a nonlinear generalization of Fermi’s golden rule, arising in the context of the theory of spontaneous emission in atomic physics. Finally, it is worth noting that we obtain a dissipative and therefore apparently time-irreversible normal form from a system which is conservative and time-reversible. There is no contradiction because the dissipation is of an internal nature; it signifies the transfer of energy from the discrete internal mode oscillations to the continuum dispersive waves which propagate to infinity. That dissipative dynamical systems emerge from conservative systems which are a coupling of a low dimensional dynamics to infinite dimensional dynamics (“masses and springs coupled to strings”) has been observed in many contexts; see, for example, and the discussion in .
The paper is organized as follows:
* Section 2 begins with a general discussion of discrete nonlinear wave equations which support kink-like (heteroclinic) structures and then specializes to a discussion of the discrete sine-Gordon (SG) and discrete $`\varphi ^4`$ systems.
* In section 3, we present the decomposition of solutions into discrete internal mode components and radiative components, enabling us to view the dynamics near a kink as the interaction of finite and infinite dimensional Hamiltonian systems.
* In section 4 the dispersive normal forms are derived and discussed for both discrete SG and discrete $`\varphi ^4`$.
* In section 5, we prove that if $`\mathrm{\Omega }`$ is an internal mode frequency and no multiple of it lies in the continuous spectrum of the kink (phonon band), then there is a family of finite energy time periodic solutions which bifurcate from the kink in the “direction” of the corresponding internal mode. The proof is based on the Poincaré continuation method and is an application of the implicit function theorem in an appropriate Banach space. This approach was used also by MacKay and Aubry , who constructed discrete breathers of nonlinear wave equations in the anti-integrable limit.
* In section 6 we combine the normal form analysis of section 4 and existence theory for periodic solutions of section 5 with observations based on numerical simulations to obtain a detailed picture of the dynamics in a neighborhood of the ground state kink.
* Finally, in section 7, we summarize our results and give directions of interest for future research.
### 1.1 Notation
* $`\mathrm{ZZ}`$ denotes the set of all integers and $`\mathrm{IR}`$ denotes the set of real numbers.
* For $`u=\{u_i\}_{i\mathrm{ZZ}}`$, $`\delta _h^2u`$ denotes the discrete Laplacian of $`u`$ defined by:
$$(\delta _h^2u)_i=h^2\left(u_{i+1}2u_i+u_{i1}\right).$$
In the case $`h=1`$ we shall use the simplified notation $`\delta ^2=\delta _1^2`$.
* The inner product of vectors $`u`$ and $`v`$ in $`l^2(\mathrm{ZZ})`$ is given by
$$u,v=\underset{j}{}\overline{u}_jv_j.$$
* The $`t`$ subscript denotes time partial derivatives, e.g. $`u_t(t)=_tu(t)`$, whereas $`i`$ denotes the lattice site numbering.
* $`l^2(\mathrm{ZZ})`$ is the Hilbert space of sequences $`\{u_i\}_{i\mathrm{ZZ}}`$ which are square summable.
* $`H^s(I)`$ denotes the Hilbert space of functions $`f`$, defined on $`I\mathrm{IR}`$ such that $`f`$ and all its derivatives of order $`s`$ are square-integrable. For the case of $`2\pi `$ periodic functions, $`I=S_{2\pi }^1`$, an equivalent norm on $`H^s`$ is given by:
$$f_{H^s}^2=\underset{n\mathrm{ZZ}}{}(1+|n|^2)^s|f_n|^2,$$
where $`f_n`$ denotes the $`n^{th}`$ Fourier coefficient of $`f`$:
$$f_n=(2\pi )^{\frac{1}{2}}_0^{2\pi }e^{2\pi int}f(t)𝑑t.$$
* For example, $`H^2(\mathrm{IR};l^2(\mathrm{ZZ}))`$ denotes the space of functions $`f(t,i)`$ which are $`H^2`$, as functions of $`t`$, with values in the space of $`l^2(\mathrm{ZZ})`$ functions.
* $`\sigma (B)`$ and $`\sigma _{cont}(B)`$ denote respectively the spectrum and the continuous (phonon) spectrum of the operator $`B`$.
* $`0<d`$ denotes the discretization parameter; $`d`$ large is the spatial continuum regime, while for $`d`$ small the effects of discreteness are strong.
## 2 Discrete and continuum wave equations
### 2.1 General Background
Two model nonlinear wave equations on lattices that have played a central role in the theory of nonlinear waves are the discrete sine-Gordon equation (SG) and the discrete $`\varphi ^4`$ model ($`\varphi ^4`$). These equations may be viewed as governing the dynamics of a chain of unit mass particles which, in equilibrium, are equally spaced, a unit distance apart. The particles are then subjected to a conservative force derived from a potential. In terms of the displacement $`u_i`$ of the $`i^{th}`$ particle from equilibrium the equations of motion are<sup>3</sup><sup>3</sup>3We use the notational conventions introduced by Peyrard and Kruskal .:
$$u_{i,tt}=(\delta ^2u)_id^2V^{}(u_i).$$
(2.1)
The case of SG and $`\varphi ^4`$ correspond, respectively, the choices of potential:
$`V(u)`$ $`=`$ $`1\mathrm{cos}u,(\mathrm{SG})`$
$`V(u)`$ $`=`$ $`{\displaystyle \frac{1}{4}}(1u^2)^2.(\varphi ^4)`$
More generally, we assume $`V`$ has three continuous derivatives and satisfies certain constraints appearing below, related to the existence of kink solutions or homoclinic orbits.
The parameter $`d`$ is a fixed constant with, $`d^2`$, having the interpretation of the ratio of on- site potential energy to elastic coupling energy.
Relation between discrete and continuum systems: The parameter $`d`$ can also be seen to play the role of the reciprocal of the lattice spacing in passing between continuum and discrete models. Consider the continuum model governing the displacement $`v(T,x)`$:
$$v_{TT}=v_{xx}V^{}(v).$$
(2.2)
Introducing the lattice spacing parameter, $`h`$, and replacing $`u_{xx}(x)`$ by $`(\delta ^2u)_i`$ we obtain the discrete nonlinear wave equation governing $`v_i=v(T,i;h)`$:
$$v_{i,TT}=(\delta _h^2v)_iV^{}(v_i).$$
The form of the discrete nonlinear wave equation we use is obtained by setting $`t=h^1T`$ and defining the time-scaled displacement $`u(t,i;h)=v(T,i;h)`$. Then, $`u(t,i;h)`$ satisfies (2.1) with $`d=h^1`$.
Taking the formal limit $`d\mathrm{}`$ or equivalently $`h0`$ we have that
$$U(T,i)\underset{d\mathrm{}}{lim}u(dT,i;d^1)=\underset{d\mathrm{}}{lim}v(T,i;d^1)$$
satisfies the associated continuum nonlinear wave equation (2.2).
Hamiltonian structure: The discrete and continuum equations we consider are infinite dimensional Hamiltonian systems with Hamiltonian energy functionals<sup>4</sup><sup>4</sup>4 There is freedom in the choice of potential $`V`$; we choose $`V`$ so that the Hamiltonian energy of the static kink is finite.:
$$=\underset{i}{}\frac{1}{2}u_{i,t}^2+\frac{1}{2}\left(u_{i+1}u_i\right)^2+d^2V(u_i)$$
(2.3)
in the discrete case and
$$=_{\mathrm{IR}}\frac{1}{2}(u_T^2+u_x^2)+V(u)dx.$$
(2.4)
Static kink solutions: Consider time-independent or static solutions of these dynamical systems. In each case, such solutions satisfy the equation obtained by setting time derivatives equal to zero. Thus, in the discrete case we have:
$$(\delta ^2K)_i=V^{}(K_i),$$
(2.5)
and in the continuum case
$$K^{\prime \prime }(x)=V^{}(K(x))$$
(2.6)
Spatially uniform solutions occur at the critical points of the potential, $`V`$. Of particular interest are the static kink solutions. These are static solutions which are heteroclinic. That is, they are static solutions which can be viewed as connections, as $`i\pm \mathrm{}`$ (respectively $`x\pm \mathrm{}`$) in the phase space of two distinct ”unstable” equilibria, values $`K_{}^\pm `$ for which
$`V(K_{}^+)`$ $`=`$ $`V(K_{}^{}),V^{}(K_{}^\pm )=0`$
$`V_{}^{\prime \prime }`$ $``$ $`V^{\prime \prime }(K_{}^+)=V^{\prime \prime }(K_{}^{})>0.`$ (2.7)
Kink solutions can also be constructed by variational methods. Consider the Hamiltonian energy functional, $``$ restricted to $`t`$-independent functions:
$`h[u]`$ $`=`$ $`{\displaystyle \underset{i}{}}\left({\displaystyle \frac{1}{2}}(u_{i+1}u_i)^2+d^2V(u_i)\right),\mathrm{discrete}\mathrm{case}`$ (2.8)
$`h[u]`$ $`=`$ $`{\displaystyle _{\mathrm{IR}}}{\displaystyle \frac{1}{2}}u_x^2+V(u)dx\mathrm{continuum}\mathrm{case}.`$ (2.9)
Then, the Euler Lagrange equations associated with $`h[u]`$ is the equation for the kink, (2.5). That is, if $`K_{gs}`$ denotes a minimizer of $`h[u]`$ then since for all $`\psi l^2(\mathrm{ZZ})`$,
$$\frac{d}{d\tau }h[K_{gs}+\tau \psi ]|{}_{\tau =0}{}^{}=0.$$
This implies that for all $`\psi l^2(\mathrm{ZZ})`$
$$\underset{i}{}\left[(\delta ^2K_{gs})_i+d^2V^{}(K_{gs,i})\right]\psi _i=0.$$
This is equivalent to (2.5). A solution constructed by minimization of $`h[u]`$ is called a ground state kink.<sup>5</sup><sup>5</sup>5 A proof that the minimum of the functional (2.8) is attained can be given using the following strategy, used in a similar problem : For positive integers, $`N`$, define $`h_N[u]`$ to be the truncated Hamiltonian, where the summation is taken over $`NnN`$. Consider $`h_N`$, restricted to vectors satisfying $`u_{\pm N}=K_{}^\pm `$. For any admissible vector, $`u=\{u_i\}_{|i|N}`$, it is possible, by rearrangement, to replace it with another, $`v`$, which is monotonically increasing and for which $`h[v]h[u]`$. Therefore the minimizer of $`h_N[]`$, $`K_{gs}^{(N)}`$, is monotonically increasing from $`K_{}^{}`$ to $`K_{}^+`$. Á priori estimates, derived from the Euler-Lagrange equation (2.5), and the boundary conditions enable one to pass to the limit and obtain a minimizer of $`h[]`$, $`K_{gs}=\{K_{gs,i}\}_{i\mathrm{ZZ}}`$, a monotonically increasing static kink.
Invariance and broken invariance:
An important structural difference between the discrete and continuum cases is that the continuum case has greater symmetry. In particular, (2.2) has the property of translation and Lorentz invariance, while (2.1) does not. A consequence of this is that static solutions of (2.2) can be translated:
$$K(x)K(xx_0)$$
to give a recentered kink and Lorentz boosted:
$$K(x)K((xvt)/\sqrt{1v^2}),|v|<1,$$
to give a traveling solution of velocity $`v`$. These symmetries are absent in the discrete models. Discrete models do however have a discrete translation symmetry but in the models considered this does not give rise to discrete traveling waves.
Dynamic stability of kinks: The issue of dynamic stability is a subtle one. First, there is the question of what is the object we expect to be stable. Also, since our systems are infinite dimensional and not all norms are equivalent, different choices of norms will measure different phenomena. Here we briefly discuss three related notions of stability we have in mind: (a) orbital Lyapunov stability, (b) asymptotic stability and (c) linear spectral stability.
(a) Orbital Lyapunov stability means stability of the shape of the kink; if at $`t=0`$ the data is nearly shaped like a kink then it remains shaped like a kink for all $`t0`$. The kink-like part of the solution will typically move so the solution at different times is near some time-dependent spatial translate of the static kink (element of the symmetry group orbit of the static kink).
In the continuum case, orbital Lyapunov stability in the space $`H^1(\mathrm{IR})`$ is a consequence of characterization of the kink as a local minimizer of $``$. A detailed proof is presented in . A simple proof of Lyapunov stability of ground state discrete kinks can be given based on the same ideas. <sup>6</sup><sup>6</sup>6The proof is based on the following idea. Since $`K_{gs}`$ is an energy minimizer, by (2.12), the second variation at the kink, $`B^2`$ is non-negative. In fact, it is strictly positive; that is, for all $`\psi l^2(\mathrm{ZZ})`$, $`\psi ,B^2\psi \omega _g^2\psi ,\psi `$, where $`\omega _g^2=\omega _g^2(d)>0`$ is the discrete eigenvalue to which the zero mode of the continuum equation perturbs due to discretization of space. Suppose we have initial data near a kink: $`u(0)=K_{gs}+v_0,_tu(0)=v_1`$ and let $`u(t)=K_{gs}+v(t)`$ be the resulting solution. Assume that $`v_0,v_1_{l^2}\epsilon `$, where $`\epsilon `$ is sufficiently small. Then, since $`K_{gs}`$ is a critical point of $`h[]`$ we have: $`\epsilon h[K_{gs}+v(t)]|_{t=0}h[K_{gs}]`$ $`=`$ $`h[K_{gs}+v(t)]h[K_{gs}]={\displaystyle \frac{1}{2}}v(t),B^2v(t)+𝒪(v(t)_{l^2}^3)`$ $`{\displaystyle \frac{\omega _g^2}{2}}v(t)_{l^2}^2cv(t)_{l^2}^3.`$ implying that $`v(t)_{l^2}𝒪(\epsilon )`$ for all $`t0`$, and $`K_{gs}`$ is stable.
(b) Orbital asymptotic stability means that if the initial data is nearly a kink then the solution converges to kink (possibly traveling, in the continuum case or pinned in the discrete case) as $`t\mathrm{}`$.
For both notions (a) and (b) we see that the stable object is the family of solitary wave solutions. That is, in the case where there is a multiparameter family of solutions (kinks, solitary waves …), it is the collection of all such that is the stable object. Concerning the terminology orbital stability, this family of solutions is related to the set of functions generated by applying elements of the equation’s symmetry group to the solution (here, static kink), and thus we can think of the group orbit as stable.
(c) Spectral stability:
As is well known, dynamic stability of the particular solution is related to the spectrum of the linearization of the dynamical system about this solution. Linearization about the kink gives an evolution equation for infinitesimal perturbations of the following form:
$$\left(_t^2+B^2\right)p=0.$$
For the discrete case:
$$B^2p_i=(\delta ^2p)_i+d^2V^{\prime \prime }(K_i)p_i,$$
and for the continuum case:
$$B^2p=p_{xx}+V^{\prime \prime }(K(x))p.$$
At a minimum we expect that stability can hold only if no solutions of the linearized evolution equation, corresponding to finite energy initial conditions, can grow as $`|t|`$ increases. In carrying out linear stability analysis we seek solutions of the form $`p_i=e^{\lambda t}P_i`$ (respectively, $`p=e^{\lambda t}P(x)`$) and obtain linear eigenvalue problems of the form:
$$\left(B^2+\lambda ^2\right)P=0.$$
(2.10)
Explicitly, we have:
$`\lambda ^2P_i`$ $`=`$ $`(\delta ^2P)_i+d^2V^{\prime \prime }(K_i)P_i,\mathrm{discrete}\mathrm{case}`$
$`\lambda ^2P`$ $`=`$ $`P_{xx}+V^{\prime \prime }(K(x))P.\mathrm{continuum}\mathrm{case}`$
We say $`\lambda `$ is in the $`l^2`$ spectrum (respectively, $`L^2`$ spectrum) of the linearization about the kink, or simply spectrum of the kink if the operator $`B^2+\lambda ^2`$ does not have a bounded inverse on $`l^2`$ (respectively $`L^2`$). The spectrum will typically consist of two parts: (i) the point spectrum consisting of isolated eigenvalues of finite multiplicity, for which the corresponding solution of (2.10) is in the Hilbert space, and due to the infinite extent of the spatial domain (ii) continuous spectrum, whose corresponding solutions (radiation or phonon modes) of (2.10) are bounded and oscillatory over the entire spatial domain.
Note that the eigenvalue equation (2.10) has the following symmtery: if $`\lambda `$ has the property that (2.10) has a nontrivial solution solution then $`\lambda ,\overline{\lambda }`$ and $`\overline{\lambda }`$ also have this property. To avoid exponential growing solutions, we must require that the linearized spectrum is a subset of the imaginary axis. If this holds we say the solution is spectrally stable.
Spectrum of the kink:
We are interested in the set of $`\lambda `$’s for which the eigenvalue equation
$$(B^2+\lambda ^2)P=0$$
has a nontrivial solution. We begin by discussing the spectrum of $`B^2`$ and then take square roots to a obtain the spectrum of the the linearization about the kink.
We now make two general remarks about the spectrum of the kink solution, one concerning the continuous spectrum and one concerning the point spectrum.
Continuous spectrum: The continuous spectrum is determined by the solutions of the constant coefficient equation obtained from (2.10) by evaluating its coefficients at spatial infinity. In the discrete case we obtain:
$$\lambda ^2P_i=(\delta ^2P)_i+d^2V_{}^{\prime \prime }P_i,$$
where $`V_{}^{\prime \prime }=lim_{i\pm \mathrm{}}V^{\prime \prime }(K_i)`$. This constant coefficient equation can be solved in terms of exponentials which yield the character of the exact continuum eigensolutions of (2.10). We let $`\lambda =i\omega `$. We find bounded oscillatory solutions of the form: $`P_n=\mathrm{exp}(\sqrt{1}kn),n\mathrm{ZZ}`$, $`k`$ real and arbitrary, where satisfies the dispersion relation:
$$\omega ^2=4\mathrm{sin}^2(k/2)+d^2V_{}^{\prime \prime }.$$
(2.11)
It follows that the continuous spectrum of $`B^2`$ is the positive interval from $`d^2V_{}^{\prime \prime }`$ to $`d^2V_{}^{\prime \prime }+4`$. Therefore, the continuous spectrum of the linearization about the kink solution consists of two intervals on the imaginary axis: $`\pm i[d^1\sqrt{V_{}^{\prime \prime }},\sqrt{4+d^2V_{}^{\prime \prime }}]`$.
Point spectrum: An important conclusion, concerning the point spectrum of kink for the continuum equation, can be made as a consequence of the equation’s symmetry group. If $`K(x)`$ is a kink then translation invariance implies that $`K^{}(x)`$ is a zero mode, a solution of the linear eigenvalue problem with eigenvalue zero, $`B^2K=0`$. This zero mode is often called the Goldstone mode. Since the eigenfunction, $`K^{}`$ does not change sign, zero is the ground state energy (lowest point in the point spectrum) of the Sturm-Liouville operator $`_x^2+V^{\prime \prime }(K(x))`$. This eigenvalue is of multiplicity one.
If one views the discrete model as a perturbation of the continuum model, due to the absence of corresponding symmetries in the discrete problem one expects the eigenvalue at zero to move, as $`d`$ decreases from infinity ($`h`$, the lattice spacing, increases from zero), to the right of zero or to the left of zero. In terms of the spectrum of the kink (plus or minus the square root of the spectrum of $`B^2`$), this means that in the discrete case the spectrum of the kink consists of either (a) a purely imaginary pair, $`\pm i\omega _s`$ (stable case) or (b) a pair of real eigenvalues, symmetrically situated about the origin (unstable case). We shall see both cases in our study of specific discrete models and we shall refer to these eigenvalues and corresponding modes loosely as Goldstone modes as well.
As we shall see in the specific models considered, other eigenvalues may appear in the gap between the upper and lower branches of the continuous spectrum. Such modes of the linearization which lie in this gap are called internal modes, and their number and location can change with $`d`$. For SG we shall find that there are one or two pairs of internal modes, while for the $`\varphi ^4`$ model there are two or three pairs of internal modes.
Finally, we remark that in the case of a ground state kink, one obtained by minimization of $`h[u]`$, the spectrum of $`B^2`$ is non-negative and the kink is spectrally stable. To see this, let $`K_{gs}`$ denote a minimizer of $`h[u]`$, which as indicated above, is a critical point of $`h[]`$. Furthermore, for all $`\psi l^2(\mathrm{ZZ})`$,
$$\frac{d^2}{d\tau ^2}h[K_{gs}+\tau \psi ]|{}_{\tau =0}{}^{}0.$$
A calculation yields that this is equivalent to:
$$\frac{1}{2}\psi ,B^2\psi 0,\mathrm{for}\mathrm{all}\psi l^2(\mathrm{ZZ}).$$
(2.12)
Therefore, $`d^2h[K_{gs}]=\frac{1}{2}B^2`$, the second variation of $`h[]`$ at the ground state kink, is a nonnegative self-adjoint operator. It follows that the spectrum of the kink, $`\pm i\sigma (B)`$, lies on the imaginary axis.
### 2.2 The Sine-Gordon equation
Consider the discrete SG equation:
$`u_{i,tt}=(u_{i+1}+u_{i1}2u_i){\displaystyle \frac{1}{d^2}}\mathrm{sin}u_i.`$ (2.13)
A static discrete $`2\pi `$ kink is a time-independent solution, $`\{K_i\}_{i\mathrm{ZZ}}`$ of (2.13) which satisfies the boundary conditions at infinity:
$$K_i0,\mathrm{as}i\mathrm{},\mathrm{and}K_i2\pi ,\mathrm{as}i\mathrm{}$$
It is well-known that there are two static kink solutions, a high energy one centered on a lattice site and a low energy one centered between two consecutive lattice sites . The low energy kink corresponds to the minimizer of the static Hamiltonian energy, $`h[u]`$, displayed in (2.8). As $`d`$ increases (the continuum limit) the energy difference between the two static kink solutions, the so-called Peierls-Nabarro (PN) barrier, tends to zero and the scaled limit $`u(i;d^1)=K_i`$, with $`i/dx`$ fixed, converges to the continuum SG static kink solution
$$K_{SG}(x)=4\mathrm{tan}^1(e^x).$$
Linear stability analysis about the high energy and low energy kinks was carried out in . From the general discussion of section 2.1, the dispersion relation defining the continuous (phonon) spectrum is $`\omega ^2=d^2+4\mathrm{sin}^2(k/2)`$. Therefore, $`B^2`$ has a band of continuous spectrum of length $`4`$, $`[d^2,4+d^2]`$, and the continuous spectrum of the kink consists of the two bands, $`\pm i[d^1,\sqrt{4+d^2}]`$ .
For the case of the high energy kink, the point spectrum of $`B^2`$ contains a negative eigenvalue derived from the zero (Goldstone) mode associated with the continuum (translation invariant) case. Since the eigenvalue parameter in (2.10) is $`\lambda ^2`$, it follows that the spectrum of the high energy kink (parametrized by $`\lambda `$) consists of two discrete real (Goldstone) eigenvalues, one positive and one negative, $`\pm \omega _{us}`$. These occur at a distance of order $`exp(\pi ^2d)`$ .
The low energy kink, minimizes the Hamiltonian. The second variation of $`h[u]`$ at the kink, $`B^2`$, is therefore a self-adjoint and nonnegative operator. Its spectrum is therefore nonnegative and so the spectrum of the kink is purely imaginary. In this case, the Goldstone mode associated with translation invariance of the continuum problem gives rise to point spectrum consisting of a complex conjugate pair of eigenvalues with order of magnitude $`\mathrm{exp}(\pi ^2d)`$. Furthermore, for $`d>d_e`$, $`d_e0.515`$ a spatially antisymmetric edge mode of energy $`\omega _e^2`$ ($`B^2e=\omega _e^2e`$) appears whose corresponding energy has a distance from the phonon band edge of order $`𝒪(d^4)`$, for $`d`$ large . This mode is not present in the continuum SG.
### 2.3 Discrete $`\varphi ^4`$ Model
Consider the discrete $`\varphi ^4`$ equation:
$`u_{i,tt}=(u_{i+1}+u_{i1}2u_i)+{\displaystyle \frac{1}{d^2}}(u_iu_i^3)`$ (2.14)
In most ways the situation is quite similar to that discussed for the discrete SG. The $`\varphi ^4`$ model has low and high energy kinks. The low energy kink is centered between lattice sites while the high energy kink is centered at a lattice site. As $`d`$ increases the energy difference, the PN barrier, tends to zero, and finally, in the scaled limit $`u(i;d^1)=K_i`$, with $`i/dx`$ fixed, converges ot the continuum $`\varphi ^4`$ kink,
$$K_{\varphi ^4}(x)=\mathrm{tanh}(x/\sqrt{2}).$$
Most qualitative properties of the spectrum of the linearization about the $`\varphi ^4`$ kink are analogous to those in the SG case. The single key difference can be traced to the spectrum associated with the continuum case. In addition to the zero (Goldstone) mode, the spectrum of the operators $`B^2`$ has internal mode of odd parity with eigenfrequency $`\omega _s^2=3/(4d^2)`$. Therefore, the spectrum of the kink has purely imaginary internal modes associated with the energies $`\omega _s=\pm i\sqrt{3}/(2d)`$. These modes are called shape modes.
Turning now to the discrete $`\varphi ^4`$ model we see that $`B^2`$ for the discrete case has the following properties:
* The dispersion relation is $`\omega ^2=2/d^2+4\mathrm{sin}^2(k/2)`$, and therefore $`B^2`$ has continuous spectrum extending from $`2/d^2`$ to $`2/d^2+4`$.
* $`B^2`$ has a positive internal Goldstone mode with energy $`\omega _g^2`$, so that $`B^2g=\omega _g^2g`$. This mode, which is traceable to the zero mode of the continuum model, is spatially even and $`\omega _g`$ is exponentially small in $`d`$ .
* $`B^2`$ has an internal odd shape mode with energy $`\omega _s^2`$, so that $`B^2s=\omega _s^2s`$. This mode, which is traceable to the internal shape mode of the continuum $`\varphi ^4`$ model, is spatially odd.
* For $`d>d_e`$ ($`d_e0.82`$) $`B^2`$ has an internal edge mode which is even and with corresponding energy $`\omega _e^2`$, with $`B^2e=\omega _e^2e`$, which emerges from the continuous spectrum. For large $`d`$ we have
$$\omega _{e}^{}{}_{}{}^{2}2d^2\frac{2}{15^2}d^6.$$
It follows that the spectrum of the discrete $`\varphi ^4`$ kink consists of pair of two or three complex conjugate pairs of internal modes at energies $`\pm i\omega _g`$, $`\pm i\omega _s`$, and for $`d>d_e`$, $`\pm i\omega _e`$.
In the next section we shall see that whether or not certain integer linear combinations of the internal mode frequencies land in the continuous spectrum determines the large time asymptotic behavior.
## 3 Resonances and radiation damping
That we can Lorentz boost a static kink and obtain a traveling kink on the spatial continuum (section 2) motivates the following question:
What happens if we attempt to cause a kink to propagate through the lattice?
The numerical and formal asymptotic investigation was initiated in ,. For example, choose as initial conditions for the discrete sine Gordon system the state $`u_i(t=0)=K(i/\sqrt{1v^2})`$, where $`K(x)`$ is the continuum kink and $`0<|v|<1`$. Dispersive radiation plays an important role in the dynamics. The kink moves approximately as particle under the influence of the sinusoidal (Peierls-Nabarro) potential. As the kink propagates, its velocity alternately increases and decreases, depending on the location of its center of mass relative to the peaks and troughs of the potential. This oscillation leads to a resonance with the continuous spectrum and a transfer of energy from the propagating kink to radiation modes. Eventually, the kink has slowed so that its energy no longer exceeds the PN barrier and then executes a damped oscillation about a fixed lattice site to which it is pinned. It is this latter process that we study here.
The goal of this study is to, by analytic and numerical means,
* identify what the asymptotic state of the system is in different regimes of the discretization parameter $`d`$,
* find the rate of approach to this state.
We begin in the setting of the discrete $`\varphi ^4`$ model, (2.14). This case is a more complicated than discrete SG; discrete $`\varphi ^4`$ has either two or three internal modes, while discrete SG which has at most two. After a detailed treatment of the $`\varphi ^4`$ case, we’ll give an outline of the analogous results for SG.
We begin with a solution which is a small perturbation about a low energy (ground state) kink. This corresponds to solving the initial value problem for (2.14) with initial conditions:
$$u_i(0)=K_i+\epsilon v_{0i},_tu_i(0)=\epsilon v_{1i},$$
where $`\epsilon `$ a small paramater. We then view the solution as
$$u_i(t)K_i+\epsilon v_i(t),$$
where $`v_i(t)`$ denotes the perturbation about the kink. Equation (2.14) implies an evolution equation for $`v_i`$: $`\left(_t^2+B^2\right)v=\mathrm{}`$ and since the data for $`v_i`$ is small, it is natural to expand $`v_i`$ in terms of the modes of the linearization about the kink. Recall that the operator $`B^2`$ is positive and self-adjoint, and has, depending on the range of $`d^{}s`$ considered, two or three internal modes $`g`$, $`s`$, and $`e`$ (if $`d>d_e`$):
$`B^2g_i=\omega _g^2g_i,B^2s_i=\omega _s^2s_i,B^2e_i=\omega _e^2e_i`$
$`0<\omega _g^2<\omega _s^2<\omega _e^2<4+2d^2,`$
where
$$B^2\psi _i\delta ^2\psi _i+d^2(3K_i^21)\psi _i.$$
We assume these internal modes to be normalized so that the set $`\{g,s,e\}`$ is orthonormal in $`l^2(\mathrm{ZZ})`$. We also introduce the operator which projects onto the orthogonal complement of the subspace of internal modes:
$$P_c\psi \psi g,\psi gs,\psi se,\psi e$$
(3.1)
$`P_c`$ is the projection onto the continuous spectral part of $`B^2`$ (phonon or radiation modes) and therefore the solutions of $`(_t^2+B^2)u=0`$ with data in the range of $`P_c`$ are expected to decay dispersively as $`t\pm \mathrm{}`$.
We expand the perturbation $`\epsilon v_i`$ in terms of the internal mode subspace and its orthogonal complement and thus
$`u_i(t)=K_i+\epsilon a(t)g_i+\epsilon b(t)s_i+\epsilon c(t)e_i+\epsilon ^2\eta _i(t),`$ (3.2)
with
$`g,\eta (t)=s,\eta (t)=e,\eta (t)=0,`$
$`P_c\eta =\eta `$
Substitution of (3.2) into (2.14) yields
$`(a_{tt}+\omega _{g}^{}{}_{}{}^{2}a)g_i+(b_{tt}+\omega _{s}^{}{}_{}{}^{2}b)s_i+(c_{tt}+\omega _{e}^{}{}_{}{}^{2}c)e_i+\epsilon _t^2\eta _i`$
$`=\epsilon B^2\eta d^2(3\epsilon K_iM_i^2+\epsilon ^2M_i^3+6\epsilon ^2K_i\eta M_i+3\epsilon ^3M_i^2\eta _i+𝒪(\epsilon ^4))`$ (3.3)
where:
$$M_ia(t)g_i+b(t)s_i+c(t)e_i.$$
(3.4)
Notice that from here on for simplicity and compactness we drop the spatial indices $`i`$.
We next project (3.3) onto the internal modes $`g,s`$ and $`e`$, as well as on the range of $`P_c`$. This yields the following coupled system of four equations for the three internal mode amplitudes and the continuous spectral part:
$`_t^2a+\omega _{g}^{}{}_{}{}^{2}a`$ $`=`$ $`{\displaystyle \frac{1}{d^2}}\left(3\epsilon g,KM^2+\epsilon ^2g,M^3+6\epsilon ^2g,KM\eta +6\epsilon ^3g,M^2\eta \right)+𝒪(\epsilon ^4)`$
$`_t^2b+\omega _{s}^{}{}_{}{}^{2}b`$ $`=`$ $`{\displaystyle \frac{1}{d^2}}\left(3\epsilon s,KM^2+\epsilon ^2s,M^3+6\epsilon ^2s,KM\eta +6\epsilon ^3s,M^2\eta \right)+𝒪(\epsilon ^4)`$
$`_t^2c+\omega _{e}^{}{}_{}{}^{2}c`$ $`=`$ $`{\displaystyle \frac{1}{d^2}}(3\epsilon e,KM^2+\epsilon ^2e,M^3+6\epsilon ^2e,KM\eta +6\epsilon ^3e,M^2\eta )+𝒪(\epsilon ^4))`$
$`_t^2\eta +B^2\eta `$ $`=`$ $`{\displaystyle \frac{1}{d^2}}P_c\left(\epsilon M^3+3KM^2+𝒪(\epsilon \eta M)+𝒪(\epsilon ^4\eta ^3)\right).`$ (3.8)
The initial data for $`a(t),b(t),c(t)`$ is:
$`a(0)`$ $`=`$ $`g,v_0,_ta(0)=g,v_1`$
$`b(0)`$ $`=`$ $`s,v_0,_tb(0)=s,v_1`$
$`c(0)`$ $`=`$ $`e,v_0,_tb(0)=e,v_1`$
$`\eta (0)`$ $`=`$ $`P_cv_0,_t\eta (0)=P_cv_1`$ (3.9)
For simplicity we shall assume:
$$\eta (0)=0,_t\eta (0)=0$$
(3.10)
corresponding to perturbations of the kink which only excite the internal modes. The general case can be treated as well; see .
The system (LABEL:eq8-3.8) may be viewed as finite dimensional Hamiltonian system governing three oscillators with amplitudes $`a,b`$ and $`c`$ and natural frequencies $`\pm \omega _g,\pm \omega _s,\pm \omega _e`$, coupled by nonlinearity to an infinite dimensional Hamiltonian wave equation for a field $`\eta `$. Systems of this type have been analyzed rigorously and the behavior of their solutions determined on short, intermediate and infinite time scales . The analysis we present uses and extends the methods of these works. Roughly speaking we transform the system (LABEL:eq8)-(3.8) into an equivalent dynamical system, which is a perturbation of a normal form for which one obtains information on: (i) what the asymptotic state of the system is, and (ii) how this asymptotic state is approached.
## 4 The dispersive normal form
Our goal in this section is to derive the normal forms which give detailed information on the dynamics in a neighborhood of the ground state static kink of SG and $`\varphi ^4`$. In particular, the normal forms anticipate the existence of time-periodic solutions. These predictions are upheld by the rigorous results of section 5. In section 6, the normal form analysis is then joined with the existence theory of section 5 and numerical simulations to fill out the picture of what the dynamics are in a neighborhood of the kink.
### 4.1 The normal form for discrete $`\varphi ^4`$
At this stage, we have in (LABEL:eq8-3.8) a formulation of the discrete $`\varphi ^4`$ dynamical as a system governing the discrete internal modes interacting, due to nonlinearity, with a system governing dispersive waves (phonons). Our goal in this section is to present and implement a method, based on , leading to a reformulation of the coupled discrete-continuum mode system (LABEL:eq8-3.8) as a perturbed dispersive normal form in which the nature of the energy transfer among modes is made explicit.
Ideally, one could solve the equation for the dispersive part, $`\eta `$, explicitly in terms of the discrete mode amplitudes: $`a,b`$, and $`c`$ and then substitute the result into the mode amplitude equations to get a closed system of equations for $`a,b`$, and $`c`$. This then could be used in the equation for $`\eta `$ to determine its behavior. Due to the system’s nonlinearity, one can’t expect to solve for the exact behavior of $`\eta `$ in terms of $`a,b`$ and $`c`$. Instead we solve perturbatively in $`\epsilon `$. Let
$$\eta \underset{j=0}{\overset{\mathrm{}}{}}\epsilon ^j\eta ^{(j)}.$$
Then,
$$_t^2\eta ^{(0)}+B^2\eta ^{(0)}=3d^2KM^2.$$
(4.1)
The function $`\eta `$ captures the leading order radiative response, to the internal mode excitations.
Note that $`\eta ^{(0)}`$ is the solution of a forced wave equation. The forcing term in (4.1) involves $`M^2`$, which by (3.4), contains squares and cubes of $`a,b`$ and $`c`$. Also, note that since the nonlinear coupling terms on the right hand sides of equations (LABEL:eq8-LABEL:eq10) are small, $`a(t),b(t)`$ and $`c(t)`$ are slow modulations of the exponentials $`e^{\pm i\omega _gt}`$, $`e^{\pm i\omega _st}`$ and $`e^{\pm i\omega _et}`$. The forcing on the right hand side of (4.1) will be resonant if nonlinearity, when acting on $`a,b`$ and $`c`$, generates frequencies which lie in the spectrum of the operator $`B`$. Since $`B`$ has a band of continuous spectrum, it can easily happen that integer linear combinations of the internal mode frequencies can lie in the continuous spectrum of $`B`$.
Figure 6 in section 4 displays the locations of the frequencies $`\omega _g,\omega _s`$ and $`\omega _e`$, their multiples and certain integer linear combinations relative to the continuous (phonon) spectrum. Figure 7 is the analogous plot for discrete SG.
The key calculation is to compute the effect of such resonances which result in the transfer of energy from the discrete to continuum modes.
Let
$`a(t)`$ $`=`$ $`A(t)\mathrm{exp}(i\omega _gt)+\overline{A}(t)\mathrm{exp}(i\omega _gt),`$ (4.2)
$`b(t)`$ $`=`$ $`B(t)\mathrm{exp}(i\omega _st)+\overline{B}(t)\mathrm{exp}(i\omega _st),`$ (4.3)
$`c(t)`$ $`=`$ $`C(t)\mathrm{exp}(i\omega _et)+\overline{C}(t)\mathrm{exp}(i\omega _et),`$ (4.4)
We further impose the constraints:
$`A_t\mathrm{exp}(i\omega _gt)+\overline{A}_t\mathrm{exp}(i\omega _gt)`$ $`=`$ $`0`$
$`B_t\mathrm{exp}(i\omega _st)+\overline{B}_t\mathrm{exp}(i\omega _st)`$ $`=`$ $`0`$
$`C_t\mathrm{exp}(i\omega _et)+\overline{C}_t\mathrm{exp}(i\omega _et)`$ $`=`$ $`0.`$
and obtain:
$`A_t`$ $`=`$ $`(2i\omega _g)^1e^{i\omega _gt}F_1`$ (4.5)
$`B_t`$ $`=`$ $`(2i\omega _s)^1e^{i\omega _st}F_2`$ (4.6)
$`C_t`$ $`=`$ $`(2i\omega _e)^1e^{i\omega _et}F_3`$ (4.7)
where $`F_i=F_{i1}+F_{i2}`$ (i=1,2,3). The $`3\times 2`$ matrix $`F_{ij}`$ is given by
$`F`$ $`=`$ $`d^2\left[\begin{array}{cc}\hfill 3\epsilon g,KM^2+\epsilon ^2g,M^3& \hfill 6\epsilon ^2g,KM\eta +3\epsilon ^3g,M^2\eta \\ \hfill 3\epsilon s,KM^2+\epsilon ^2s,M^3& \hfill 6\epsilon ^2s,KM\eta +3\epsilon ^3s,M^2\eta \\ \hfill 3\epsilon e,KM^2+\epsilon ^2e,M^3& \hfill 6\epsilon ^2e,KM\eta +3\epsilon ^3e,M^2\eta \end{array}\right]`$ (4.11)
We first solve (4.1) and obtain the expression for $`\eta ^{(0)}`$ in terms of $`a(t),b(t)`$ and $`c(t)`$ or equivalently $`A(t),B(t)`$ and $`C(t)`$
$`\eta ^{(0)}=3d^2{\displaystyle _0^t}{\displaystyle \frac{\mathrm{sin}(B(t\tau ))}{B}}P_cKM^2𝑑\tau `$ (4.13)
Recall that the dependence on the internal mode amplitudes is through $`M`$, defined in (3.4). Substitution of (4.13) into equations (4.5 -4.7) yields a closed system for the internal mode amplitudes through order $`\epsilon ^3`$. The key terms in this equation come from certain resonances and our goal now is to show how they arise.
The first column of terms in $`F`$ involves the interactions among discrete modes. They do not generate frequencies contributing to any resonant forcing. The second column contains terms which couple the discrete bound state part to the continuum radiation modes. In order to identify all the resonances one has to explicitly expand out the $`F_{i2}`$ terms and to use equation (4.13). We do not carry out the detailed computations in all detail here. Rather, we illustrate the key ideas and methodology by considering prototypical terms. A complete and rigorous implementation of these ideas in another nonlinear wave context is presented in .
We focus on the $`𝒪(\epsilon ^2)`$ term in $`F_{12}`$ and some of its contributions to the equation for $`A(t)`$. In particular we shall consider all resonant contributions and a sample nonresonant contribution. By equation (4.5) we must then consider the expression:
$$3i\epsilon ^2d^2\omega _g^1e^{i\omega _gt}Kg,M\eta ^{(0)}.$$
(4.14)
Consider $`\eta ^{(0)}`$.
$`\eta ^{(0)}=3d^2{\displaystyle _0^t}{\displaystyle \frac{\mathrm{sin}(B(t\tau ))}{B}}P_cKM^2𝑑\tau ={\displaystyle \frac{3}{2iB}}d^2{\displaystyle _0^t}e^{iB(t\tau )}P_cKM^2𝑑\tau +\mathrm{}`$ (4.15)
$`=`$ $`{\displaystyle \frac{3i}{2B}}d^2e^{iBt}{\displaystyle _0^t}e^{iB\tau }P_cKg^2a^2(\tau )𝑑\tau +\mathrm{}`$
$`=`$ $`{\displaystyle \frac{3i}{2B}}e^{iBt}{\displaystyle _0^t}e^{iB\tau }\left(A^2(\tau )e^{2i\omega _g\tau }+2|A(\tau )|^2+\overline{A}^2(\tau )e^{2i\omega _g\tau }\right)P_cKg^2𝑑\tau +\mathrm{}`$
$`=`$ $`{\displaystyle \frac{3i}{2B}}e^{iBt}{\displaystyle _0^t}A^2(\tau )e^{i(B2\omega _g)\tau }P_cKg^2𝑑\tau +{\displaystyle \frac{3i}{2B}}e^{iBt}{\displaystyle _0^t}e^{i(B+2\omega _g)\tau }\overline{A}^2(\tau )P_cKg^2𝑑\tau +\mathrm{}`$
$``$ $`\eta _{res}^{(0)}+\eta _{nr}^{(0)}+\mathrm{}`$
We have only kept two terms in the above calculation: the only term leading to a resonant contribution in the $`A`$ equation and one (of many) nonresonant terms.
We wish to expand $`\eta ^{(0)}`$ using the following formula, which follows by straightforward integration by parts:
$`e^{iBt}{\displaystyle _0^t}e^{i(B\zeta )\tau }P_c\alpha (\tau )𝑑\tau `$ (4.16)
$`=`$ $`ie^{iBt}(B\zeta )^1P_ce^{i\zeta t}\alpha (\tau )ie^{iBt}(B\zeta )^1P_c\alpha (0)`$
$`ie^{iBt}{\displaystyle _0^t}e^{i(B\zeta )\tau }(B\zeta )^1P_c_\tau \alpha (\tau )d\tau `$
with $`\alpha (\tau )=A^2(\tau )P_cKg^2`$, for example, and that $`\tau `$-derivatives of $`A`$ are of order $`\epsilon `$. For the term $`\eta _{res}^{(0)}`$, $`\zeta =2\omega _g`$, which for $`d`$ in the range $`[0.54,0.6364]`$ lies in the continuous spectrum; see Table 2 and Figure 6 below. Such resonant terms are of paramount interest and govern energy transfer. To treat these resonant terms, an appropriate modification of (4.16) is required. We use the following:
Let $`\kappa =sgn(t)`$, which is equal to $`+1`$ if $`t>0`$ and $`1`$ for $`t<0`$.
$`e^{iBt}{\displaystyle _0^t}e^{i(B\zeta )\tau }P_c\alpha (\tau )𝑑\tau =`$
$`ie^{iBt}(B\zeta +i\kappa 0)^1P_ce^{i\zeta t}\alpha (\tau )ie^{iBt}(B\zeta +i\kappa 0)^1P_c\alpha (0)`$
$`ie^{iBt}{\displaystyle _0^t}e^{i(B\zeta )\tau }(B\zeta +i\kappa 0)^1P_c_\tau \alpha (\tau )d\tau ,`$ (4.17)
where
$$(B\zeta +i\kappa 0)^1=\underset{ϵ0}{lim}(B\zeta +i\kappa ϵ)^1.$$
###### Remark 4.1
(1) The sense in which the operator expansions (4.16) and (4.17) are correct is in a distributional sense. That is, equality holds when multiplying both sides by a smooth function with spatial support which is compact and integrating both sides over all space.
(2) Formula (4.17) is proved by first writing the integral on the left hand side as
$$_0^t\mathrm{exp}(i(B\zeta +i\kappa ϵ)\tau )P_c\alpha (\tau )𝑑\tau .$$
For any $`ϵ>0`$ (4.16) can be used with $`\zeta `$ replaced by $`\zeta iϵ`$. We then pass to the limit as $`ϵ0`$.
(3) The choice of regularization, $`+i0(\kappa =1)`$ for $`t>0`$ and $`i0(\kappa =1)`$ for $`t<0`$ is connected with the condition that the latter two terms in (4.17) consist of outgoing radiation at spatial infinity, and is therefore time-decaying in an appropriate local energy sense . In this article, we take into account only the first term in (4.17). Reference contains a fully detailed and rigorous treatment in a related context. Henceforth, we shall for simplicity assume $`t>0`$ and therefore work with the $`+i0`$ regularization.
(4) If $`\zeta `$ does not lie in the continuous spectrum the formula (4.17) reduces to (4.16).
We now continue the expansion of $`\eta ^{(0)}`$ using (4.16) to study $`\eta _{nr}^{(0)}`$ and (4.17) to study $`\eta _{res}^{(0)}`$:
$`\eta _{res}^{(0)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}d^2B^1(B2\omega _g+i0)^1A(t)^2P_cKg^2+\mathrm{}`$
$`\eta _{nr}^{(0)}`$ $`=`$ $`{\displaystyle \frac{3}{2}}d^2B^1(B+2\omega _g)^1\overline{A}(t)^2P_cKg^2+\mathrm{}`$
Substitution of $`\eta ^{(0)}=\eta _{res}^{(0)}+\eta _{nr}^{(0)}+\mathrm{}`$ into (4.14) we have:
$`3i\epsilon ^2d^2\omega _g^1e^{i\omega _gt}Kg,M\eta ^{(0)}`$
$`=3i\epsilon ^2d^2\omega _g^1e^{i\omega _gt}\left(Ae^{i\omega _gt}+\overline{A}e^{i\omega _gt}\right)Kg^2,\eta ^{(0)}+\mathrm{}`$
$`={\displaystyle \frac{9}{2}}i\epsilon ^2d^4\omega _g^1|A|^2AKg^2,B^1(B2\omega _g+i0)^1P_cKg^2+\mathrm{}`$
$`{\displaystyle \frac{9}{2}}i\epsilon ^2d^4\omega _g^1\overline{A}^3e^{3i\omega _gt}Kg^2,B^1(B+2\omega _g)^1P_cKg^2+\mathrm{}`$
$`\epsilon ^2\left(\mathrm{\Gamma }_{2\omega _g}+i\mathrm{\Lambda }_{2\omega _g}\right)|A(t)|^2A(t)+\rho (t)\overline{A}^3(t)`$ (4.18)
To calculate $`\mathrm{\Lambda }`$ and $`\mathrm{\Gamma }`$, we apply a generalization to self-adjoint operators of the well-known distributional identity:
$$\underset{ϵ0}{lim}(\xi \pm iϵ)^1=\mathrm{P}.\mathrm{V}.\xi ^1i\pi \delta (\xi ),$$
(4.19)
where $`\delta (\xi )`$ is the Dirac delta mass at $`\xi =0`$ and $`\mathrm{P}.\mathrm{V}.`$ denotes the principal value integral.
Therefore, using (4.19) we obtain:
$`\epsilon ^2\mathrm{\Gamma }_{2\omega _g}`$ $``$ $`\epsilon ^2\mathrm{\Gamma }(2\omega _g)={\displaystyle \frac{9\pi }{4\omega _g^2}}\epsilon ^2d^4Kg^2,\delta (B2\omega _g)Kg^2,`$ (4.20)
$`\epsilon ^2\mathrm{\Lambda }_{2\omega _g}`$ $``$ $`\epsilon ^2\mathrm{\Lambda }(2\omega _g)={\displaystyle \frac{9}{2}}\epsilon ^2d^4Kg^2,P.V.(B2\omega _g)P_cKg^2.`$ (4.21)
###### Remark 4.2
(1) Had we done this calculation for the $`B`$ (respectively, $`C`$) equation, we’d have obtained as a coefficient of $`|B|^2B`$ (respectively, $`|C|^2C`$) the quantity $`\epsilon ^2(\mathrm{\Gamma }_{2\omega _s}+i\mathrm{\Lambda }_{2\omega _s})`$ (respectively, $`\epsilon ^2(\mathrm{\Gamma }_{2\omega _e}+i\mathrm{\Lambda }_{2\omega _e})`$) (2) If $`\zeta \sigma _{cont}(B)`$, then $`\mathrm{\Gamma }_\zeta `$ is always nonnegative and, generically, is strictly positive. It is the analogue of Fermi’s golden rule wwhich arises in the context of the theory of spontaneous emission ,. Apart from a positive constant prefactor, $`\mathrm{\Gamma }_\zeta `$ is the square of the Fourier transform of $`Kg^2`$ relative to the continuous spectral part of $`B`$ evaluated at $`\zeta `$.
The above calculations yield the following information on the structure of the equation for $`A(t)`$:
$$A_t=\epsilon ^2\left(\mathrm{\Gamma }_{2\omega _g}+i\mathrm{\Lambda }_{2\omega _g}\right)|A|^2A+\mathrm{}.$$
(4.22)
Extensive calculations, of which the above are representative, yield a system of the form:
$`A_t`$ $`=`$ $`\epsilon ^2\left(\alpha _1|A|^2+\alpha _2|B|^2+\alpha _3|C|^2\right)A`$
$`\epsilon ^4\left(\alpha _4|A|^4+\alpha _5|A|^2|B|^2+\alpha _6|B|^4+\alpha _7|B|^2|C|^2+\alpha _8|C|^4\right)A`$
$`+`$ $`\epsilon \mathrm{\Xi }_A(t,A,B,C;\epsilon )`$
$`B_t`$ $`=`$ $`\epsilon ^2(\beta _|A|^2+\beta _2|B|^2+\beta _3|C|^2)B`$
$`\epsilon ^4\left(\beta _4|A|^4+\beta _5|A|^2|B|^2+\beta _6|B|^4+\beta _7|B|^2|C|^2+\beta _8|C|^4\right)B`$
$`+`$ $`\epsilon \mathrm{\Xi }_B(t,A,B,C;\epsilon )`$
$`C_t`$ $`=`$ $`\epsilon ^2\left(\gamma _1|A|^2+\gamma _2|B|^2+\gamma _3|C|^2\right)C`$ (4.23)
$`\epsilon ^4\left(\gamma _4|A|^4+\gamma _5|A|^2|B|^2+\gamma _6|B|^4+\gamma _7|B|^2|C|^2+\gamma _8|C|^4\right)C`$
$`+`$ $`\epsilon \mathrm{\Xi }_C(t,A,B,C;\epsilon )`$
The coefficients $`\alpha _j`$, $`\beta _j`$ and $`\gamma _j`$ are, in general, complex numbers. The terms $`\mathrm{\Xi }_A`$, $`\mathrm{\Xi }_A`$ and $`\mathrm{\Xi }_A`$ involve bounded and oscillatory complex exponentials in $`t`$ multiplying monomials in $`A,B,C`$ of cubic or higher degree. Also, in order to obtain the fifth degree terms it is necessary to construct $`\eta `$ through second order in $`\epsilon `$ (and therefore the continuous spectral part of the perturbation about the kink, $`\epsilon ^2\eta `$, through $`\epsilon ^4`$. Coupling to $`\eta `$ is neglected as is the dynamical equation for $`\eta `$. The preceding calculation gives $`\alpha _1=\mathrm{\Gamma }_g+i\mathrm{\Lambda }_g`$ plus a further term contributed by the $`𝒪(\epsilon ^2)`$ entry of $`F_{11}`$ in (LABEL:eq:Fmatrix). An exhaustive tabulation of all coefficients $`\alpha _j,\beta _j,\gamma _j`$ would be, to put it midly, a very lengthy exercise. Since the principal effect we seek to illuminate is that of nonlinear resonant coupling on the internal mode amplitudes $`|A|,|B|`$ and $`|C|`$, we only tabulate those coefficients up to the order considered which may play a role.
For this it is convenient to introduce the notation:
$$G(\zeta )=P_cB^1(B\zeta +i0)^1P_c$$
(4.24)
Table 1A: Principal $`\varphi ^4`$ g-mode coefficients
| Term | Coefficient ($`\alpha _j`$) |
| --- | --- |
| $`\left|A\right|^2A`$ | $`\frac{9}{2}id^4\omega _g^1Kg^2,G\left(2\omega _g\right)Kg^2`$ |
| $`\left|B\right|^2A`$ | $`9id^4\omega _g^1Kgs,G\left(\omega _g+\omega _s\right)Kgs`$ |
| $`\left|C\right|^2A`$ | $`9id^4\omega _g^1Kge,G\left(\omega _g+\omega _e\right)Kge`$ |
| $`\left|A\right|^4A`$ | $`\frac{3}{4}id^4\omega _g^1g^3,G\left(3\omega _g\right)g^3`$ |
| $`\left|B\right|^4A`$ | $`\frac{9}{4}id^4\omega _g^1gs^2,G\left(\omega _g+2\omega _s\right)gs^2`$ |
| $`\left|C\right|^4A`$ | $`\frac{9}{4}id^4\omega _g^1ge^2,G\left(\omega _g+2\omega _e\right)ge^2`$ |
| $`\left|A\right|^2\left|B\right|^2A`$ | $`\frac{9}{2}id^4\omega _g^1g^2s,G\left(2\omega _g+\omega _s\right)g^2s`$ |
| $`\left|A\right|^2\left|C\right|^2A`$ | $`\frac{9}{2}id^4\omega _g^1g^2e,G\left(2\omega _g+\omega _e\right)g^2e`$ |
| $`\left|B\right|^2\left|C\right|^2A`$ | $`9id^4\omega _g^1gse,G\left(\omega _g+\omega _s+\omega _e\right)gse`$ |
Table 1B: Principal $`\varphi ^4`$ s-mode coefficients
| Term | Coefficient ($`\beta _j`$) |
| --- | --- |
| $`\left|B\right|^2B`$ | $`\frac{9}{2}id^4\omega _s^1Ks^2,G\left(2\omega _s\right)Ks^2`$ |
| $`\left|A\right|^2B`$ | $`9id^4\omega _s^1Kgs,G\left(\omega _g+\omega _s\right)Kgs`$ |
| $`\left|C\right|^2B`$ | $`9id^4\omega _s^1Kse,G\left(\omega _s+\omega _e\right)Kse`$ |
| $`\left|B\right|^4B`$ | $`\frac{3}{4}id^4\omega _s^1s^3,G\left(3\omega _s\right)s^3`$ |
| $`\left|A\right|^4B`$ | $`\frac{9}{4}id^4\omega _s^1sg^2,G\left(2\omega _g+\omega _s\right)sg^2`$ |
| $`\left|C\right|^4B`$ | $`\frac{9}{4}id^4\omega _s^1se^2,G\left(\omega _s+2\omega _e\right)se^2`$ |
| $`\left|A\right|^2\left|B\right|^2B`$ | $`\frac{9}{2}id^4\omega _s^1gs^2,G\left(\omega _g+2\omega _s\right)s^2g`$ |
| $`\left|A\right|^2\left|B\right|^2B`$ | $`\frac{9}{2}id^4\omega _s^1gs^2,G\left(2\omega _s\omega _g\right)s^2g`$ |
| $`\left|B\right|^2\left|C\right|^2B`$ | $`\frac{9}{2}id^4\omega _s^1s^2e,G\left(2\omega _s+\omega _e\right)s^2e`$ |
| $`\left|A\right|^2\left|C\right|^2B`$ | $`9id^4\omega _s^1gse,G\left(\omega _s+\omega _e\omega _g\right)gse`$ |
| $`\left|A\right|^2\left|C\right|^2B`$ | $`9id^4\omega _s^1gse,G\left(\omega _g+\omega _s+\omega _e\right)gse`$ |
Table 1C: Principal $`\varphi ^4`$ e-mode coefficients
| Term | Coefficient ($`\gamma _j`$) |
| --- | --- |
| $`\left|C\right|^2C`$ | $`\frac{9}{2}id^4\omega _e^1Ke^2,G\left(2\omega _e\right)Ke^2`$ |
| $`\left|B\right|^2C`$ | $`9id^4\omega _e^1Kse,G\left(\omega _s+\omega _e\right)Kse`$ |
| $`\left|A\right|^2C`$ | $`9id^4\omega _e^1Kge,G\left(\omega _g+\omega _e\right)Kge`$ |
| $`\left|C\right|^4C`$ | $`\frac{3}{4}id^4\omega _e^1e^3,G\left(3\omega _e\right)e^3`$ |
| $`\left|B\right|^4C`$ | $`\frac{9}{4}id^4\omega _e^1s^2e,G\left(2\omega _s+\omega _e\right)s^2e`$ |
| $`\left|A\right|^4C`$ | $`\frac{9}{4}id^4\omega _e^1g^2e,G\left(2\omega _g+\omega _e\right)g^2e`$ |
| $`\left|C\right|^2\left|B\right|^2C`$ | $`\frac{9}{2}id^4\omega _e^1se^2,G\left(\omega _s+2\omega _e\right)se^2`$ |
| $`\left|C\right|^2\left|B\right|^2C`$ | $`\frac{9}{2}id^4\omega _e^1se^2,G\left(2\omega _e\omega _s\right)se^2`$ |
| $`\left|C\right|^2\left|A\right|^2C`$ | $`\frac{9}{2}id^4\omega _e^1ge^2,G\left(\omega _g+2\omega _e\right)ge^2`$ |
| $`\left|C\right|^2\left|A\right|^2C`$ | $`\frac{9}{2}id^4\omega _e^1ge^2,G\left(2\omega _e\omega _g\right)ge^2`$ |
| $`\left|B\right|^2\left|A\right|^2C`$ | $`9id^4\omega _e^1gse,G\left(\omega _g+\omega _s+\omega _e\right)gse`$ |
| $`\left|B\right|^2\left|A\right|^2C`$ | $`9id^4\omega _e^1gse,G\left(\omega _s+\omega _e\omega _g\right)gse`$ |
Before discussing the information contained in these tables, we note that the system (4.23) can be further simplified. Using a near-identity transformation:
$$(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C})=(A,B,C)+𝒪\left(|A|^2+|B|^2+|C|^2\right)$$
(4.25)
the system (4.23) can be transformed into a new system for $`\stackrel{~}{A},\stackrel{~}{B}`$ and $`\stackrel{~}{C}`$ of a very similar form, but with the following modifications:
* The real parts of the coefficients are the same but the imaginary parts may be modified.
* The $`\epsilon ^2\mathrm{\Xi }`$ terms are now replaced by terms of order $`\epsilon ^6`$
We refer to the system governing $`\stackrel{~}{A}`$, $`\stackrel{~}{B}`$ and $`\stackrel{~}{C}`$, obtained in the manner, after neglecting the $`\mathrm{\Xi }`$ terms, as a dispersive normal form.
While complicated in its details, there is a simple way to think about this normal form. We introduce the internal mode powers <sup>7</sup><sup>7</sup>7Strictly speaking, by (4.25) $`P`$, $`Q`$ and $`R`$ are only approximately equal, respectively, to the Goldstone, shape and edge mode powers.:
$$P=|\stackrel{~}{A}|^2,Q=|\stackrel{~}{B}|^2,R|\stackrel{~}{C}|^2.$$
(4.26)
The equations for the powers are:
$`P_t`$ $`=`$ $`2\epsilon ^2\left(\alpha _1^rP+\alpha _2^rQ+\alpha _3^rR\right)P`$ (4.27)
$`+2\epsilon ^4\left(\alpha _4^rP^2+\alpha _5^rPQ+\alpha _6^rQ^2+\alpha _7^rQR+\alpha _8^rR^2\right)P`$
$`Q_t`$ $`=`$ $`2\epsilon ^2\left(\beta _1^rP+\beta _2^rQ+\beta _3^rR\right)Q`$ (4.28)
$`+2\epsilon ^4\left(\beta _4^rP^2+\beta _5^rPQ+\beta _6^rQ^2+\beta _7^rQR+\beta _8^rR^2\right)Q`$
$`R_t`$ $`=`$ $`2\epsilon ^2\left(\gamma _1^rP+\gamma _2^rQ+\gamma _3^rR\right)R`$ (4.29)
$`+2\epsilon ^4\left(\gamma _4^rP^2+\gamma _5^rPQ+\gamma _6^rQ^2+\gamma _7^rQR+\gamma _8^rR^2\right)R,`$
where $`\alpha _j^r,\beta _j^r`$ and $`\gamma _j^r`$ denote the real parts of the coefficients $`\alpha _j^r,\beta _j^r`$ and $`\gamma _j^r`$ appearing in (4.23). The real parts do not change under the near identity change of variables: $`A\stackrel{~}{A},B\stackrel{~}{B},C\stackrel{~}{C}`$, and are therefore given by the real parts of the coefficients displayed in Table 1. In general, as the reader may note, a contribution to the above mentioned internal mode power equations of the form $`|A|^{2m_1}|B|^{2m_2}|C|^{2m_3}`$ comes from a frequency combination $`m_1\omega _g+m_2\omega _s+m_3\omega _e`$ landing in the band of continuous spectrum through a term
$$(M(K,g,s,e),G(m_1\omega _g+m_2\omega _s+m_3\omega _e)M(K,g,s,e))$$
(4.30)
where $`M`$ is the appropriate monomial combination of the relevant spatial parts.
As the discreteness parameter, $`d`$, varies the spectrum of the kink (the continuous spectrum and the number and location of the internal modes) changes. As indicated in Figure 6, for different values of $`d`$ various integer linear combinations of the internal mode frequencies, the simplest of which are those appearing as arguments of $`G()`$ in Table 1, may lie in the continuous spectrum.
How does this influence the character of the normal form?
Let $`\zeta `$ denote one such linear combination of frequencies. By (4.19):
$`G(\zeta )=P_cB^1(B\zeta )^1P_c,\zeta \sigma _{cont}(B)`$
$`G(\zeta )=P_c\mathrm{P}.\mathrm{V}.B^1(B\zeta )^1P_ci{\displaystyle \frac{\pi }{\zeta }}P_c\delta (B\zeta )P_c,\zeta \sigma _{cont}(B).`$ (4.31)
Each of the coefficients of the system listed in Table 1 is a positive multiple of an expression of the form $`if,G(\zeta )f`$, where $`f`$ is spatially localized. Therefore, if $`\zeta `$ does not lie in the continuous spectrum of $`B`$, the associated coefficient of the system for the powers $`P,Q,R`$, $`\alpha ^r=\mathrm{}\alpha ,\beta ^r=\mathrm{}\beta `$ or $`\gamma ^r=\mathrm{}\gamma `$, will be zero. On the other hand, if $`\zeta `$ does lie in the continuous spectrum of $`B`$, this coefficient will be of the form:
$`\mathrm{\Gamma }_\zeta `$ $``$ $`{\displaystyle \frac{\pi }{\zeta }}P_cf,\delta (B\zeta )P_cf`$
$`=`$ $`{\displaystyle \frac{\pi }{\zeta }}\left|_B[f](\zeta )\right|^2,`$
where $`_B`$ denotes the Fourier transform with respect to the continuous spectral part of the operator $`B`$. Generically, one has $`\mathrm{\Gamma }_\zeta `$ is strictly negative. Therefore, such resonances are associated with nonlinear damping of energy in the internal modes. This does not contradict the Hamiltonian character of the equations of motion. The $`\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C}`$ system is coupled to the dispersive system governing $`\eta `$; damping of the discrete mode amplitudes implies a transfer of energy from the discrete to continuum modes. The information contained in Figure 6 enables us to determine, for each $`d`$, which combinations of harmonics appear in the phonon band. Then Tables 1A, 1B and 1C, together with (4.27-4.29) give us the precise form of the internal mode power equations, from which we can ascertain the detailed behavior of solutions in a neighborhood of the static kink.
In the Table 2 (below) we present the form of the internal mode power equations for different ranges of the parameter, $`d`$. There are numerous changes in the form of these equations as $`d`$ varies, so for clarity, we indicate only the key transitions. These key transitions occur across values of $`d`$ where there is a topological change in the phase portrait of the system governing the internal mode powers. Such changes are found due to a change in the nature of the set of equilibria, e.g. going from a system with one line of equilibria to one where there are two lines of equilibria, or due to a change in the number of internal modes (jump in the dimensionality of the phase portrait). The latter transition occurs at $`d=d_e,d_e0.82`$, when a point eigenvalue emerges from the edge of the continuous spectrum and appears as a third internal mode, $`e`$, an edge mode, with corresponding frequency $`\omega _e`$.
Table 2: $`\varphi ^4`$ internal mode power equations
| Regimes of $`d`$ | Resonances | System Form |
| --- | --- | --- |
| $`I:d<0.5398`$ | $`\{2\omega _s\omega _g,\mathrm{}\}`$ | $`P_t=0,Q_t=\epsilon ^4\left\{PQ^2\right\}`$ |
| $`II:0.5398d<0.6145`$ | $`2\omega _g,\omega _g+\omega _s`$ | $`P_t=\epsilon ^2P\{P,Q\},`$ |
| $`2\omega _g\sigma _{cont}`$ | $`2\omega _s\omega _g`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2\left\{Q\right\}\}`$ |
| $`III:0.6145d<0.6364`$ | $`2\omega _g,2\omega _s\omega _g`$ | $`P_t=\epsilon ^2\left\{P^2\right\},`$ |
| | | $`Q_t=\epsilon ^4\left\{PQ^2\right\}`$ |
| $`IV:0.6364d<0.6679`$ | $`\omega _g+\omega _s`$ | $`P_t=\epsilon ^2\left\{PQ\right\}`$, |
| $`2\omega _g\sigma _{cont}`$ | $`2\omega _s\omega _g`$ | $`Q_t=\epsilon ^2\left\{PQ\{1,\epsilon ^2\left\{Q\right\}\}\right\}`$ |
| $`V:0.6679d<d_e`$ | $`\{3\omega _g,4\omega _g,2\omega _s\omega _g\}`$ | $`P_t=\epsilon ^4P^3\{1,\epsilon ^2P\}\epsilon ^2PQ\{1,\epsilon ^2P\},`$ |
| $`d_e0.82`$ | $`\{\omega _g+\omega _s,2\omega _g+\omega _s\}`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2\{P,Q\}\}`$ |
| $`VI:d_ed<0.9229`$ | $`\{4\omega _g,5\omega _g,2\omega _s\omega _g\}`$ | $`P_t=\epsilon ^2P\{Q\{1,\epsilon ^2P\},R\{1,\epsilon ^2P\},\epsilon ^4P^3\}+\mathrm{}`$ |
| $`\omega _e`$ appears | $`\omega _g+\omega _s,2\omega _g+\omega _s,2\omega _e\omega _s`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2\{P,Q,R\}\}+\mathrm{}`$ |
| | $`\omega _g+\omega _e,2\omega _g+\omega _e,\omega _e+\omega _s\omega _g`$ | $`R_t=\epsilon ^2PR\{1,\epsilon ^2\{P,Q\}\}\epsilon ^4QR^2+\mathrm{}`$ |
| $`VII:0.9229d<d_{}`$ | $`\left\{n\omega _g\right\}_{5n38},2\omega _s`$ | $`P_t=\epsilon ^{2n2}\left\{P^n\right\}+\mathrm{}`$ |
| $`d_{}1.2234`$ | $`\omega _g+\omega _s,\omega _g+\omega _e,2\omega _g+\omega _s,`$ | $`Q_t=\left\{\epsilon ^2Q^2\right\}+\mathrm{}`$ |
| $`2\omega _s\sigma _{cont}`$ | $`2\omega _g+\omega _e,2\omega _g+2\omega _s,\omega _g+2\omega _s,`$ | $`R_t=\epsilon ^2R\{Q,P\}+\mathrm{}`$ |
| | $`\omega _g+\omega _s+\omega _e,\omega _s+\omega _e,2\omega _e\omega _s`$ | |
| | $`2\omega _e\omega _g,\omega _e+\omega _s\omega _g,2\omega _s\omega _g`$ | |
| $`VIII:d_{}d`$ | $`n\omega _g\left(5nN\right),2\omega _s,2\omega _e`$ | $`P_t=\epsilon ^{2n2}\left\{P^n\right\}+\mathrm{}`$ |
| $`2\omega _e\sigma _{cont}`$ | | $`Q_t=\left\{\epsilon ^2Q^2\right\}+\mathrm{}`$ |
| | | $`R_t=\left\{\epsilon ^2R^2\right\}+\mathrm{}`$ |
Remark on notation: We illustrate the notation of the table with the following example. Consider the regime V. This regime can be broken into three subregimes illustrated in the following table:
Table 3: internal mode power equations for subregime $`V`$
| Subregime of $`d`$ | Resonances | System Form |
| --- | --- | --- |
| $`a)0.6679d<0.743`$ | $`3\omega _g,\omega _g+\omega _s`$ | $`P_t=\epsilon ^2C_1PQ\epsilon ^4C_2P^3,`$ |
| | $`2\omega _s\omega _g`$ | $`Q_t=\epsilon ^2C_3PQ\epsilon ^4C_4PQ^2`$ |
| $`b)0.743d<0.7622`$ | $`3\omega _g,4\omega _g`$ | $`P_t=\epsilon ^2C_1PQ\epsilon ^4C_2P^3\epsilon ^6C_5P^4,`$ |
| | $`\omega _g+\omega _s,2\omega _s\omega _g`$ | $`Q_t=\epsilon ^2C_3PQ\epsilon ^4C_4PQ^2`$ |
| $`c)0.7622d<0.7687`$ | $`3\omega _g,\omega _g+\omega _s,2\omega _s\omega _g`$ | $`P_t=\epsilon ^2C_1PQ\epsilon ^4C_2P^3\epsilon ^4C_6P^2Q\epsilon ^6C_5P^4`$ |
| | $`4\omega _g,2\omega _g+\omega _s`$ | $`Q_t=\epsilon ^2C_3PQ\epsilon ^4C_4PQ^2\epsilon ^4C_7P^2Q`$ |
| $`c)0.7687d<d_e`$ | $`\omega _g+\omega _s,2\omega _g+\omega _s`$ | $`P_t=\epsilon ^2C_1PQ\epsilon ^4C_6P^2Q\epsilon ^6C_5P^4`$ |
| $`d_e0.82`$ | $`4\omega _g,2\omega _s\omega _g`$ | $`Q_t=\epsilon ^2C_3PQ\epsilon ^4C_4PQ^2\epsilon ^4C_7P^2Q`$ |
Although the details of the system form change between regimes, the topological character of the phase portrait and therefore the qualitative nature of the solutions does not change; each phase portrait has $`P=0`$ as a stable line of equilibria. For regime III, Table 2 is read as follows: in this regime some or all resonances occur from among each of the indicated sets: $`\{3\omega _g,4\omega _g\}`$, and $`\{\omega _g+\omega _s,2\omega _s\omega _g,2\omega _g+\omega _s\}`$, giving rise to terms in the following manner:
$$k\omega _g+l\omega _s+m\omega _e\epsilon ^{2n2}P^kQ^lR^m,k,l,m\mathrm{ZZ}_+,k+l+m=n.$$
(4.32)
Thus the sets in curly brackets in Table 2 are to be viewed as columns of a “menu” from which one (or the dynamical system) chooses all or some items (resonant combinations) depending on the subregime of $`d`$. For a given subregime, this choice of subset gives rise to linear combination with nonnegative coefficients, $`C_j`$, of the corresponding monomials in $`P,Q`$ and $`R`$ in the power equations. Therefore $`P_t,Q_t,R_t0`$. Generically , we have $`C_j>0`$.
It is straightforward to analyze the sets of equilibria and their stability for each of the systems in Table 2. These take the form of constant vectors with at most one nonzero component. These states are dynamically stable. If we take the view that the system of power equations determines the nonlinear dynamics near the static kink we anticipate that:
* the zero solution of the internal mode power equations corresponds to the ground state static kink solution discrete nonlinear equation. We denote the ground state kink by $`K_{gs}`$.
* a nonzero equilibrium state with $`Q0`$ corresponds to a time periodic solution:
$$u_iK_{gs,i}+\mathrm{cos}(\omega _st)s_i.$$
We call such a periodic solution of the full nonlinear dynamical system, which would have the same symmetry as the kink, a wobbling kink, which we designate by $`sW`$ or simply $`W`$.
* a nonzero equilibrium state with $`P0`$ power equations corresponds to a time periodic solution:
$$u_iK_{gs,i}+\mathrm{cos}(\omega _gt)g_i.$$
We call such a periodic solution a g-wobbling kink. We denote this state by $`gW`$.
* a nonzero equilibrium state with $`R0`$ power equations corresponds to a time periodic solution:
$$u_iK_{gs,i}+\mathrm{cos}(\omega _et)g_i.$$
Such a periodic solution of the full nonlinear dynamical system, would not have the same symmetry as the kink. We call such a periodic solution an e-wobbling kink. We denote this state by $`eW`$.
Below, we present a table of the kinds of static and periodic states anticipated by the normal form / internal mode power equation analysis.
Table 4: $`\varphi ^4`$ normal form, equilibria and anticipated coherent structures
| Regime of $`d`$ | Equilibria | Coherent structures |
| --- | --- | --- |
| I: $`d<0.5398`$ | $`\{(P,0):P0\}`$ | $`K_{gs},W,gW`$ |
| | $`\{(0,Q):Q0\}`$ | |
| II-III: $`0.5398d<0.6364`$ | $`\{(0,Q):Q0\}`$ | $`K_{gs},W`$ |
| IV: $`0.6364d<0.6679`$ | $`\{(0,Q):Q0\}`$ | $`K_{gs},W,gW`$ |
| | $`\{(P,0):P0\}`$ | |
| V: $`0.6679d<d_e`$ | $`\{(0,Q):Q0\}`$ | $`K_{gs},W`$ |
| $`d_e0.82`$ | | |
| VI: $`d_ed<0.9229`$ | $`\{(0,Q,0):Q0\}`$ | $`K_{gs},W,eW`$ |
| | $`\{(0,0,R):R0\}`$ | |
| VII: $`0.9229d<1.2234`$ | $`\{(0,0,R):R0\}`$ | $`K_{gs},eW`$ |
| $`d_{}1.2234`$ | | |
| VIII: $`d_{}d`$ | $`\left\{(0,0,0)\right\}`$ | $`K_{gs}`$ |
### 4.2 The normal form for discrete sine-Gordon
The procedure for deriving the normal form for the internal mode amplitudes and the internal mode power equations presented in sections 3 and 4.1 can be applied to the discrete sine-Gordon equation as well. The implementation is actually simpler because for discrete SG there are only two internal modes: $`g`$ and $`e`$; see section 2. Therefore, the decomposition of the solution is:
$$u_i(t)=K_i+\epsilon a(t)g_i+\epsilon b(t)e_i+\epsilon ^2\eta _i(t),$$
(4.33)
where
$`g,\eta (t)=e,\eta (t)=0`$
$`P_c\eta \eta g,\eta (t)ge,\eta (t)e=\eta (t)`$
As in the previous subsection the amplitude equations for the slowly varying modulation functions can be obtained:
$`A_t`$ $`=`$ $`\epsilon ^2\left(\alpha _1|A|^2+\alpha _2|B|^2\right)A+\epsilon ^4\alpha _3|B|^4A+\mathrm{\Xi }_A(A,B,t)`$
$`B_t`$ $`=`$ $`\epsilon ^2\left(\beta _1|A|^2+\beta _2|B|^2\right)A+\epsilon ^4\beta _3|B|^4A+\mathrm{\Xi }_B(A,B,t)`$ (4.34)
The coefficients $`\alpha _j,\beta _j`$ can be evaluated along the lines detailed in the pervious section and are tabulated below.
Table 5A : Principal SG g-mode coefficients
| $`\left|A\right|^2A`$ | $`\frac{1}{8}d^4\omega _{g}^{}{}_{}{}^{1}\mathrm{sin}Kg^2,G\left(2\omega _g\right)\mathrm{sin}Kg^2`$ |
| --- | --- |
| $`\left|A\right|^4A`$ | $`\frac{3}{4\left(3!\right)^2}d^4\omega _g^1\mathrm{cos}Kg^3,G\left(3\omega _g\right)\mathrm{cos}Kg^3`$ |
| $`\left|B\right|^2A`$ | $`\frac{1}{2}d^4\omega _g^1\mathrm{sin}Kge,G\left(\omega _g+\omega _e\right)\mathrm{sin}Kge`$ |
| $`\left|B\right|^4A`$ | $`\frac{9}{4\left(3!\right)^2}d^4\omega _{g}^{}{}_{}{}^{1}\mathrm{cos}Kge^2,G\left(\omega _g+2\omega _e\right)\mathrm{cos}Kge^2`$ |
| $`\left|A\right|^2\left|B\right|^2A`$ | $`\frac{18}{4\left(3!\right)^2}d^4\omega _g^1\mathrm{cos}Kg^2e,G(2\omega _g+\omega _e)\mathrm{cos}Kg^2e>`$ |
Table 5B : Discrete SG e-mode coefficients
| $`\left|B\right|^2B`$ | $`\frac{1}{8}d^4\omega _e^1\mathrm{sin}Ke^2,G\left(\omega _e\right)\mathrm{sin}Ke^2`$ |
| --- | --- |
| $`\left|B\right|^4B`$ | $`\frac{3}{4\left(3!\right)^2}d^4\omega _e^1\mathrm{cos}Ke^3,G\left(3\omega _e\right)\mathrm{cos}Ke^3`$ |
| $`\left|A\right|^2B`$ | $`\frac{1}{2}d^4\omega _e^1\mathrm{sin}Kge,G\left(\omega _g+\omega _e\right)\mathrm{sin}Kge`$ |
| $`\left|A\right|^4B`$ | $`\frac{9}{4\left(3!\right)^2}d^4\omega _e^1\mathrm{cos}Keg^2,G\left(\omega _e+2\omega _g\right)\mathrm{cos}Keg^2`$ |
| $`\left|A\right|^2\left|B\right|^2B`$ | $`\frac{18}{4\left(3!\right)^2}d^4\omega _e^1\mathrm{cos}Ke^2g,G\left(\omega _g+2\omega _e\right)\mathrm{cos}Ke^2g`$ |
| $`\left|A\right|^2\left|B\right|^2B`$ | $`\frac{18}{4\left(3!\right)^2}d^4\omega _e^1\mathrm{cos}Ke^2g,G\left(2\omega _e\omega _g\right)\mathrm{cos}Ke^2g`$ |
As with discrete $`\varphi ^4`$, the details of the internal mode power equations change as $`d`$ varies due to the different types of resonances with the continuous spectrum which may occur. For discrete SG there are essentially only two regimes, and one value of the parameter $`d`$ across which there is a topological change in the phase portraits. Since the detailed picture is simpler we tabulate it in greater detail.
Figure 7 displays the variation of the internal mode frequencies $`(\omega _g,\omega _e)`$, certain multiples of them and certain other integer linear combinations of them. As with discrete $`\varphi ^4`$, transitions in the structure of the normal form occur across values of $`d`$ for which there is a change in the set of integer linear combinations which lie in the band of continuous spectrum.
The precise normal form and power equations can be worked out using the coefficient Table 5 and the expression for $`G(\zeta )`$, (4.31). Notice that only the leading order terms are given in these internal mode power equations.
Table 6: SG internal mode power equations
| Regime of $`d`$ | Resonances | System Form |
| --- | --- | --- |
| I: $`d<d_e0.515`$ | None | $`P_t=0`$ |
| II: $`d_ed<.565`$ | $`2\omega _e\omega _g,\left\{3\omega _g\omega _e\right\}`$ | $`P_t=\epsilon ^6\left\{P^3Q\right\},Q_t=\epsilon ^4\left\{PQ^2\right\}`$ |
| III: $`0.565d<0.65`$ | $`2\omega _g,2\omega _e\omega _g`$ | $`P_t=\epsilon ^2\left\{P^2\right\}`$, $`Q_t=\epsilon ^4\left\{PQ^2\right\}`$ |
| $`0.65d<0.7`$ | $`2\omega _g,\omega _g+\omega _e`$ | $`P_t=\epsilon ^2P\{P,Q\}`$, |
| | | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2Q\}`$ |
| $`0.7d<0.76`$ | $`2\omega _g,3\omega _g,\omega _g+\omega _e,`$ | $`P_t=\epsilon ^2P\{P,Q,\epsilon ^4P^2\}`$, |
| | $`2\omega _e\omega _g`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2Q\}`$ |
| $`0.76d<0.785`$ | $`3\omega _g,\omega _g+\omega _e,`$ | $`P_t=\epsilon ^2P\left\{Q\right\}\epsilon ^2\left\{P^2\right\}`$, |
| | $`2\omega _e\omega _g`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2Q\}`$ |
| $`0.785d<0.8`$ | $`3\omega _g,4\omega _g`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^2P^2,\epsilon ^6P^4\}`$, |
| | $`\omega _g+\omega _e,2\omega _e\omega _g`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2PQ\}`$ |
| $`0.8d<0.847`$ | $`3\omega _g,4\omega _g,2\omega _e\omega _g`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^2P^2,\epsilon ^2PQ,\epsilon ^4P^3\}`$, |
| | $`\omega _g+\omega _e,2\omega _g+\omega _e`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2Q,\epsilon ^2P\}`$ |
| $`0.847d<d_{}`$ | $`3\omega _g,4\omega _g,5\omega _g,2\omega _e\omega _g`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^2P^2,\epsilon ^4P^3,\epsilon ^2PQ,\epsilon ^6P^4\}`$, |
| $`d_{}0.86`$ | $`\omega _g+\omega _e,2\omega _g+\omega _e`$ | $`Q_t=\epsilon ^2PQ\{1,\epsilon ^2Q,\epsilon ^2P\}`$ |
| IV: $`d_{}d<0.9`$ | $`3\omega _g,4\omega _g,5\omega _g,2\omega _e`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^2P^2,\epsilon ^4P^3,\epsilon ^2PQ\}`$, |
| | $`\omega _g+\omega _e,2\omega _g+\omega _e,2\omega _e\omega _g`$ | $`Q_t=\epsilon ^2Q\{P,Q,\epsilon ^2PQ,\epsilon ^2P^2\}`$ |
| $`0.9d<0.99`$ | $`4\omega _g,5\omega _g,2\omega _e`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^4P^3,\epsilon ^2PQ,\epsilon ^4P^4\}`$, |
| | $`\omega _g+\omega _e,2\omega _g+\omega _e,2\omega _e\omega _g`$ | $`Q_t=\epsilon ^2Q\{P,Q,\epsilon ^2PQ,\epsilon ^2P^2\}`$ |
| $`0.99d<1.003`$ | $`5\omega _g,2\omega _e,2\omega _e\omega _g`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^2PQ,\epsilon ^6P^4\}`$, |
| | $`\omega _g+\omega _e,2\omega _g+\omega _e`$ | $`Q_t=\epsilon ^2Q\{Q,\epsilon ^2PQ,\epsilon ^2P^2\}`$ |
| $`1.003d`$ | $`n\omega _g\left(n5\right),2\omega _e,\omega _g+\omega _e`$ | $`P_t=\epsilon ^2P\{Q,\epsilon ^2PQ,\epsilon ^{2n2}P^{n1},\epsilon ^2Q^2,\mathrm{}\}`$, |
| | $`2\omega _g+\omega _e,\omega _g+2\omega _e`$ | $`Q_t=\epsilon ^2Q\{P,Q,\epsilon ^4P^2,\epsilon ^4PQ,\mathrm{}\}`$ |
The inferred coherent structures are displayed in the following table.
Table 7: SG normal form, equilibria and anticipated coherent structures
| Regime of $`d`$ | Equilibria | Coherent structures |
| --- | --- | --- |
| I: $`d<d_e`$ | $`P0`$ | $`K_{gs},gW`$ |
| II: $`d_ed<0.565`$ | $`\{(P,0),(0,Q):P0,Q0\}`$ | $`K_{gs},W,gW`$ |
| III: $`0.565d<d_{}`$ | $`\{(0,Q):Q0\}`$ | $`K_{gs},W`$ |
| $`d_{}0.86`$ | | |
| IV: $`d_{}d`$ | $`\left\{(0,0)\right\}`$ | $`K_{gs}`$ |
## 5 Time periodic solutions
### 5.1 Existence of wobbling kinks
From the discussion of the previous section, we expect that the dispersive normal form we have derived for discrete nonlinear wave equations (in particular, discrete SG and $`\varphi ^4`$) models the dynamics in a neighborhood of the static kink solution. This normal form captures the effects of nonlinear resonant interactions of the internal and continuum mode fluctuations about the kink. Tables 4 and 7 indicate for discrete $`\varphi ^4`$ discrete SG the anticipated existence of periodic solutions which bifurcate from the static kink solution. In this section we prove, under suitable nonresonance hypotheses, the existence of such time periodic solutions.
In tables 3 and 5 we see that these normal forms fall into two categories (i) those for which the origin is the only equilibrium ($`d`$ sufficiently large, $`d>d_{}`$) (ii) $`d`$ is such that the normal form admits one or two lines of equilibria. In case (i) we see that the zero solution is asymptotically stable. By (3.2) and (4.33) this suggests that the ground state kink is a stable attractor. At $`d_{}`$ there is a bifurcation in the phase portrait; for $`d<d_{}`$, zero is no longer an isolated equilibrium. Rather, the normal form now has a line of stable equilibria passing through the origin. This then suggests an approximate solution of the nonlinear wave equation of the form:
$$u_jK_j+\epsilon B_{eq}\mathrm{cos}(\mathrm{\Omega }t)\xi _{\mathrm{\Omega },j},j\mathrm{ZZ}$$
(5.1)
where $`\xi _\mathrm{\Omega }`$ denotes an internal mode $`\mathrm{\Omega }`$ its corresponding frequency; $`B\chi _\mathrm{\Omega }=\mathrm{\Omega }\chi _\mathrm{\Omega }`$ and $`\epsilon `$ is small. If such a family of periodic solutions exists, we say this family bifurcates from the kink solution in the direction of the internal mode, $`\xi _\mathrm{\Omega }`$. Such directions for bifurcation are available for $`d<d_{}`$.
For $`d<d_{}`$ are there bifurcating time periodic solutions?
In view of the analysis of the regime $`d>d_{}`$, one must be cautious about reaching a conclusion about the large time behavior on the basis of the above approximation. In fact we learn from the regime $`d>d_{}`$ that radiation damping, due to resonant coupling of oscillations to the continuum, may lead to the slow decay of a solution which appears to be time-periodic on shorter time scales.
On the other hand, given the approximate periodic solution (5.1) it is certainly natural to attempt a perturbative (in $`ϵ`$) construction of a true time-periodic solution. The approach we use is the Poincaré continuation, which is standard in the persistence theory of periodic solutions of systems of ordinary differential equations . As in typical perturbation expansions, one expects a hierarchy of linear inhomogeneous problems governing contributions to the expansion at various orders in $`ϵ`$. Solvability at each order requires that one can arrange for the solution to each inhomogeneous problem to be time periodic of some common period. That this can be achieved in the problem at hand follows from the following nonresonance condition, verified in our numerical investigation of the spectrum of the kink:
(NR) Assume that no integer multiple of the internal mode frequency $`\mathrm{\Omega }`$ lies in the continuous spectrum of $`B`$.
A precise statement of the result is as follows:
###### Theorem 5.1
Let $`K=\{K_i\}_{i\mathrm{ZZ}}`$ denote a ground state kink. Assume that the linear operator, $`B`$, acting on the space $`l^2(\mathrm{ZZ})`$, has a simple eigenvalue, $`\mathrm{\Omega }`$, with corresponding eigenfunction, $`\xi _\mathrm{\Omega }`$ satisfying $`B\xi _\mathrm{\Omega }=\mathrm{\Omega }\xi _\mathrm{\Omega }`$. Also assume the nonresonance condition (NR). Then, in a neighborhood of $`K`$, there is a curve of solutions $`ϵy(t;ϵ)`$ passing through the ground state kink with the following properties:
* There is a number $`ϵ_0>0`$ such that $`y(t;ϵ)`$ is defined for all $`ϵ<ϵ_0`$.
* $`y(t;0)=K`$
* $`y(t;ϵ)`$ has period $`2\pi /\omega (ϵ)`$ in $`t`$, where $`\omega ^2(ϵ)=\mathrm{\Omega }^2(1+𝒪(ϵ))`$ is a smooth function with $`\omega (0)=\mathrm{\Omega }`$.
* $`y(t;ϵ)l^2(\mathrm{ZZ})`$
* $`y_i(t;ϵ)=K_i+ϵ\mathrm{cos}(\omega (ϵ)t)\xi _{\mathrm{\Omega },i}+𝒪(ϵ^2).`$
The next result, a simple consequence of the proof of Theorem 5.1, yields the class of solutions which we call wobbling kinks.
###### Corollary 5.1
For discrete SG, the periodic solution bifurcating from the kink in the direction of the spatially odd edge mode ($`Be=\omega _ee`$) has the same symmetry as the SG kink around its center. For discrete $`\varphi ^4`$ the periodic solution bifurcating from the kink in the direction of the spatially odd shape mode ($`Bs=\omega _ss`$) has the same symmetry as the $`\varphi ^4`$ kink.
We seek a solution in the form:
$$u_i=K_i+b(t)\xi _{\mathrm{\Omega },i}+\eta _i,i\mathrm{ZZ}.$$
(5.2)
Substitution of (5.2) into (2.14) and then projection of the resulting equation separately onto the eigenvector $`\xi _\mathrm{\Omega }`$ and its orthogonal complement yields the coupled system for the shape mode amplitude and the radiation components <sup>8</sup><sup>8</sup>8Here we employ the Taylor expansion for the function $`W=V^{}`$:
$$W(K+M)=W(K)+W^{}(K)M+_0^1(1\theta )W^{\prime \prime }(K+\theta M)𝑑\theta M^2$$
$`_t^2b+\mathrm{\Omega }^2b`$ $`=`$ $`d^2\xi _\mathrm{\Omega },𝒱[b\xi _\mathrm{\Omega }+\eta ](b\xi _\mathrm{\Omega }+\eta )^2`$ (5.3)
$`_t^2\eta +B^2\eta `$ $`=`$ $`d^2P_c𝒱[b\xi _\mathrm{\Omega }+\eta ](b\xi _\mathrm{\Omega }+\eta )^2,`$ (5.4)
where
$$𝒱[M]=_0^1(1\theta )V^{\prime \prime \prime }(K+\theta M)𝑑\theta .$$
(5.5)
###### Remark 5.1
It is possible to formulate the proof of Theorem 5.1, as essentially a consequence results in (Chapter 14, Theorem 2.1), applied to (5.3-5.4); a simple generalization of the result to an infinite dimensional setting is required. In this context, hypothesis (NR), corresponds to the hypothesis that $`1`$ be a simple Floquet multiplier of the unperturbed linear problem obtained by setting the right hand sides of (5.3-5.4) equal to zero. In this paper we present a direct proof of Theorem 5.1 in order to make our study elementary and self-contained.
We shall be considering small perturbations of the static kink, and therefore introduce rescaled internal mode and radiation components:
$$b=ϵb_1,\eta =ϵ\eta _1,$$
where $`ϵ`$ is a small parameter. Since the nonlinear wave equations considered are autonomous, we can’t a priori specify the period, we leave the period unspecified at this stage and introduce a new time variable, $`\tau `$, with respect to which the sought-for periodic solution is $`2\pi `$ periodic. Let
$$t=\omega (ϵ)^1\tau ,\omega (0)=\mathrm{\Omega },$$
(5.6)
and
$$b_1(t)=\beta _1(\tau ),\beta _1(\tau +2\pi )=\beta _1(\tau ).$$
Then, equations (5.3-5.4) become:
$`\left(\omega ^2(ϵ)_\tau ^2+\mathrm{\Omega }^2\right)\beta _1(\tau )`$ $`=`$ $`ϵ_\beta (\beta _1,\eta _1;ϵ)`$
$`_\beta (\beta _1,\eta _1;ϵ)`$ $``$ $`d^2\xi _\mathrm{\Omega },𝒱[ϵ(\beta _1\xi _\mathrm{\Omega }+\eta _1)](\beta _1\xi _\mathrm{\Omega }+\eta _1)^2`$
$`\left(\omega ^2(ϵ)_\tau ^2+B^2\right)\eta _1(\tau ,i)`$ $`=`$ $`ϵ_\eta (\beta _1,\eta _1;ϵ)`$
$`_\eta (\beta _1,\eta _1;ϵ)`$ $``$ $`d^2P_c𝒱[ϵ(\beta _1\xi _\mathrm{\Omega }+\eta _1)](\beta _1\xi _\mathrm{\Omega }+\eta _1)^2`$ (5.8)
Let
$`\omega ^2(ϵ)=\mathrm{\Omega }^2(1+ϵ\sigma (ϵ)),`$ (5.9)
where $`\sigma (ϵ)`$ is to be determined. Then (LABEL:eq41) becomes
$`\left(_\tau ^2+1\right)\beta _1={\displaystyle \frac{ϵ}{1+ϵ\sigma (ϵ)}}[\sigma (ϵ)\beta _1+\mathrm{\Omega }^2_\beta (\beta _1,\eta _1;ϵ)].`$ (5.10)
We extract the leading order behavior by defining:
$$\beta _1=\mathrm{cos}(\tau )+ϵ\beta _2,$$
where
$$(_{\tau }^{}{}_{}{}^{2}+1)\beta _2=\frac{1}{1+ϵ\sigma (ϵ)}[\sigma (ϵ)\mathrm{cos}(\tau )+ϵ\sigma (ϵ)\beta _2+\mathrm{\Omega }^2_\beta (\mathrm{cos}(\tau )+ϵ\beta _2,\eta _1;ϵ)$$
(5.11)
By the Fredholm alternative, equation (5.11) has an even $`2\pi `$ periodic solution if and only if
$$P_{\mathrm{cos}}(\sigma \mathrm{cos}(\tau )+ϵ\sigma \beta _2+\mathrm{\Omega }^2_\beta )=0.$$
(5.12)
Here,
$`P_{\mathrm{cos}}g=\pi ^1{\displaystyle _0^{2\pi }}\mathrm{cos}(\mu )g(\mu )𝑑\mu `$ (5.13)
is the projection onto the even $`2\pi `$ periodic null space of $`_\tau ^2+1`$. The orthogonality constraint (5.12) can be rewritten as
$$\pi \sigma +ϵ\sigma _0^{2\pi }\mathrm{cos}(\tau )\beta _2(\tau )𝑑\tau +\omega _s^2_0^{2\pi }\mathrm{cos}(\tau )_\beta (\mathrm{cos}(\tau )+ϵ\beta _2(\tau ),\eta _1;ϵ)𝑑\tau =0$$
(5.14)
We have therefore reformulated the problem of finding a periodic solution of the discrete nonlinear wave equation, or equivalently (5.3-5.4), as the problem of finding $`2\pi `$ periodic solutions in $`\tau `$, $`(\beta _2(ϵ),\eta _1(ϵ),\sigma (ϵ))`$, of the system:
$$F(\beta _2,\eta _1,\sigma ;ϵ)=0,$$
where $`F=(F_1,F_2,F_3)^t`$, is defined by:
$`F_1(\beta _2,\eta _1,\sigma ;ϵ)=(_\tau ^2+1)\beta _2`$
$`(1+ϵ\sigma )^1\left[\sigma \mathrm{cos}(\tau )+ϵ\sigma \beta _2+\mathrm{\Omega }^2_\beta (\mathrm{cos}(\tau )+ϵ\beta _2,\eta _1;ϵ)\right]`$ (5.15)
$`F_2(\beta _2,\eta _1,\sigma ;ϵ)=\left(\mathrm{\Omega }^2(1+ϵ\sigma )_\tau ^2+B^2\right)\eta _1(\tau ,i)ϵ_\eta (\mathrm{cos}(\tau )+ϵ\beta _2,\eta _1;ϵ)`$
(5.16)
$`F_3(\beta _2,\eta _1,\sigma ;ϵ)=(1+ϵ\sigma )^1[\pi \sigma +ϵ\sigma {\displaystyle _0^{2\pi }}\mathrm{cos}(\tau )\beta _2d\tau `$
$`+\mathrm{\Omega }^2{\displaystyle _0^{2\pi }}\mathrm{cos}(\tau )_\beta (\mathrm{cos}(\tau )+ϵ\beta _2(\tau ),\eta _1;ϵ)d\tau ]`$ (5.17)
We view $`F`$ as a mapping of $`(\zeta ,ϵ)𝒳F(\zeta ,ϵ)𝒴`$. Here, $`𝒳`$ and $`𝒴`$ are defined by:
$`𝒳`$ $`:`$ $`\zeta =(\beta ,\eta ,\sigma )\mathrm{such}\mathrm{that}`$
$`\beta =\beta (\tau )H^2,even,\mathrm{and}2\pi \mathrm{periodic}`$
$`\eta =\eta (\tau ,)H^2even,\mathrm{and}2\pi \mathrm{periodic}\mathrm{in}\tau `$
$`\mathrm{with}\mathrm{values}\mathrm{in}\mathrm{the}\mathrm{space}\mathrm{of}l^2(\mathrm{ZZ})\mathrm{functions}\mathrm{and}\sigma \mathrm{IR}`$
$`𝒴`$ $`:`$ $`\zeta =(\beta ,\eta ,\rho )\mathrm{such}\mathrm{that}`$
$`\beta =\beta (\tau )L^2,even,\mathrm{and}2\pi \mathrm{periodic}`$
$`\eta =\eta (\tau ,)L^2even,\mathrm{and}2\pi \mathrm{periodic}\mathrm{in}\tau `$
$`\mathrm{with}\mathrm{values}\mathrm{in}\mathrm{the}\mathrm{space}\mathrm{of}l^2(\mathrm{ZZ})\mathrm{functions}\mathrm{and}`$
$`\rho ={\displaystyle _0^{2\pi }}\beta (\mu )\mathrm{cos}(\mu )𝑑\mu .`$
We find a particular $`\zeta ^{(0)}𝒳`$ for which $`F(\zeta ^{(0)};0)=0`$, and then seek to construct curve of solutions $`ϵ\zeta (ϵ),\zeta (0)=\zeta ^{(0)}`$, for all sufficiently small $`ϵ`$ using the implicit function theorem . To find $`\zeta ^{(0)}`$ we set $`ϵ=0`$ and consider the system $`F(\zeta ;0)=0`$. Taking $`\eta _1=\eta _1^{(0)}0`$ and $`\sigma =\sigma ^{(0)}=0`$, we find that $`\beta _2^{(0)}`$ satisfies the equation:
$$\left(_\tau ^2+1\right)\beta _2^{(0)}=d^2\mathrm{\Omega }^2\xi _\mathrm{\Omega },𝒱[0]\xi _\mathrm{\Omega }^2\mathrm{cos}^2(\tau ),$$
which has the solution
$$\beta _2^{(0)}=\frac{1}{2}d^2\mathrm{\Omega }^2\xi _\mathrm{\Omega }𝒱[0]\xi _\mathrm{\Omega }^2\left(1\frac{1}{3}\mathrm{cos}(2\tau )\right).$$
Thus, $`\zeta ^{(0)}=(\beta _2^{(0)},0,0)𝒳`$ satisfies $`F(\zeta ^{(0)};0)=0`$. We shall now use the implicit function theorem to continue this solution to show that this solution deforms uniquely to nearby solutions for $`ϵ`$ sufficiently small and nonzero.
To apply the implicit function theorem, it suffices to check that $`d_\zeta F(\zeta ^{(0)};0)`$ is bounded and invertible. A computation yields:
$$d_\zeta F(\zeta ^{(0)};0)=\left[\begin{array}{ccc}\hfill _{\tau }^{}{}_{}{}^{2}+1& \hfill 0& \hfill \mathrm{cos}(\tau )\\ \hfill 0& \hfill \mathrm{\Omega }^2_\tau ^2+B^2& \hfill 0\\ \hfill 0& \hfill 0& \hfill \pi \end{array}\right]$$
(5.18)
Invertibility of $`d_\zeta F(\zeta ^{(0)};0)`$ can be shown by solving the system of inhomogeneous equations:
$`d_\zeta F\delta \zeta =R`$ (5.19)
where
$$\delta \zeta =\left[\begin{array}{c}\hfill \delta \beta _2\\ \hfill \delta \eta _1\\ \hfill \delta \sigma \end{array}\right],R=\left[\begin{array}{c}\hfill B\\ \hfill E\\ \hfill \mathrm{\Sigma }\end{array}\right]\mathrm{and}\mathrm{\Sigma }_0^{2\pi }\mathrm{cos}(\tau )B(\tau )𝑑\tau .$$
The third equation implies:
$$\delta \sigma =\pi ^1\mathrm{\Sigma }.$$
Substitution into the first equation yields the equation an $`\delta \beta _2`$:
$$\left(_\tau ^2+1\right)\delta \beta _2=B\pi ^1_0^{2\pi }\mathrm{cos}(\mu )B(\mu )𝑑\mu \mathrm{cos}(\tau ),$$
which has an even $`2\pi `$ periodic solution. Finally the second equation can be rewritten as
$`(\mathrm{\Omega }^2_{\tau }^{}{}_{}{}^{2}+B^2)\delta \eta _1=E={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\mathrm{cos}(n\tau )e_n`$ (5.20)
which can then be solved by Fourier series: $`\delta \eta _1=_{n=0}^{\mathrm{}}g_n\mathrm{cos}(n\tau )`$. We obtain
$`g_n=(B^2n^2\mathrm{\Omega }^2)^1e_n.`$ (5.21)
By (NR), $`e_n`$ is well-defined for each $`n`$. Note also that there is a strictly positive minimum distance of the set $`\{n\mathrm{\Omega }\}`$ to $`\sigma (B)`$. Therefore,
$$g_n_{l^2(\mathrm{ZZ})}\mathrm{dist}((n\mathrm{\Omega })^2,\sigma (B^2))^1E_n_{l^2(\mathrm{ZZ})},$$
from which we have boundedness of the inverse: $`\left(\mathrm{\Omega }^2_\tau ^2+B^2\right)^1`$:
$$\delta \eta _1_{H^2(\mathrm{IR};l^2(\mathrm{ZZ}))}CE_{L^2(\mathrm{IR};l^2(\mathrm{ZZ}))}.$$
By the implicit function theorem there is a number $`ϵ_0>0`$ such that for $`ϵ<ϵ_0`$ the nonlinear wave equation has a unique periodic solution of the form:
$`y_i(\tau ,ϵ)=K_i+(ϵ\mathrm{cos}(\tau )+ϵ^2\beta _2(\tau ))\xi _{\mathrm{\Omega },i}+ϵ\eta _1(i,\tau ;ϵ)`$ (5.22)
with $`\tau =\omega (ϵ)t`$, $`\omega ^2(ϵ)=\mathrm{\Omega }^2(1+ϵ\sigma (ϵ))`$ and $`\beta _2,\eta _1`$ $`2\pi `$ periodic functions of $`\tau `$. Note that $`\eta _1=𝒪(ϵ)`$, by the equation $`F_2=0`$, so that we have (5.1). This completes the proof of Theorem 5.1.
To prove Corollary 5.1, we note that by hypothesis, the direction of bifurcation is spatially odd. It is simple to check that the entire proof goes through with the spaces $`𝒳`$ and $`𝒴`$ additionally constrained to consist of functions $`\eta `$, such that $`\eta (,i)=\eta (,i)`$.
### 5.2 Do quasiperiodic solutions bifurcate from the kink?
If we attempt using the above method of proof to construct quasiperiodic solutions, we do not succeed. To be specific, suppose we seek to construct a quasiperiodic solution which is generated by the two internal modes $`\xi _{\mathrm{\Omega }_1}`$ and $`\xi _{\mathrm{\Omega }_1}`$. Formally, we seek $`u_i(t)=u(i,t)`$ in the form:
$`u(i,t)=K(i)+ϵ{\displaystyle \underset{j=1}{\overset{2}{}}}\alpha _j\xi _{\mathrm{\Omega }_j}(i)\mathrm{cos}(\mathrm{\Omega }_jt)+ϵ\eta _1(i,t)`$ (5.23)
Substitution into the nonlinear wave equation gives
$`\left(_t^2+B^2\right)\eta _1=`$ $`d^2ϵ𝒱\left[ϵ{\displaystyle \underset{j=1}{\overset{2}{}}}\alpha _j\xi _{\mathrm{\Omega }_j}\mathrm{cos}(\mathrm{\Omega }_jt)+ϵ\eta _1\right]`$
$`\times \left({\displaystyle \underset{j=1}{\overset{2}{}}}\alpha _j\xi _{\mathrm{\Omega }_j}\mathrm{cos}(\mathrm{\Omega }_jt)+\eta _1\right)^2.`$
Expansion of $`\eta _1(i,t)`$ as a power series in $`ϵ`$,
$$\eta _1(i,t)=\underset{j=0}{\overset{\mathrm{}}{}}ϵ^j\eta _1^{(j)}(i,t),$$
leads to a hierarchy of inhomogeneous equations of the form:
$$\left(_t^2+B^2\right)\eta _1^{(j)}=𝒮^{(j)}(t)$$
(5.24)
where as $`j`$ increases $`𝒮^{(j)}(t)`$ contains a finite sum of terms with time-frequencies of the form $`n_1\mathrm{\Omega }_1+n_2\mathrm{\Omega }_2`$ with $`|n_1|N_1^{(j)}`$, $`|n_2|N_2^{(j)}`$ and $`N_1^{(j)},N_2^{(j)}`$ increasing with $`j`$. Thus it is natural to solve for each $`\eta _1^{(j)}`$ as a truncated multiple Fourier series:
$$\eta _1^{(j)}=\underset{|n_1|N_1^{(j)},n_2N_2^{(j)}}{}e^{i(n_1\mathrm{\Omega }_1+n_2\mathrm{\Omega }_2)t}g_{n_1,n_2}^{(j)}.$$
Substitution into (5.24) yields:
$$\left(B^2(n_1\mathrm{\Omega }_1+n_2\mathrm{\Omega }_2)^2\right)g_{n_1,n_2}^{(j)}=𝒢_{n_1,n_2}^{(j)},$$
(5.25)
where $`𝒢_{n_1,n_2}^{(j)}`$ denotes the term in $`𝒮^{(j)}(t)`$ which is proportional to $`\mathrm{exp}(i(n_1\mathrm{\Omega }_1+n_2\mathrm{\Omega }_2)t)`$. In order to ensure the general solvability of (5.25) we need that for any $`n_1,n_2\mathrm{ZZ}`$,
$$n_1\mathrm{\Omega }_1+n_2\mathrm{\Omega }_2\pm \sigma (B).$$
Although it is nongeneric for a frequency in the point spectrum (internal mode frequency) to be hit, since the continuous spectrum (phonon band) is an interval, generically one has $`n_1\mathrm{\Omega }_1+n_2\mathrm{\Omega }_2\sigma _{cont}(B)`$, for infinitely many choices of $`n_1,n_2`$. For such choices the operator on the right hand side of (5.25) is not invertible and these resonant frequencies are an obstruction to solvability. A closer look at the set of resonances and their contribution on the dispersive normal form, would give insight into the lifetime of such eventually decaying quasiperiodic oscillations.
## 6 Large time behavior in a neighborhood of $`K_{gs}`$
In this section we combine the normal form analysis of section 4, and the existence theory for time periodic solutions of section 5 with numerical simulation to get a more detailed picture of the large time behavior of discrete $`\varphi ^4`$ and SG in a neighborhood of the ground state kink. We shall make repeated use of Tables 2 and 4 for $`\varphi ^4`$ and of Tables 6 and 7 for SG in which the normal form / power equations and coherent structures are tabulated. In interpreting these tables we recall that $`\epsilon `$, introduced in (3.2), measures the size of the component of the perturbation about the kink in the internal mode subspace. In full simulations of the evolution equation $`\epsilon `$ is typically of order $`10^1`$, and therefore some of the decay phenomena anticipated by our analysis, occur on very large time scales (e.g. $`\tau \epsilon ^2,\epsilon ^4`$ or longer) and are difficult to simulate accurately.
In each of the various $`d`$ ranges we proceed as follows: we mention the relevant coherent structures (the high energy kink is omitted because it is unstable). These are all anticipated by the normal form analysis and their existence is established rigorously in section 5. We then discuss the numerically observed large time behavior for different classes of initial conditions. This gives evidence of the attracting nature of the periodic solutions constructed in section 5 in various regimes of the discreteness parameter $`d`$.
As mentioned in the introduction, this is related to the final stages of the evolution of a propagating kink, pinned to a particular lattice site, and its damped oscillation to an asymptotic state within the Peierls-Nabarro potential.
### 6.1 Discrete $`\varphi ^4`$
Regimes II, III and V: The coherent structures of interest are: $`K_{gs}`$ and $`W`$. Numerical simulations show that the wobbling kink, $`W`$, is a local attractor.
Regimes I,IV: The coherent structures of interest are: $`K_{gs}`$ , $`W`$ and $`gW`$.
Experiment 1 (data given by $`K_{gs}`$ plus a small multiple of the Goldstone (even) $`g`$-mode): solution appears to approach $`gW`$.
Experiment 2 (data given by $`K_{gs}`$ plus a small multiple of the shape (odd) $`s`$-mode): solution appears to approach $`W`$.
Experiment 3 (data given by a general small perturbations of $`K_{gs}`$): $`gW`$ and $`W`$ appear to have basins of attraction. Although an exact determination of this basin is not analytically tractable, its projection onto the internal mode subspace, the span of $`\{g,s\}`$, can be approximated using the internal mode power equations and the explicit information on coefficients from Table 1. In particular, we find
$$\frac{dQ}{dP}=\frac{\omega _g}{\omega _s}$$
from which one gets the prediction that for data with
$$|B(0)|^2<\frac{\omega _g}{\omega _s}|A(0)|^2,$$
solutions asymptotically approach a g-wobbler, $`gW`$, while for
$$|B(0)|^2>\frac{\omega _g}{\omega _s}|A(0)|^2,$$
solutions asymptotically approach a wobbler, $`W`$.
Regime VI: The phase portrait of the normal form jumps from dimension $`2`$ to dimension $`3`$ due to the appearence of a (spatially even) edge mode. In addition to the ground state kink, $`K_{gs}`$, there are time-periodic (wobbling solutions) $`W`$ and $`eW`$. Numerical simulations indicate the presence of very long lived quasiperiodic oscillations about the kink. These appear to be oscillations of the form represented in the first two terms of (5.23). Strong evidence of the eventual decay of such oscillations can be seen as follows:
* In section 5.2 we have shown that an attempt to construct quasiperiodic solutions breaks down due to a high order resonance, i.e. $`n_1\omega _s+n_2\omega _e\sigma _{cont}(B)`$, where $`|n_1|+|n_2|`$ is large. In this particular case, we have $`\omega _s+2\omega _e\sigma _{cont}(B)`$; see Table 2.
* Such resonances correspond to obstructions in the power equations to equilibrium solutions of the form: $`(0,Q,R)`$ with both $`Q`$ and $`R`$ nonzero. The obstructing term is the term $`\epsilon ^4QR^2`$ in the $`R`$ equation.
* Linearization of the power equations about any equilibrium point $`(0,Q_{eq},0)`$ or $`(0,0,R_{eq})`$ shows that each line of equilibria is asymptotically stable, corresponding to the conclusion that quasiperiodic oscillations will damp with a resulting asymptotically periodic state, $`W`$ or $`eW`$, depending on initial conditions. The time scale of decay of these oscillations, $`\tau `$, is set by the order in $`\epsilon `$ of the obstructing terms, e.g. $`\tau \epsilon ^4`$.
Regime VII: The coherent structures of interest are: $`K_{gs}`$ and $`eW`$.
Experiment 1: For data which is an odd perturbation of $`K_{gs}`$. $`K_{gs}`$ is the attractor.
Experiment 2: For general data, $`eW`$ appears to be a local attractor.
Regime VIII: The coherent structure of interest is $`K_{gs}`$.
Experiment 1 (data given by $`K_{gs}`$ plus an small perturbation in the direction of the (odd) shape mode): approach to $`K_{gs}`$ at a rate $`𝒪(t^{\frac{1}{2}})`$.
Experiment 2 (data given by $`K_{gs}`$ plus an small perturbation in the direction of the (even) edge mode): approach to $`K_{gs}`$ at a rate $`𝒪(t^{\frac{1}{2}})`$.
Experiment 3 (data given by $`K_{gs}`$ plus an small perturbation in the direction of the (even) Goldstone mode): perturbation from $`K_{gs}`$ appears to decay but at a different rate and on much longer time scales than in Experiments 1 and 2.
### 6.2 Discrete sine-Gordon
Following the model of the previous subsection, we indicate relevant coherent structures and briefly discuss the dynamics in a neighborhood of the ground state kink for the discrete sine-Gordon equation.
Regime I: The coherent structures of interest are: $`K_{gs}`$ and $`gW`$. For general data, $`gW`$ appears to be a local attractor.
Regime II: The coherent structures of interest are: $`K_{gs}`$, $`W`$ and $`gW`$. The dynamical behavior is as, for example, in regime IV for the $`\varphi ^4`$ model.
Regime III: The coherent structures of interest are $`K_{gs}`$ and $`W`$. Numerical simulations show that $`W`$ is a local attractor.
Regime IV: For all $`dd_{}`$, there are positive integer multiples of both $`\omega _g`$ and $`\omega _e`$ which fall in the phonon band, $`\sigma _{cont}(B)`$. Therefore, hypothesis (NR) of section 5 is not satisfied for either internal mode, and the existence theory of periodic solutions breaks down. A general perturbation about $`K_{gs}`$ will excite the internal modes and these resonances are responsible for transfer of energy from the internal modes to $`K_{gs}`$ and to radiation modes. The asymptotic state is observed to the ground state kink, with the details of the damped oscillatory approach to it, being predicted by the various $`d`$ dependent normal forms.
Within each subregime of regime IV, the power equations are, up to terms which can be shown to be negligible for large $`t`$, of the form
$`P_t`$ $`=`$ $`\epsilon ^{2n}P^n`$
$`Q_t`$ $`=`$ $`\epsilon ^2PQ.`$
It follows that
$$P(t)=P_0(1+(n1)P_0^{n1}\epsilon ^{2n}t)^{\frac{1}{n1}}$$
and that $`Q(t)`$ is decaying at exponentially. These predicted power laws for $`P(t)`$ (on a time scale of order $`\epsilon ^{2n}`$, and the implied very rapid decay of $`Q`$ are, to a good degree of approximation, observed in numerical simulations of the discrete sine-Gordon equation. In particular, using the form of the near-identity transformation ($`4.24`$), we obtain for the decay rates (of initial conditions with general data):
* $`0.86d_{}d<0.9`$$`n=3,P(t)=𝒪(t^{\frac{1}{2}})`$ $`|a(t)|,|b(t)|𝒪(t^{\frac{1}{4}})`$
* $`0.9d<0.99`$$`n=4,P(t)=𝒪(t^{\frac{1}{3}})`$, $`|a(t)|,|b(t)|𝒪(t^{\frac{1}{6}})`$
* $`0.99d<1.003`$$`n=5,P(t)=𝒪(t^{\frac{1}{4}})`$, $`|a(t)|,|b(t)|𝒪(t^{\frac{1}{8}})`$
Note, however, that for purely odd initial perturbations around the ground state kink in this regime, the edge mode projection decays as $`t^{1/2}`$, as dictated by the $`\epsilon ^2Q^2`$ term.
## 7 Summary and open questions
Summary: In this paper we have considered the large time behavior of solutions to discrete nonlinear wave equations, particularly the discrete sine-Gordon and discrete $`\varphi ^4`$ models, for initial conditions which are a small perturbation of a stable (ground state) kink.
(1) We have proved in the regime where the discreteness parameter is sufficiently small, corresponding to sufficiently large lattice spacing or weak coupling of neighboring oscillators, the existence of various classes of time periodic solutions. These include the, in the case of the $`\varphi ^4`$ the wobbling kinks ($`W`$), which have the same spatial symmetry as the kink, and which were anticipated in previous numerical studies as well as $`e`$-wobblers ($`eW`$) and $`g`$-wobblers ($`gW`$), periodic solutions which do not respect the spatial symmetry of the kink. In the SG case, it is the $`eW`$’s that respect the spatial symmetry while the $`gW`$’s do not.
(2) We have used the methods of scattering theory and ( Hamiltonian dispersive) normal forms to study the final stages of pinning of a kink to a particular lattice site. For large values of the discreteness parameter (“near” the continuum regime), the asymptotic state is a static kink. In contrast, for small values of the discreteness parameter, the above time periodic states can be attracting orbits. This is in sharp contrast to the behavior of solutions for the corresponding continuum equation in which a (possibly moving) kink is the attractor. Our analysis makes clear the relation of broken (Lorentz) invariance to these contrasting dynamics. Our results give a systematic clarification of the physicist’s heuristic picture of the dynamics of the center of mass of the kink as the effective (radiation-) damped motion of a massive particle in the Peierls-Nabarro potential. The approach to asymptotic state is via a slow damped (periodic or quasiperiodic) oscillation. The damping of the oscillation can indeed be very, very slow, and the lifetime is deducible from the normal form.
The methods we use are very systematic and general, and apply to many situations where the dynamical system can be viewed as the interaction of two subsystems: a finite dimensional subsystem, here governing the kink and its internal oscillations, and an infinite dimensional dynamical system, here governing dispersive radiation. In particular, one problem of related interest in which these methods can be readily generalized is the behavior of pulse-like breathing modes in the discrete non-linear Schrödinger equation $`i\psi _{i,t}=k\delta ^2\psi _i|\psi _i|^2\psi _i`$. Using the monochromatic gauge symmetry of this equation, we can convert the time periodic problem of the breathing solutions into a static problem by looking for solutions of the form $`\psi _i=\mathrm{exp}(\sqrt{1}\omega t)u_i`$ and studying their stability in the frame rotating with the same frequency $`\omega `$ (the so-called rotating wave approximation). The excitation and nonlinear Lyapunov stability of such states (see for example footnote 6 in section 2) is established in This notion of stability, being defined in terms of conserved integrals, is insensitive to the radiative behavior and the asymptotic approach to a ground state. For a recent account of the (numerical) construction and stability of such modes, see for instance . Then, once the problem has been posed on the rotating wave frame, it has become a static problem amenable to the same techniques for the study of dispersive waves and the radiation losses of the pulse-like coherent structure as the kink-like structures studies in this work. Exponentially small effects in the solutions for this class of systems arising for problems with a perturbed form of nonlinearity have been explored using perturbation theory techniques . These techniques are applicable to a restricted regime of parameter space because of the adiabatic approximation they entail. However, to capture, primarily, radiation effects caused by the discretization in the dynamics of continuum-like coherent structures our technique is obviously most appropriate (and clearly by no means restricted in its applicability) as illustrated by the analysis given above. What’s more, the phenomena present in such an analysis are in general of the power-law type rather than exponential.
Open questions: There are many related open problems of which we now mention several.
* Although the normal form analysis anticipates the existence and stability properties of various time-periodic solutions, we only have a rigorous theory concerning their existence. Linear stability theory would require an infinite dimensional Floquet analysis, and nonlinear stability theory would require the analogue of our normal form in a neighborhood of a periodic, rather than static, solution. Numerically, the asymptotic stability of the time periodic structures can be concluded from their persistence for long time integrations (to the extent of our computational power). The periodic structure construction and stability can also be performed using a limit-cycle type of technique similar to the one considered in . The rigorous stability theory however is still a challenging open mathematical problem.
* A rigorous infinite time theory, as in , for discrete nonlinear wave equations with kinks is a challenging open question. The problem is challenging even in the continuum case, where the kink is the sole attractor. In this case, the normal form yields the correct asymptotic rate of decay . However, due to the strength of the interaction (in space dimension one) of the dispersive and bound state components, the deviation of the full solution, $`u(t)`$, from the kink plus its internal mode modulations is not asymptotically free, i.e. not a free dispersive solution of the linear Klein-Gordon equation. Rather, it requires a solution-dependent logarithmic in time corrected phase , characteristic of long range scattering problems.
* Of interest would be more detailed long time simulations which would further elucidate the large time dynamics and assist in mapping out basins of attraction for various periodic states.
* Our study gives a detailed picture of the final stages of the pinning of a kink to a particular lattice site. However, the early stages of a kink plus large perturbation propagating in a lattice have some of the same features of the regime we consider. Namely , as the kink’s center of mass moves under the influence of the Peierls-Nabarro potential, it executes repeated acceleration and slowing with an approximate PN frequency. This oscillation appears to resonate with continuous radiation modes, resulting in an emission of energy, deceleration of the kink and eventual capture of the kink by a particular lattice site. It would be of interest to extend and apply the ideas of the current paper to this problem.
* In this paper we have analyzed the problem of radiative effects of discreteness for the $`2\pi `$ kinks within the P-N barrier (for a recent review on the effects of discreteness as well as the physical applications of a number of models similar to the ones considered here see ). It is well-known however that in the SG lattices multiple kinks ($`2n\pi `$ kinks in general) can also exist. For these kinks first mentioned in the static stability picture was analyzed in . There are a number of open problems regarding the behavior of such modes:
1. Their motion prior to pinning is only very slightly effected by radiative phenomena . This phenomenon hasn’t been accounted for satisfactorily to the best of our knowledge to date. In fact, this forms a part of a more general problem concerning the motion of coherent structures in lattices. It is known that under certain conditions and for potential different than the ones considered here there exist lattice systems with travelling wave solutions. Hence, it would be very interesting to elucidate the general conditions under which such lattice systems support travelling wave solutions (or equivalently under which conditions solutions to the advance-delay equations of the travelling wave frame exist).
2. When trapped (eventually) in the PN barrier these moving $`4\pi `$ kinks also radiate in a way that can be captured by the analysis of this paper. The problem there will have more (as shown in ) internal modes bifurcating from the continuum and hence will be more complicated but it can still be analyzed in the way presented above.
Acknowledgements
Thanks are expressed to N.J. Balmforth, A.R. Bishop, J.L. Lebowitz, C.K.R.T. Jones, Y.G. Kevrekidis, H. Segur and A. Soffer for stimulating discussions and for suggesting some important references. P.G.K. gratefully acknowledges assistanship support from the Computational Chemodynamics Laboratory of Rutgers University, fellowship support from the “Alexander S. Onasis” Public Benefit Foundation and also partial support from DIMACS. |
warning/0003/cs0003033.html | ar5iv | text | # Smodels: A System for Answer Set ProgrammingThe work has been funded by the Academy of Finland (Project 43963) and by Helsinki Graduate School in Computer Science and Engineering.
## General Information
The Smodels system is written in C++ and the source code, test cases and documentation are available at http://www.tcs.hut.fi/Software/smodels/. In order to compile the system a C++ compiler is needed as well as other standard tools such as make, tar, bison, and perl. The system has been developed under Linux and should work as is on any platform having the appropriate GNU tools installed. It has been used on a wide range of hardware (PC/SPARC/Alpha) mostly running Unix. The total number of lines of codes is about 20000.
## Description of the System
The Smodels system implements the stable model semantics for normal logic programs extended by built-in functions as well as cardinality and weight constraints for domain-restricted programs. In this section we briefly discuss the syntax, implementation techniques and use of the system. More information can be found at the home page http://www.tcs.hut.fi/Software/smodels/.
As input the Smodels system takes logic program rules basically in Prolog style syntax:
```
ancestor(X,Y) :- ancestor(X,Z),
parent(Z,Y),
person(X).
ancestor(X,Y) :- parent(X,Y).
son(X,Y) :- parent(Y,X), male(X).
daughter(X,Y) :- parent(Y,X), female(X).
person(X) :- male(X).
person(X) :- female(X).
parent(jack, jill). parent(joan, jack).
male(jack). female(jill). female(joan).
```
However, in order to support efficient implementation techniques and extensions the programs are required to be *domain-restricted* where the idea is the following. No proper function symbols are allowed (but we do allow built-in functions) and the predicate symbols in the program are divided into two classes, *domain predicates* and *non-domain predicates*. Domain predicates are predicates that are defined non-recursively. In the program above all predicates except `ancestor` are domain predicates. The predicate `ancestor` is not a domain predicate because it depends recursively on itself.
The main intuition of domain predicates is that they are used to define the set of terms over which the variables range in each rule of a program $`P`$. All rules of $`P`$ have to be domain-restricted in the sense that every variable in a rule must appear in a domain predicate which appears positively in the rule body. For instance, in the first rule of the program above all variables appear in domain predicates `parent` and `person` in the body of the rule.
In addition to normal logic program rules, Smodels supports rules with cardinality and weight constraints. The idea is that, e.g., a cardinality constraint
$$\mathtt{1\; \{\; a,b,not\; c\; \}\; 2}$$
holds in a stable model if at the least 1 but at most 2 of the literals in the constraint are satisfied in the model and a weight constraint
$$\mathtt{10\; [\; a=10,\; b=10,\; not\; c=10\; ]\; 20}$$
holds if the sum of weights of the literals satisfied in the model is between 10 and 20 (inclusive). With built-in functions for integer arithmetic (included in the system), these kinds of rules allow compact and fairly straightforward encodings of many interesting problems. For example, the N queens problem can be captured using rules as a program `queens.lp` as follows:
```
1 { q(X,Y):d(X) } 1 :- d(Y).
1 { q(X,Y):d(Y) } 1 :- d(X).
:- d(X), d(Y), d(X1), d(Y1),
q(X,Y), q(X1,Y1),
X != X1, Y != Y1,
abs(X - X1) == abs(Y - Y1).
d(1..n).
```
where `d(1..n)` is a domain predicate giving the dimension of the board (from 1 to integer $`n`$ which can be specified during run time). The first rule says that for each row $`y`$, a stable model contains exactly one atom $`q(x,y)`$ where for $`x`$, $`d(x)`$ holds and similarly in the second rule for all columns. The third rule is an integrity constraint saying that there cannot be two queens on the same diagonal. Now each stable model corresponds to a legal configuration of $`n`$ queens on a $`n\times n`$ board, i.e., $`q(x,y)`$ is in a stable model iff $`(x,y)`$ is a legal position for a queen.
Stable models of a domain-restricted logic program with variables are computed in three stages. First, the program is transformed into a ground program without variables. Second, the rules of the ground program are translated into primitive rules, and third, a stable model is computed using a Davis-Putnam like procedure (?). The first two stages have been implemented in the program lparse, which functions as a front end to smodels which in turn implements the third stage.
In the first stage lparse automatically determines the domain predicates and then using database techniques evaluates the domain predicates and creates a ground program which has exactly the same stable models as the original program with variables. Then the rules are compiled into primitive rules (?).
The smodels procedure is a Davis-Putnam like backtracking search procedure that finds the stable models of a set of primitive rules by assigning truth values to the atoms of the program. Moreover, it uses the properties of the stable model semantics to infer and propagate additional truth values. Since the procedure is in effect traversing a binary search tree, the number of nodes in the search space is in the worst case on the order of $`2^n`$, where $`n`$ can be taken to be the number of atoms that appear in a constraint in a head of a rule or that appear as a negative literal in a recursive loop of the program.
Hence, in order to compute stable models, one uses the two programs lparse, which translates logic programs into an internal format, and smodels, which computes the models, see Figure 1.
For instance, a solution to the 8 queens problem given the program `queens.lp` above, would be typically computed by a command line:
```
lparse -c n=8 -d none queens.lp | smodels
```
where `-c n=8` option instructs to use the value 8 for the constant $`n`$ and `-d none` option instructs to remove all domain predicates from the rules as soon as they have been evaluated. The command line produces output:
```
Answer: 1
Stable Model: q(4,1) q(2,2) q(7,3) q(5,4)
q(1,5) q(8,6) q(6,7) q(3,8)
```
### A Further Example
The graph coloring problem may be encoded to a program `ncolor.lp` using the following two Smodels rules:
```
1 { col(N, C) : color(C) } 1 :- node(N).
:- col(X, C), col(Y,C), edge(X,Y), color(C).
```
Here the predicate `col(N,C)` denotes that the color of the node `N` is `C`. The first rule states that each node has exactly one color and the second rule states that two adjacent nodes may not have the same color. Suppose we have a fully connected three node graph and we want to find all 3-colorings of it where the first node is colored red. We can encode the problem instance to a program `graph.lp`:
```
node(a ; b; c).
edge(a,b). edge(a,c). edge(b,c).
color(red ; green ; blue).
compute { col(a, red) }.
```
The first two lines define the graph and the third line defines the three colors. The `compute` statement tells Smodels that we are interested only in those models where `col(a, red)` is true, i.e., the node `a` has color `red`. We can find all stable models that satisfy the compute statement with the command line:
```
lparse -d none ncolor.lp graph.lp | smodels 0
```
This input produces the output:
```
Answer: 1
Stable Model: col(c,blue) col(b,green)
col(a,red)
Answer: 2
Stable Model: col(c,green) col(b,blue)
col(a,red)
```
More information about the syntax and use of the system can be found in the lparse user’s manual at http://www.tcs.hut.fi/Software/smodels/lparse/ and about the implementation in (??).
## Applying the System
### Methodology
An interesting application methodology for Smodels is based on answer set programming (??) which has emerged as a viable approach to declarative logic-based knowledge representation. It is based on the stable model semantics of logic programs and can be seen as a novel form of constraint programming where constraints are expressed as rules. The underlying idea is to encode an application problem using logic program rules so that the solutions to the problem are captured by the stable models of the rules. The solution of the N queens problem in the previous section illustrates nicely main ideas of answer set programming.
### Specifics
It is important that the system is based on an implementation-independent declarative semantics. This makes it much easier to develop applications because one does not have to worry too much about the internal implementation-specific aspects of the system. Hence, the system is relatively easy to learn to use. On the other hand, declarative semantics provides much more flexibility in developing implementation approaches and in optimizing different parts of the implementation. We have taken advantage of this and developed methods which are substantially different from usual implementation methods for logic programming (Prolog) but still work efficiently in new kinds of applications where Prolog style systems are not appropriate.
Smodels implements the stable model semantics for range-restricted function free normal programs. It can also compute well-founded models for these programs. Smodels supports built-in functions, e.g., for integer arithmetic. Basic Smodels extends normal logic programs with cardinality and weight constraints (??). On top of this core engine more involved systems can be built. As an example, we have implemented total and partial stable model computation for disjunctive logic programs (?).
The semantics for logic programs with cardinality and weight constraints supported by the core engine is an interesting compromise: it is rather simple to learn, its complexity stays in NP in the ground case (like for propositional logic) but it seems strictly more expressive than propositional logic (or other standard constraint satisfaction formalisms).
With our work on Smodels we hope to demonstrate that nonmonotonic reasoning techniques are useful conceptual tools as well as bring computational advantages which can lead to new interesting applications and to developments of novel implementation techniques and tools.
### Users and Usability
The basic semantics of Smodels rules seems to be simple enough that it can be explained to a non-expert or a student in a relatively short amount of time sufficiently so that one can start using Smodels.
Smodels has been employed in a number of areas including planning (???), model checking (?), reachability analysis (??), product configuration (??), dynamic constraint satisfaction (?), feature interaction (?), and logical cryptanalysis (?).
In order to make Smodels more flexible an API has been added to lparse and smodels. Hence, they can be used through the API and embedded into a C/C++ program as libraries. Furthermore, it is possible to define new built-in functions for the front-end lparse in order to accommodate new applications. These new functions are written in C/C++ and they are dynamically linked to lparse when needed.
As a further usability feature, lparse performs simple analysis of the program and warns about constructs that are often erroneous. For example, lparse detects if a variable is accidentally mistyped as a constant or vice versa. These warnings can be enabled with command line options.
## Evaluating the System
### Benchmarks
We have compiled quite a large collection of families of benchmarks that we use to evaluate new developments and improvements in our system and which we can employ to compare our system to other competing ones and which can be used also by others (see, http://www.tcs.hut.fi/Software/smodels/tests/). Similar testing methodology (e.g., generating test cases from graph problems) has also been used by the groups in Kentucky (DeReS/TheoryBase system (?)) and in Vienna (dlv system (?)). These kinds of tests can be used for measuring the base level performance and to compare different systems. However, it is unclear what is a portable way of representing such benchmarks. We use Prolog style syntax which is quite generally accepted but it seems that each system has its specialties.
### Comparison
Smodels can already compete with special purpose systems and we have even cases where it outperforms commercial top edge tools, e.g., in a verification application it has performance better than that of one of the most efficient commercial integer programming tools (CPLEX) (?).
### Problem Size
We believe that Smodels is no longer a mere prototype but it can handle realistic size problems. We have applications where programs with hundreds of thousands of non-trivial ground rules are treated efficiently, see e.g. (??). |
warning/0003/hep-ph0003244.html | ar5iv | text | # References
. 24 March 2000
Anomalous $`t`$-dependence in diffractive electroproduction of 2S radially excited light vector mesons at HERA
J. Nemchik
Institute of Experimental Physics, Slovak Academy of Sciences,
Watsonova 47, 04353 Košice, Slovakia
Abstract
Within the color dipole gBFKL dynamics applied to the diffraction slope, we predict an anomalous $`t`$ dependence of the differential cross section as a function of energy and $`Q^2`$ for production of radially excited $`V^{}(2S)`$ light vector mesons in contradiction with a well known standard monotonous $`t`$\- behaviour for $`V(1S)`$ mesons. The origin of this phenomenon is based on the interplay of the nodal structure of $`V^{}(2S)`$ radial wave function with the energy and dipole size dependence of the color dipole cross section and of the diffraction slope. We present how a different position of the node in $`V^{}(2S)`$ wave function leads to a different form of anomalous $`t`$\- behaviour of the differential cross section and discuss a possibility how to determine this position from the low energy and HERA data.
The main goal of this paper is a demonstration of further salient features of the node effect coming from the nodal structure of radial wave function for $`V^{}(2S)`$ vector mesons in conjunction with the gBFKL phenomenology of the diffraction slope leading to anomalous $`t`$ dependence of the differential cross section for $`V^{}(2S)`$ production in contrast with the standard monotonous $`t`$ behaviour of $`d\sigma (\gamma ^{}V)/dt`$ for $`V(1S)`$ production.
Diffractive photo- and electroproduction of ground state $`V(1S)`$ and radially excited $`V^{}(2S)`$ vector mesons,
$$\gamma ^{}pV(V^{})pV=\rho ,\mathrm{\Phi },\omega ,J/\mathrm{\Psi },\mathrm{{\rm Y}}\mathrm{}(V^{}=\rho ^{},\mathrm{\Phi }^{},\omega ^{},\mathrm{\Psi }^{},\mathrm{{\rm Y}}^{}\mathrm{}),$$
(1)
at high c.m.s. energy $`W=\sqrt{s}`$ intensively studied by the recent experiments at HERA represents one of a main source for a further development of the pomeron physics. The pomeron exchange in diffractive leptoproduction of vector mesons at high energies has been intensively studied within the framework of perturbative QCD (pQCD).
The standard approach to the pQCD is based on the BFKL equation , which represents the integral equation for the leading-log$`s`$ (LLs) evolution of the gluon distribution, formulated in the scaling approximation of the infinite gluon correlation radius, $`R_c\mathrm{}`$, (massless gluons) and of the fixed running coupling, $`\alpha _S=const`$. Later, however, a novel $`s`$-channel approach to the $`LLs`$ BFKL equation (running gBFKL approach) has been developed in terms of the color dipole cross section $`\sigma (\xi ,r)`$ (hereafter $`r`$ is the transverse size of the color dipole, $`\xi =log(\frac{W^2+Q^2}{m_V^2+Q^2})`$ is the rapidity variable) and incorporates consistently the asymptotic freedom (AF) (i.e. the running QCD coupling $`\alpha _S(r)`$) and the finite propagation radius $`R_c`$ of perturbative gluons. The freezing of $`\alpha _S(r)`$, $`\alpha _S(r)\alpha _S^{fr}`$, and the gluon correlation radius $`R_c`$ represents the nonperturbative parameters, which describe the transition from the soft (nonperturbative, infrared) to the hard (perturbative) region.
The details of the gBFKL phenomenology of diffractive electroproduction of light vector mesons are presented in the paper . The color dipole phenomenology of the diffraction slope for photo- and electroproduction of heavy vector mesons has been developed in the paper . The analysis of the diffractive production of light and heavy vector mesons at $`t=0`$ within the gBFKL phenomenology shows that the $`1S`$ vector meson production amplitude probes the color dipole cross section at the dipole size $`rr_S`$ (scanning phenomenon ), where the scanning radius can be expressed through the scale parameter $`A`$
$$r_S\frac{A}{\sqrt{m_V^2+Q^2}},$$
(2)
where $`Q^2`$ is the photon virtuality, $`m_V`$ is the vector meson mass and $`A6`$. Thus, changing $`Q^2`$ and the mass of the produced vector meson, one can probe the dipole cross section $`\sigma (\xi ,r)`$, and the dipole diffraction slope $`B(\xi ,r)`$, and measure so the effective intercept $`\mathrm{\Delta }_{eff}(\xi ,r)=\mathrm{log}\sigma (\xi ,r)/\xi `$ and the local Regge slope $`\alpha _{eff}^{}(\xi ,r)=\frac{1}{2}B(\xi ,r)/\xi `$ in a very broad range of the dipole sizes, $`r`$. This fact allows also to study the transition between the perturbative (hard) and nonperturbative (soft) regimes.
The experimental investigation of the electroproduction of the radially excited ($`2S`$) vector mesons can extend an additional information on the dipole cross section and on the dipole diffraction slope. The presence of the node in the $`2S`$ radial wave function leads to a strong cancellation of the dipole size contributions to the production amplitude from the region above and below the node position, $`r_n`$, in the $`2S`$ radial wave function . For this reason, the amplitudes of the electroproduction of the $`1S`$ and $`2S`$ vector mesons probe $`\sigma (\xi ,r)`$ and $`B(\xi ,r)`$ in a different way. The onset of strong node effect has been demonstrated in Ref. in electroproduction of radially excited light vector mesons leading to an anomalous $`Q^2`$ and energy dependence of the production cross section. The node effect is much weaker for the electroproduction of $`2S`$ heavy vector mesons but still leads to a slightly different $`Q^2`$ and energy dependence of the production cross section for $`\mathrm{\Psi }^{}`$ vs. $`J/\mathrm{\Psi }`$ and to a nonmonotonic $`Q^2`$ dependence of the diffraction slope at small $`Q^2\text{ }<5`$ GeV<sup>2</sup> for $`\mathrm{\Psi }^{}`$ production . Only for $`\mathrm{{\rm Y}}^{}`$ production, the node effect is negligible small and gives approximately the same $`Q^2`$ and energy behaviour of the production cross section and practically the same diffraction slope at $`t=0`$ for $`\mathrm{{\rm Y}}`$ and $`\mathrm{{\rm Y}}^{}`$ production . Therefore, it is very important to explore farther the salient features of the node effect with conjunction with the emerging gBFKL phenomenology of the diffraction slope especially in production of $`V^{}(2S)`$ light vector mesons where the node effect is expected to be very strong.
There are two main reasons which affect the cancellation pattern in the diffraction slope for $`2S`$ state. The first reason is connected with the $`Q^2`$ behaviour of the scanning radius $`r_S`$ (see (2)); for the electroproduction of $`V^{}(2S)`$ light vector mesons at moderate $`Q^2`$ when the scanning radius $`r_S`$ is close to $`r_n`$, due to $`r^2`$ behaviour of $`B(\xi ,r)`$ , even a slight variation of $`r_S`$ with $`Q^2`$ strongly changes the cancellation pattern and leads to an anomalous $`Q^2`$ dependence of the forward diffraction slope, $`B(t=0)`$ . The second reason is due to different energy dependence of $`\sigma (\xi ,r)`$ at different dipole sizes $`r`$ coming from the gBFKL dynamics leading also to an anomalous energy dependence of $`B(t=0)`$ for the $`V^{}(2S)`$ production. This nonmonotonous energy and $`Q^2`$ dependence of the diffraction slope for production of light vector mesons will be detaily studied elsewhere .
The effects mentioned above are sensitive to the form of the dipole cross section. In Ref. we presented the first direct determination of the color dipole cross section from the data on the photo- and electroproduction of $`V(1S)`$ vector mesons. So extracted dipole cross section is in a good agreement with the dipole cross section obtained from gBFKL analysis . This fact confirms a very reasonable choice of the nonperturbative component of the dipole cross section corresponding to a soft nonperturbative mechanism contribution to the scattering amplitude.
In the present paper we concentrate on the production of $`2S`$ radially excited light vector mesons, where the node in the radial wave function in conjunction with the subasymptotic energy dependence of $`B(\xi ,r)`$ leads to a strikingly different $`t`$ dependence of the differential cross section at different energies and $`Q^2`$ for the production of $`V^{}(2S)`$ vs. $`V(1S)`$ vector mesons. As was mentioned above, due to a large value of the scale parameter in (2), the large-distance contributions to the production amplitude from the semiperturbative and nonperturbative region of color dipoles, $`r\text{ }>R_c`$, becomes substantial especially for light vector mesons. Only the virtual $`\rho ^0`$ and $`\varphi ^0`$ photoproduction at $`Q^2\text{ }>100`$ GeV<sup>2</sup> can be treated as a purely perturbative process, when the production amplitude is dominantly contributed from the perturbative region, $`r\text{ }<R_c`$.
Thus, in this paper we present the $`Q^2`$ and energy dependence of the $`t`$\- behaviour of the differential cross section for electroproduction of the ground state and radially excited (2S) light vector meson. and study how the position of the node in the radial wave function for (2S) vector mesons can be extracted from the data. We present an exact prescription how the experimental measurement of the $`t`$ dependent differential cross section for $`V^{}(2S)`$ production could distinguish between the undercompensation and overcompensation scenarios of the $`2S`$ production amplitude (see below). The explicit form of that $`t`$\- behaviour is connected with the position of the node in radial wave function for $`V^{}(2S)`$ vector mesons.
In the mixed $`(𝐫,z)`$ representation, the high energy meson is considered as a system of color dipole described by the distribution of the transverse separation $`𝐫`$ of the quark and antiquark given by the $`q\overline{q}`$ wave function, $`\mathrm{\Psi }(𝐫,z)`$, where $`z`$ is the fraction of meson’s lightcone momentum carried by a quark. The Fock state expansion for the relativistic meson starts with the $`q\overline{q}`$ state and the higher Fock states $`q\overline{q}g\mathrm{}`$ become very important at high energy $`\nu `$. The interaction of the relativistic color dipole of the dipole moment, $`𝐫`$, with the target nucleon is quantified by the energy dependent color dipole cross section, $`\sigma (\xi ,r)`$, satisfying the gBFKL equation for the energy evolution. This reflects the fact that in the leading-log $`\frac{1}{x}`$ approximation the effect of higher Fock states can be reabsorbed into the energy dependence of $`\sigma (\xi ,r)`$. The dipole cross section is flavor independent and represents the universal function of $`r`$ which describes various diffractive processes in unified form. At high energy, when the transverse separation, $`𝐫`$, of the quark and antiquark is frozen during the interaction process, the scattering matrix describing the $`q\overline{q}`$-nucleon interaction becomes diagonal in the mixed $`(𝐫,z)`$-representation ($`z`$ is known also as the Sudakov light cone variable). This diagonalization property is held even when the dipole size, $`𝐫`$, is large, i.e. beyond the perturbative region of short distances. The detailed discussion about the space-time pattern of diffractive electroproduction of vector mesons is presented in .
Following an advantage of the $`(𝐫,z)`$-diagonalization of the $`q\overline{q}N`$ scattering matrix, the imaginary part of the production amplitude for the real (virtual) photoproduction of vector mesons with the momentum transfer $`𝐪`$ can be represented in the factorized form
$$\mathrm{Im}(\gamma ^{}V,\xi ,Q^2,𝐪)=V|\sigma (\xi ,r,z,𝐪)|\gamma ^{}=\underset{0}{\overset{1}{}}𝑑zd^2𝐫\sigma (\xi ,r,z,𝐪)\mathrm{\Psi }_V^{}(𝐫,z)\mathrm{\Psi }_\gamma ^{}(𝐫,z)$$
(3)
whose normalization is $`d\sigma /dt|_{t=0}=||^2/16\pi .`$ In Eq. (3), $`\mathrm{\Psi }_\gamma ^{}(𝐫,z)`$ and $`\mathrm{\Psi }_V(𝐫,z)`$ represent the probability amplitudes to find the color dipole of size, $`r`$, in the photon and quarkonium (vector meson), respectively. The color dipole distribution in (virtual) photons was derived in . $`\sigma (\xi ,r,z,𝐪)`$ is the scattering matrix for $`q\overline{q}N`$ interaction and represents the above mentioned color dipole cross section for $`𝐪=0`$. The color dipole cross section $`\sigma (\xi ,r)`$ depends only on the dipole size $`r`$. For small $`𝐪`$ considered in this paper, one can safely neglect the $`z`$-dependence of $`\sigma (\xi ,r,z,𝐪)`$ for light and heavy vector meson production and set $`z=\frac{1}{2}`$. This follows partially from the analysis within double gluon exchange approximation leading to a slow $`z`$ dependence of the dipole cross section.
The energy dependence of the dipole cross section is quantified in terms of the dimensionless rapidity, $`\xi =\mathrm{log}\frac{1}{x_{eff}}`$, and $`x_{eff}`$ is the effective value of the Bjorken variable
$$x_{eff}=\frac{Q^2+m_V^2}{Q^2+W^2}\frac{m_V^2+Q^2}{2\nu m_p},$$
(4)
where $`m_p`$ and $`m_V`$ is the proton mass and mass of vector meson, respectively. Hereafter, we will write the energy dependence of the dipole cross section in both variables, either in $`\xi `$ or in $`x_{eff}`$ whenever convenient.
The production amplitudes for the transversely (T) and the longitudinally (L) polarized vector mesons with the small momentum transfer, $`𝐪`$, can be written in more explicit form
$`\mathrm{Im}_T(x_{eff},Q^2,𝐪)={\displaystyle \frac{N_cC_V\sqrt{4\pi \alpha _{em}}}{(2\pi )^2}}`$
$`{\displaystyle }d^2𝐫\sigma (x_{eff},r,𝐪){\displaystyle _0^1}{\displaystyle \frac{dz}{z(1z)}}\{m_q^2K_0(\epsilon r)\varphi (r,z)[z^2+(1z)^2]\epsilon K_1(\epsilon r)_r\varphi (r,z)\}`$
$`={\displaystyle \frac{1}{(m_V^2+Q^2)^2}}{\displaystyle \frac{dr^2}{r^2}\frac{\sigma (x_{eff},r,𝐪)}{r^2}W_T(Q^2,r^2)}`$ (5)
$`\mathrm{Im}_L(x_{eff},Q^2,𝐪)={\displaystyle \frac{N_cC_V\sqrt{4\pi \alpha _{em}}}{(2\pi )^2}}{\displaystyle \frac{2\sqrt{Q^2}}{m_V}}`$
$`{\displaystyle }d^2𝐫\sigma (x_{eff},r,𝐪){\displaystyle _0^1}dz\{[m_q^2+z(1z)m_V^2]K_0(\epsilon r)\varphi (r,z)_r^2\varphi (r,z)\}`$
$`={\displaystyle \frac{1}{(m_V^2+Q^2)^2}}{\displaystyle \frac{2\sqrt{Q^2}}{m_V}}{\displaystyle \frac{dr^2}{r^2}\frac{\sigma (x_{eff},r,𝐪)}{r^2}W_L(Q^2,r^2)}`$ (6)
where
$$\epsilon ^2=m_q^2+z(1z)Q^2,$$
(7)
$`\alpha _{em}`$ is the fine structure constant, $`N_c=3`$ is the number of colors, $`C_V=\frac{1}{\sqrt{2}},\frac{1}{3\sqrt{2}},\frac{1}{3},\frac{2}{3},\frac{1}{3}`$ for $`\rho ^0,\omega ^0,\varphi ^0,J/\mathrm{\Psi },\mathrm{{\rm Y}}`$ production, respectively and $`K_{0,1}(x)`$ are the modified Bessel functions. The detailed discussion and parameterization of the lightcone radial wave function $`\varphi (r,z)`$ of the $`q\overline{q}`$ Fock state of the vector meson is given in .
The terms $`ϵK_1(ϵr)_r\varphi (𝐫,z)`$ for $`T`$ polarization and $`K_0(ϵr)_r^2\mathrm{\Phi }(𝐫,z)`$ for $`L`$ polarization in the integrands of (5) and (6) represent the relativistic corrections which become important at large $`Q^2`$ and for the production of light vector mesons. For the production of heavy quarkonia, the nonrelativistic approximation can be used with a rather high accuracy .
For small dipole size and $`𝐪=0`$, in the leading-log $`\frac{1}{x}`$ approximation, the dipole cross section can be related to the gluon structure function $`G(x,q^2)`$ of the target nucleon through
$$\sigma (x,r)=\frac{\pi ^2}{3}r^2\alpha _s(r)G(x,q^2),$$
(8)
where the gluon structure function enters at the factorization scale, $`q^2\frac{B}{r^2}`$ with the parameter $`B10`$ .
The weight functions, $`W_T(Q^2,r^2)`$ and $`W_L(Q^2,r^2)`$, introduced in (5) and (6) have a smooth $`Q^2`$ behaviour and are very convenient for the analysis of the scanning phenomenon. They are sharply peaked at $`rA_{T,L}/\sqrt{Q^2+m_V^2}`$. At small $`Q^2`$ the values of the scale parameter $`A_{T,L}`$ are close to $`A6`$, which follows from $`r_S=3/\epsilon `$ with the nonrelativistic choice $`z=\frac{1}{2}`$. In general, $`A_{T,L}6`$ and increases slowly with $`Q^2`$ . For production of light vector mesons the relativistic corrections play an important role especially at large $`Q^2m_V^2`$, and lead to $`Q^2`$ dependence of $`A_{L,T}`$ coming from the large-size asymmetric $`q\overline{q}`$ configurations: $`A_L(\rho ^0;Q^2=0)6.5,A_L(\rho ^0;Q^2=100\mathrm{GeV}^2)10,A_T(\rho ^0;Q^2=0)7,A_T(\rho ^0,Q^2=100\mathrm{GeV}^2)12`$ . Due to an extra factor $`z(1z)`$ in the integrand of (6) in comparison with (5), the contribution from asymmetric $`q\overline{q}`$\- configurations to the longitudinal meson production is considerably smaller.
The integrands in Eqs. (5) and (6) contain the dipole cross section, $`\sigma (\xi ,r,𝐪)`$. As was mentioned, due to a very slow onset of the pure perturbative region (see Eq. (2)), one can easily anticipate a contribution to the production amplitude coming from the semiperturbative and nonperturbative $`r\text{ }>R_c`$. Following the simplest assumption about an additive property of the perturbative and nonperturbative mechanism of interaction, we can represent the contribution of the bare pomeron exchange to $`\sigma (\xi ,r,𝐪)`$ as a sum of the perturbative and nonperturbative component<sup>1</sup><sup>1</sup>1 additive property of a such decomposition of the dipole cross section has been detaily discussed in
$$\sigma (\xi ,r,𝐪)=\sigma _{pt}(\xi ,r,𝐪)+\sigma _{npt}(\xi ,r,𝐪),$$
(9)
with the parameterization of both components at small $`𝐪`$
$$\sigma _{pt,npt}(\xi ,r,𝐪)=\sigma _{pt,npt}(\xi ,r,𝐪=0)\mathrm{exp}\left(\frac{1}{2}B_{pt,npt}(\xi ,r)𝐪^\mathrm{𝟐}\right).$$
(10)
Here $`\sigma _{pt,npt}(\xi ,r,𝐪=0)=\sigma _{pt,npt}(\xi ,r)`$ represent the contribution of the perturbative and nonperturbative mechanisms to the $`q\overline{q}`$-nucleon interaction cross section, respectively, $`B_{pt,npt}(\xi ,r)`$ are corresponding diffraction slopes.
A small real part of production amplitudes can be taken in the form
$$\mathrm{Re}(\xi ,r)=\frac{\pi }{2}\frac{}{\xi }\mathrm{Im}(\xi ,r).$$
(11)
and can be easily included in the production amplitudes (5),(6) using substitution
$$\sigma (x_{eff},r,𝐪)\left(1i\frac{\pi }{2}\frac{}{logx_{eff}}\right)\sigma (x_{eff},r)=\left[1i\alpha _V(x_{eff},r)\right]\sigma (x_{eff},r,𝐪)$$
(12)
The formalism for calculation of $`\sigma _{pt}(\xi ,r)`$ in the leading-log $`s`$ approximation was developed in . The nonperturbative contribution, $`\sigma _{npt}(\xi ,r)`$, to the dipole cross section was used in Refs. where we assume that this soft nonperturbative component of the pomeron is a simple Regge pole with the intercept, $`\mathrm{\Delta }_{npt}=0`$. The particular form together with assumption of the energy independent $`\sigma _{npt}(\xi =\xi _0,r)=\sigma _{npt}(r)`$ ($`\xi _0`$ corresponds to boundary condition for the gBFKL evolution, $`\xi _0=log(1/x_0)`$, $`x_0=0.03`$) allows one to successfully describe the proton structure function at very small $`Q^2`$, the real photoabsorption and diffractive real and virtual photoproduction of light and heavy vector mesons. A larger contribution of the nonperturbative pomeron exchange to $`\sigma _{tot}(\gamma p)`$ vs. $`\sigma _{tot}(\gamma ^{}p)`$ can, for example, explain a much slower rise with energy of the real photoabsorption cross section, $`\sigma _{tot}(\gamma p)`$, in comparison with $`F_2(x,Q^2)\sigma _{tot}(\gamma ^{}p)`$ observed at HERA . Besides, the reasonable form of this soft cross section, $`\sigma _{npt}(r)`$, was confirmed in the process of the first determination of the dipole cross section from the data on vector meson electroproduction . The so extracted dipole cross section is in a good agreement with the dipole cross section obtained from the gBFKL dynamics . Thus, this nonperturbative component of the pomeron exchange plays a dominant role at low NMC energies in the production of the light vector mesons, where the scanning radius, $`r_S`$ (2), is large. However, the perturbative component of the pomeron become more important with the rise of energy also in the nonperturbative region of the dipole sizes.
If one starts with the familiar impact-parameter representation for amplitude of elastic scattering of the color dipole
$$\mathrm{Im}(\xi ,r,\stackrel{}{q})=2d^2\stackrel{}{b}\mathrm{exp}(i\stackrel{}{q}\stackrel{}{b})\mathrm{\Gamma }(\xi ,\stackrel{}{r},\stackrel{}{b}),$$
(13)
then the diffraction slope $`B=2d\mathrm{log}\mathrm{Im}/dq^2|_{q=0}`$ equals
$$B(\xi ,r)=\frac{1}{2}\stackrel{}{b}^2=\lambda (\xi ,r)/\sigma (\xi ,r),$$
(14)
where
$$\lambda (\xi ,r)=d^2\stackrel{}{b}\stackrel{}{b}^2\mathrm{\Gamma }(\xi ,\stackrel{}{r,}\stackrel{}{b}).$$
(15)
The generalization of the color dipole factorization formula (3) to the diffraction slope of the reaction $`\gamma ^{}pVp`$ reads:
$$B(\gamma ^{}V,\xi ,Q^2)\mathrm{Im}(\gamma ^{}V,\xi ,Q^2,\stackrel{}{q}=0)=\underset{0}{\overset{1}{}}𝑑zd^2\stackrel{}{r}\lambda (\xi ,r)\mathrm{\Psi }_V^{}(r,z)\mathrm{\Psi }_\gamma ^{}(r,z).$$
(16)
The diffraction cone in the color dipole gBFKL approach for production of vector mesons has been detaily studied in . Here we only present the salient feature of the color diffraction slope reflecting the presence of the geometrical contribution from beam dipole - $`r^2/8`$ and the contribution from the target proton size - $`R_N^2/3`$:
$$B(\xi ,r)=\frac{1}{8}r^2+\frac{1}{3}R_N^2+2\alpha _{𝐈𝐏}^{}(\xi \xi _0)+𝒪(R_c^2),$$
(17)
where $`R_N`$ is the radius of the proton. For electroproduction of light vector mesons the scanning radius is larger than the correlation one $`r\text{ }>R_c`$ even for $`Q^2\text{ }<50`$ GeV<sup>2</sup> and one recovers a sort of additive quark model, in which the uncorrelated gluonic clouds build up around the beam and target quarks and antiquarks and the term $`2\alpha _{𝐈𝐏}^{}(\xi \xi _0)`$ describe the familiar Regge growth of diffraction slope for the quark-quark scattering. The geometrical contribution to the diffraction slope from the target proton size, $`\frac{1}{3}R_N^2`$, persists for all the dipole sizes, $`r\text{ }>R_c`$ and $`r\text{ }<R_c`$. The last term in (17) is also associated with the proton size and is negligibly small.
The soft pomeron and diffractive scattering of large color dipole has been detaily studied in the paper . Here we assume the conventional Regge rise of the diffraction slope for the soft pomeron ,
$$B_{npt}(\xi ,r)=\mathrm{\Delta }B_d(r)+\mathrm{\Delta }B_N+2\alpha _{npt}^{^{}}(\xi \xi _0),$$
(18)
where $`\mathrm{\Delta }B_d(r)`$ and $`\mathrm{\Delta }B_N`$ stand for the contribution from the beam dipole and target nucleon size. As a guidance we take the experimental data on the pion-nucleon scattering , which suggest $`\alpha _{npt}^{}=0.15`$ GeV<sup>-2</sup>. In (18) the proton size contribution is
$$\mathrm{\Delta }B_N=\frac{1}{3}R_N^2,$$
(19)
and the beam dipole contribution has been proposed to have a form
$$B_d(r)=\frac{r^2}{8}\frac{r^2+aR_N^2}{3r^2+aR_N^2},$$
(20)
where $`a`$ is a phenomenological parameter, $`a1`$. We take $`\mathrm{\Delta }B_N=4.8\mathrm{GeV}^2`$. Then the pion-nucleon diffraction slope is reproduced with reasonable values of the parameter $`a`$ in the formula (20): $`a=0.9`$ for $`\alpha _{npt}^{}=0.15`$ GeV<sup>-2</sup> .
Using the expressions (5) and (6) for the $`T`$ and $`L`$ production amplitudes in conjunction with Eqs. (9) and (10), we can calculate the differential cross section of vector meson electroproduction as a function of $`t`$.
Following the simple geometrical properties of the gBFKL diffraction slope, $`B(\xi ,r)`$, (see Eq. (17) and ), one can express its energy dependence through the energy dependent effective Regge slope, $`\alpha _{eff}^{}(\xi ,r)`$
$$B_{pt}(\xi ,r)\frac{1}{3}<R_N^2>+\frac{1}{8}r^2+2\alpha _{eff}^{}(\xi ,r)(\xi \xi _0).$$
(21)
The effective Regge slope, $`\alpha _{eff}^{}(\xi ,r)`$, varies with energy differently at different size of the color dipole ; at fixed scanning radius and/or $`Q^2+m_V^2`$, it decreases with energy. At fixed rapidity $`\xi `$ and/or $`x_{eff}`$ (4), $`\alpha _{eff}^{}(\xi ,r)`$ rises with $`r\text{ }<1.5`$ fm. At fixed energy, it is a flat function of the scanning radius. At the asymptotically large $`\xi `$ ($`W`$), $`\alpha _{eff}^{}(\xi ,r)\alpha _{𝐈𝐏}^{}=0.072`$ GeV<sup>-2</sup>. At the lower and HERA energies, the subasymptotic $`\alpha _{eff}^{}(\xi ,r)(0.150.20)`$ GeV<sup>-2</sup> and is very close to $`\alpha _{soft}^{}`$ known from the Regge phenomenology of soft scattering. It means, that the gBKFL dynamics predicts a substantial rise with the energy and dipole size, $`r`$, of the diffraction slope, $`B(\xi ,r)`$, in accordance with the energy and dipole size dependence of the effective Regge slope, $`\alpha _{eff}^{}(\xi ,r)`$ and due to a presence of the geometrical component, $`r^2`$, in (17) and (18).
Now we will concentrate on the production of radially excited $`2S`$ light vector mesons and will study the differential cross section $`d\sigma /dt`$ as a function of $`t`$. The most important feature of the production of $`V^{}(2S)`$ vector mesons is the node effect \- the $`Q^2`$ and energy dependent cancellations from the soft (large size) and hard (small size) contributions, i.e. from the region above and below the node position, $`r_n`$, to the $`V^{}(2S)`$ production amplitude. The strong $`Q^2`$ dependence of these cancellations comes from the scanning phenomenon (2) when the scanning radius $`r_S`$ for some value of $`Q^2`$ is close to $`r_nR_V`$ ($`R_V`$ is the vector meson radius). The energy dependence of the node effect is due to the different energy dependence of the dipole cross section at small ($`r<R_V`$) and large ($`r>R_V`$) dipole sizes. The strong node effect in production of radially excited light vector mesons leading to an anomalous $`Q^2`$ and energy dependence of the production cross section was demonstrated in Ref. <sup>2</sup><sup>2</sup>2 Manifestations of the node effect in electroproduction on nuclei were discussed earlier, see and Note, that the predictive power is weak and is strongly model dependent in the region of $`Q^2`$ and energy where the node effect becomes exact.
There are several reasons to expect that, for the production of $`2S`$ light vector mesons, the node effect depends on the polarization of the virtual photon and of the produced vector meson . First, the wave functions of $`T`$ and $`L`$ polarized (virtual) photon are different. Second, different regions of $`z`$ contribute to the $`_T`$ and $`_L`$. Third, different scanning radii for production of $`T`$ and $`L`$ polarized vector mesons and different energy dependence of $`\sigma (\xi ,r)`$ at these scanning radii lead to a different $`Q^2`$ and energy dependence of the node effect in production of $`T`$ and $`L`$ polarized $`V^{}(2S)`$ vector mesons. Not so for the nonrelativistic limit of heavy quarkonia where the node effect is very weak and is approximately polarization independent. There is a weak polarization dependence of the node effect which is the most marginal for $`\mathrm{\Psi }^{}`$ production and this weak node effect still leads to a nonmonotonic $`Q^2`$ dependence of the diffraction slope. For $`\mathrm{{\rm Y}}^{}`$ production the node effect is negligibly small and is polarization independent with very high accuracy.
There are two possible scenarios for the node effect which can occur in the $`2S`$ production amplitude; the undercompensation and the overcompensation scenario . In the undercompensation case, the $`2S`$ production amplitude $`V^{}(2S)|\sigma (\xi ,r)|\gamma ^{}`$ is dominated by the positive contribution coming from small dipole sizes, $`r\text{ }<r_n`$ ($`r_n`$ is the node position), and the $`V(1S)`$ and $`V^{}(2S)`$ photoproduction amplitudes have the same sign. This scenario corresponds namely to the production of $`2S`$ heavy vector mesons, $`\mathrm{\Psi }^{}(2S)`$ and $`\mathrm{{\rm Y}}^{}(2S)`$. In the overcompensation case, the $`2S`$ production amplitude $`V^{}(2S)|\sigma (\xi ,r)|\gamma ^{}`$ is dominated by the negative contribution coming from large dipole sizes, $`r\text{ }>r_n`$, and the $`V(1S)`$ and $`V^{}(2S)`$ photoproduction amplitudes have the opposite sign.
The anomalous properties of the diffraction slope come from the expression (16) and will be presented elsewhere . The matrix element on l.h.s of (16) represents the well known production amplitude $`V(V^{})|\sigma (\xi ,r)|\gamma ^{}`$. As was mentioned, the $`1S`$ production amplitude is dominated by contribution from dipole size corresponding to the scanning radius $`r_S3/ϵ`$ (2) with the scale parameter $`A6`$ at $`Q^2=0`$ slightly dependent on $`Q^2`$ . However, due to $`r^2`$ behaviour of the slope parameter (see (17) and (18)), the integrand of the matrix element on the r.s.h of Eq. (16), $`V(1S)|\sigma (\xi ,r)B(\xi ,r)|\gamma ^{}`$, is $`r^5\mathrm{exp}(ϵr)`$ and is peaked by $`rr_B=5/ϵ=5/3r_S`$.
The node of the radial wave function of the $`2S`$ states leads to peculiarities in $`t`$ dependence of the differential cross section. Following the simple geometrical properties of the diffraction slope (17), (18), because of $`r^2`$ behaviour, the large size negative contribution to the production amplitude from the region above the node position corresponds to larger value of the diffraction slope than small size contribution from the region below the node position. It means, that the negative contribution to the $`V^{}(2S)`$ production amplitude coming from the region above the node position, $`r\text{ }>r_n`$, has a steeper $`t`$ dependence than the positive contribution coming from the small size dipoles, $`r\text{ }>r_n`$. It can be expressed in a somewhat demonstrative form as a $`t`$\- dependent production amplitude:
$$(t)=\alpha \mathrm{exp}(\frac{1}{2}B_1t)\beta \mathrm{exp}(\frac{1}{2}B_2t),$$
(22)
where $`\alpha `$ and $`\beta `$ represent the contribution to the matrix element from the region below and above the node position, respectively. In (22) $`B_1`$ and $`B_2`$ are the effective diffraction slopes, which correspond to integration over dipole size $`r`$ from 0 to the position of the node $`r_n`$ and above the node position. Thus, the inequality $`\alpha >\beta `$ corresponds to the undercompensation whereas $`\alpha <\beta `$ to the overcompensation regime. The destructive interference of these two amplitudes results in a decrease of the effective diffraction slope for the $`V^{}(2S)`$ meson production for small $`t`$ in contrary with the familiar increase for the $`V(1S)`$ meson production. Such a situation is shown in Fig. 1, where we present the model predictions for the differential cross section as a function of $`t`$ for production of $`V(1S)`$ and $`V^{}(2S)`$ mesons at different c.m.s. energies $`W`$ and at $`Q^2=0`$. Real photoproduction measures the purely transverse cross section. As was mentioned in the paper , using our wave functions, at $`W\text{ }<150`$ GeV for $`\rho ^{}(2S)`$ production and at $`W\text{ }<20`$ GeV for $`\varphi ^{}(2S)`$ production, the forward production amplitude (3) is in undercompensation regime whereas the matrix element $`<V^{}(2S)|\sigma (\xi ,r)B(\xi ,r)|\gamma >`$ on r.h.s. of Eq. (16) is in the overcompensation regime. As the result we predict the negative valued diffraction slope at $`t=0`$ and $`Q^2=0`$ For this reason and due to destructive interference of two contributions to the production amplitude (22) with different $`t`$ dependencies, the differential cross section firstly rises with $`t`$, flattens at $`t(0.10.2)`$ GeV<sup>2</sup> having a maximum. At large $`t`$, the node effect is weak in $`t`$\- dependent production amplitude because of a steeper $`t`$ dependence from the large size dipoles and the differential cross section falls down following the differential cross section for $`V(1S)`$ production. The position of the maximum can be roughly evaluated from (22) as follows:
$$t_{max}\frac{1}{BA}log\left[\frac{\beta ^2}{\alpha ^2}\frac{B^2}{A^2}\right],$$
(23)
with the supplementary condition
$$\frac{\beta }{\alpha }>\frac{A}{B}$$
(24)
where $`A=2B_1`$ and $`B=2B_2`$, $`A<B`$. If the condition (24) is not fulfilled the differential cross section $`d\sigma /dt`$ for production of $`V^{}(2S)`$ vector mesons has no maximum and has a standard monotonous $`t`$\- behaviour like for production of $`V(1S)`$ mesons.
The nonmonotonous $`t`$\- behaviour of the differential cross section for $`\rho ^{}(2S)`$ and $`\varphi ^{}(2S)`$ production in the photoproduction limit is strikingly different from the familiar decrease with $`t`$ of the differential cross section for the $`\rho ^0(1S)`$ and $`\varphi ^0(1S)`$ real photoproduction. Here we can not insist on the precise form of the $`t`$ dependence of the differential cross section, the main emphasis is on the likely pattern of the $`t`$ dependence coming from the node effect.
At larger energies, $`W\text{ }>150`$ GeV for the $`\rho ^{}(2S)`$ photoproduction and $`W\text{ }>30`$ GeV for $`\varphi ^{}(2S)`$ photoproduction, the node effect becomes weaker and we predict the positive valued diffraction slope at $`t=0`$ because of positive valued matrix elements $`<V^{}(2S)|\sigma (\xi ,r)|\gamma >`$ (3) and $`<V^{}(2S)|\sigma (\xi ,r)B(\xi ,r)|\gamma >`$ on the r.h.s. of Eq. (16). For this reason, the nonmonotonous $`t`$ dependence of the differential cross section is changed for the monotonous one, but still the effective diffraction slope decreases towards small $`t`$ in contrary to the familiar increase for the $`\rho ^0(1S)`$ and $`\rho ^0(1S)`$ photoproduction (see Fig. 1).
Because of a possible overcompensation scenario for the longitudinally polarized $`\rho ^{}(2S)`$ and $`\varphi ^{}(2S)`$ mesons in the forward direction and at small $`Q^2`$ (see Ref. ), we present in Fig. 2 the model predictions for the differential cross sections as a function of $`t`$ at different energies $`W`$ and at fixed $`Q^2=0.5`$ GeV<sup>2</sup> for the production of $`T`$, $`L`$ polarized and polarization unseparated $`\rho ^{}(2S)`$ and $`\varphi ^{}(2S)`$ mesons. As it was mentioned above, at $`Q^2=0.5`$ GeV<sup>2</sup>, the node effect becomes weaker, the amplitude for $`\rho _T^{}(2S)`$ and $`\varphi _T^{}(2S)`$ production at $`t=0`$ is in undercompensation regime and the corresponding slope parameter $`B(V_T^{}(2S))`$ is positive valued because of the positive valued matrix element $`<V^{}(2S)|\sigma (r,z)B(r,z)|\gamma ^{}>`$ on the r.s.h. of Eq. (16). For this reason, we predict the standard decrease of $`d\sigma \left(\gamma ^{}V_T^{}(2S)\right)/dt`$ with $`t`$ (see bottom boxes in Fig. 2). The above mentioned maximum of $`d\sigma /dt`$ for the undercompensation regime is absent due to a weaker node effect and because the condition (24) is not fulfilled.
However, for $`Q^2\text{ }<0.5`$ GeV<sup>2</sup>, the amplitude for $`\rho _L^{}(2S)`$ and $`\varphi _L^{}(2S)`$ production in forward direction, $`(t=0)`$ (and the matrix element $`<V_L^{}(2S)|\sigma (r,z)B(r,z)|\gamma ^{}>`$ as well), is still in overcompensation regime with the positive valued diffraction slope $`B(V_L^{}(2S))`$ at small energies $`W\text{ }<20`$ GeV. It follows in anomalous $`t`$ dependence of $`d\sigma \left(\gamma ^{}V_L^{}(2S)\right)/dt`$ shown if Fig. 2 (middle boxes). With the increase of $`t`$, because of the above mentioned interference of two different contributions to the production amplitude with different $`t`$ dependencies, one encounters the exact cancellation of the large and small distance contributions. This fact corresponds to the exact node effect at some $`tt_{min}`$. Thus, the differential cross section firstly falls down rapidly with $`t`$, have a minimum at $`tt_{min}`$, following by a rise when the overcompensation scenario of $`t`$\- dependent production amplitude is changed for the undercompensation one and the slope parameter becomes to be negative. At larger $`t`$, further pattern of $`t`$\- behaviour is practically the same as the nonmonotonous $`t`$ dependence of $`d\sigma \left(\gamma ^{}V_T^{}(2S)\right)/dt`$ at $`Q^2=0`$ (see Fig. 1).
The position of the minimum, $`t_{min}`$, in differential cross section is model dependent and can be roughly estimated from (22)
$$t_{min}\frac{1}{BA}log\left[\frac{\beta ^2}{\alpha ^2}\right].$$
(25)
With our wave functions we find $`t_{min}0.03`$ GeV<sup>2</sup> for $`\rho _L^{}(2S)`$ production and $`t_{min}0.05`$ GeV<sup>2</sup> for $`\varphi _L^{}(2S)`$ production at $`Q^2=0.5`$ GeV<sup>2</sup> and at $`W=5`$ GeV. However, we can not exclude a possibility that this minimum will take a place at larger $`t`$. At $`Q^2<0.5`$ GeV<sup>2</sup>, $`t_{min}`$ will be also located at larger values of $`t`$. At higher energy, the position of $`t_{min}`$ is shifted to a smaller value of $`t`$ unless the exact node effect is reached at $`t=0`$. At still larger energy, when longitudinally polarized $`2S`$ production amplitude is in undercompensation regime, this minimum disappears and we predict the pattern of $`t`$\- behaviour of $`L`$ differential cross section very similar to one like nonmonotonous $`t`$ dependence of $`d\sigma \left(\gamma V_T^{}(2S)\right)/dt`$ in the photoproduction limit described in Fig. 1. These predicted anomalies can be tested at HERA measuring the diffractive electroproduction of $`2S`$ radially excited light vector mesons in the separate polarizations, $`T`$ and $`L`$.
Conclusions
We study the diffractive photo- and electroproduction of ground state $`V(1S)`$ and radially excited $`V^{}(2S)`$ vector mesons within the color dipole gBFKL dynamics with the main emphasis related to the differential cross section $`d\sigma /dt`$, which is connected with the diffraction slope. There are two main consequences of vector meson production coming from the gBFKL dynamics. First, the energy dependence of the $`1S`$ vector meson production is controlled by the energy dependence of the dipole cross section which is steeper for smaller dipole sizes. The energy dependence of the diffraction slope for $`V(1S)`$ production is given by the effective Regge slope with a small variation with energy. Second the $`Q^2`$ dependence of the $`1S`$ vector meson production is controlled by the shrinkage of the transverse size of the virtual photon and the small dipole size dependence of the color dipole cross section. The $`Q^2`$ behaviour of the diffraction slope is given by the simple geometrical properties, $`r^2`$, coming from the color dipole gBFKL phenomenology of the slope parameter.
The diffraction slope for the production of $`2S`$ light vector mesons shows very interesting and anomalous behaviour as function of c.m.s. energy $`W`$ and $`Q^2`$ and will be detaily analysed elsewhere . As a consequence of the node in $`2S`$ radial wave function, we predict a strikingly different $`t`$ dependence of the differential cross section for production of $`V^{}(2S)`$ vs. $`V(1S)`$ mesons. The origin is in destructive interference of the large distance negative contribution to the production amplitude from the region above the node position with a steeper $`t`$\- dependence and small distance positive contribution to the production amplitude from the region below the node position with a weaker $`t`$\- dependence. As a result, at $`Q^2=0`$ (when the $`T`$ polarized $`V_T^{}(2S)`$ mesons are only produced) as a consequence of the undercompensation scenario for $`T`$ polarized forward production amplitude, we predict a nonmonotonous $`t`$\- dependence of $`d\sigma \left(\gamma V_T^{}(2S)\right)/dt`$ and a decreasing effective diffraction slope for $`V_T^{}(2S)`$ mesons towards to negative values at small $`t`$ in contrary with the familiar increase for the $`V(1S)`$ mesons. The differential cross section $`d\sigma \left(\gamma V_T^{}(2S)\right)/dt`$ firstly rises with $`t`$ having a maximum at $`tt_{max}`$ given by Eq.(23). At large $`t`$ when the node effect is weaker $`d\sigma (\gamma V_T^{}(2S))/dt`$ has the standard monotonous $`t`$\- behaviour like for production of $`V(1S)`$ vector mesons. The position of the maximum is model dependent and is shifted to smaller values of $`t`$ with rising energy and $`Q^2`$ due to a weaker node effect.
For production of $`L`$ polarized $`V_L^{}(2S)`$ mesons, there is overcompensation at $`t=0`$ leading to an exact cancellation of the positive contribution from large size dipoles and the negative contribution from small size dipoles to the production amplitude and to a minimum of the differential cross section at some value of $`tt_{min}`$. The position of $`t_{min}`$ is given by Eq. (25), is energy dependent and leads to a complicated anomalous $`t`$ dependence of $`d\sigma \left(\gamma ^{}V_L^{}(2S)\right)/dt`$ at fixed $`Q^2`$. Thus, $`d\sigma \left(\gamma ^{}V_L^{}(2S)\right)/dt`$ firstly falls down with $`t`$ having a minimum at $`tt_{min}`$ when the overcompensation scenario is changed for the undercompensation one. The following pattern of $`t`$\- behaviour is then the same like for $`d\sigma (\gamma V_T^{}(2S))/dt`$ at $`Q^2=0`$. These anomalies are also energy and $`Q^2`$\- dependent and can be studied at HERA.
The experimental investigation of $`t`$\- dependent differential cross section for real photoproduction ($`Q^2=0`$) of $`V^{}(2S)`$ mesons at fixed target and HERA experiments, offers an unique possibility to make a choice between the undercompensation and overcompensation scenarios. The presence of the minimum in $`t`$\- dependent $`d\sigma (\gamma V^{}(2S))/dt`$ in a broad energy region from small to large energies, corresponds to the overcompensation scenario, whereas its absence corresponds to the undercompensation scenario.
The position of the node in the radial $`(2S)`$ wave function can be tested also by the vector meson data with separate polarizations $`(L)`$ and $`(T)`$ at $`Q^2>0`$. The existence of the minimum in $`t`$\- dependent differential cross section is connected again with the overcompensation scenario in $`(2S)`$ production amplitude whereas the undercompensation scenario reflects the maximum of $`d\sigma /dt`$ and/or the standard monotonous $`t`$\- behaviour.
Figure captions:
* \- The color dipole model predictions for the differential cross sections $`d\sigma (\gamma ^{}V(V^{}))/dt`$ for the real photoproduction ($`Q^2=0`$) of the $`\rho ^0,\rho ^{}(2S),\varphi ^0`$ and $`\varphi ^{}(2S)`$ at different values of the c.m.s. energy $`W`$.
* \- The color dipole model predictions for the differential cross sections $`d\sigma _{L,T}(\gamma ^{}V^{})/dt`$ for transversely (T) (top boxes) and longitudinally (L) (middle boxes) polarized radially excited $`\rho ^{}(2S)`$, $`\varphi ^{}(2S)`$ and for the polarization-unseparated $`d\sigma (\gamma ^{}V^{})/dt=d\sigma _T(\gamma ^{}V^{})/dt+ϵd\sigma _L(\gamma ^{}V^{})/dt`$ for $`ϵ=1`$ (bottom boxes) at $`Q^2=0.5`$ GeV<sup>2</sup> and different values of the c.m.s. energy $`W`$. |
warning/0003/astro-ph0003152.html | ar5iv | text | # Bound systems in an expanding universe
## I INTRODUCTION AND SUMMARY
Currently there are two general relativistic descriptions of spacetime in popular use. For planetary systems and other gravitationally bound structures which are small on the scale of the universe, there is a static description of the behavior of spacetime. On the other hand, for large-scale behavior, there is a time dependent description which is appropriate as a description of phenomena such as the observed red-shift of distant galaxies.
The classic question is, “How can these two disparate descriptions of spacetime possibly be reconciled with each other?” The current standard answer is that these two are meshed together on a spherical surface surrounding a mass concentration which grows with time in the static metric, but stays at a fixed coordinate radius in the non-static metric. This is the “Schwarzschild solution in a cosmological model” picture, or it is also called the “Swiss cheese model.” This later name refers to the fact the in this picture, the background material is removed inside the spherical boundary. The mass removed depends on the mass of the central concentration and the curvature of space. The replacement of the mass interior to the sphere by a concentration of mass at the center is based on Birkhoff’s theorem which says that in a homogeneous, zero pressure cosmological model, as long as the material inside a sphere is spherically symmetric, we can replace it with a compact mass at the center with no change on its exterior effects. To make this matching work, the pressure of the exterior solution (Friedmann-Lemaître metric class of the general Robertson-Walker line elements) must vanish, as does the pressure of the interior (the exterior Schwarzschild metric) solution. This condition places a restriction on the form of the universal expansion factor of the overall universe. Although the matching conditions can be met, as we shall see in the third section, an additional problem arises when the dynamics are considered.
In the second section, for the convenience of the reader, I gather together the necessary classical equations from general relativity for the study at hand. One non-classical item in this section is, instead of the usual $`3+1`$ spacetime split into 3 space and 1 time dimensions, I split spacetime into one radial coordinate, and time plus the two angles of spherical coordinates. The change allows the direct computation of the relevant extrinsic curvatures.
In the third section, I compute the stress-energy tensors and the extrinsic curvatures for both the Schwarzschild and the Friedman-Lemaître metrics used in the “Swiss cheese model.” By a reparameterization of the Schwarzschild metric, both the intrinsic and the extrinsic curvature can be made to be continuous. I remark that the stress-energy tensor, for certain parameter choices, displays no pressure discontinuity but only a cosmic fluid density discontinuity. That discontinuity is in line with the “Swiss cheese model” idea that there are holes in the cosmic fluid.
In the fourth section, I show that it can happen, for trajectories which approach the metric interface an near grazing angles, that the subsequent trajectories are discontinuous functions of their initial conditions. Those which enter the inner or Schwarzschild metric region can be bound in a finite sized, closed orbit, while those which do not, travel on a parabolic trajectory to infinity. To emphasis, this case is not the same as in Newtonian orbit theory where ellipses of progressively larger size blend into parabolas, but here the ellipse is just finite in size!
In the fifth section, I introduce an alternative metric. This metric is basically an adaptation of the Schwarzschild metric in curved space. I compute the stress-energy tensor. It shows an isotropic pressure, and no mass-density flux. It is only second order in magnitude in both the Hubble constant and the inverse radius of curvature of the universe. The same statement is true of the Friedmann-Lemaître stress energy tensor. In addition I have computed the extrinsic curvature. Both the stress-energy tensor and the extrinsic curvature are continuous, outside the Schwarzschild radius of course, as they come from infinitely differentiable expressions.
In the sixth section, I compute the equations of motion for a freely moving test particle in the alternative metric. I then transform them to a coordinate system at rest at the center of the mass concentration. The flat space, slowly moving particle, weak gravitational field limit of these equations of motion are also given. The only correction to Newton’s equations of motion in this limit is a term proportional to the square of Hubble’s constant, $`H_0`$.
In the final section, I gives some examples of the dynamics found using my alternative metric.
## II METRICS
It is useful to review the properties of a few metrics. The line elements will all be of the general form,
$`ds^2`$ $`=`$ $`e^\mu [(dx^1)^2+(dx^2)^2+(dx^3)^2]+e^\nu (dx^4)^2`$ (1)
$`=`$ $`g_{ij}dx^idx^j,`$ (2)
where the Einstein summation convention is used, and where
$$r^2=(x^1)^2+(x^2)^2+(x^3)^2,\mu =\mu (r,t),\nu =\nu (r,t),$$
(3)
appropriate to the non-static, spherically symmetric case. Eq. 2 defines the metric tensor $`g_{ij}`$. It will be of interest to know what Einstein field equation is satisfied for each of these metrics. The equation is,
$$R_i^j\frac{1}{2}Rg_i^j+\mathrm{\Lambda }g_i^j=8\pi T_i^j,$$
(4)
where $`T_{ij}`$ is the stress-energy tensor, and $`R_{ij}`$ is the Ricci tensor, which is a contraction of Riemann’s four index curvature tensor. The Ricci tensor can be expressed in terms of the three index Christoffel symbols $`\mathrm{\Gamma }`$ as,
$$R_{km}=(\mathrm{\Gamma }_{km}^i)_{,i}(\mathrm{\Gamma }_{ki}^i)_{,m}+\mathrm{\Gamma }_{ni}^i\mathrm{\Gamma }_{km}^n\mathrm{\Gamma }_{mn}^i\mathrm{\Gamma }_{ki}^n,$$
(5)
where the notation $`)_{,i}`$ means take the partial derivative with respect to $`x^i`$. In turn, the Christoffel symbols are defined in terms of the metric tensor as,
$$\mathrm{\Gamma }_{ij}^m=\frac{1}{2}g^{mk}(g_{ki,j}+g_{kj.i}g_{ij,k}),\mathrm{where}g^{mk}g_{kj}=\delta _j^m,$$
(6)
defines $`g^{km}`$, and $`\delta _j^m`$ is the Kronecker delta function. Finally,
$$R=g^{ij}R_{ij},$$
(7)
is the contraction of the Ricci tensor.
The $`T^{44}`$ element is the mass-energy density. The $`T^{4\beta }`$ is the mass-flux density through an area perpendicular to the direction $`\beta `$ per unit time. (Greek indices run from 1 to 3 while Roman indices run from 1 to 4.) The $`T^{\alpha \alpha }`$ element is the pressure in the $`\alpha `$ direction, and the $`T^{\alpha \beta }`$ element is the flux density of the $`\alpha `$ component of momentum in the $`\beta `$ direction. It is manifest that eq. 4 allows the direct computation of the stress energy tensor from the metric tensor. Tolman gives the results for the form of the line element
$$ds^2=e^\mu (dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\varphi ^2)+e^\nu dt^2,$$
(8)
which is the spherical coordinate version or eq. 2. The non-vanishing elements of the stress-energy tensor are
$`8\pi T_1^1`$ $`=`$ $`e^\mu \left({\displaystyle \frac{\mu ^2}{4}}+{\displaystyle \frac{\mu ^{}\nu ^{}}{2}}+{\displaystyle \frac{\mu ^{}+\nu ^{}}{r}}\right)`$ (10)
$`+e^\nu \left(\ddot{\mu }+\frac{3}{4}\dot{\mu }^2{\displaystyle \frac{\dot{\mu }\dot{\nu }}{2}}\right)\mathrm{\Lambda }`$
$`8\pi T_2^2`$ $`=`$ $`8\pi T_3^3=e^\mu \left({\displaystyle \frac{\mu ^{\prime \prime }}{2}}+{\displaystyle \frac{\nu ^{\prime \prime }}{2}}+{\displaystyle \frac{\nu ^2}{4}}+{\displaystyle \frac{\mu ^{}+\nu ^{}}{2r}}\right)`$ (12)
$`+e^\nu \left(\ddot{\mu }+\frac{3}{4}\dot{\mu }^2{\displaystyle \frac{\dot{\mu }\dot{\nu }}{2}}\right)\mathrm{\Lambda }`$
$`8\pi T_4^4`$ $`=`$ $`e^\mu \left(\mu ^{\prime \prime }+{\displaystyle \frac{\mu ^2}{4}}+{\displaystyle \frac{2\mu ^{}}{r}}\right)+\frac{3}{4}e^\nu \dot{\mu }^2\mathrm{\Lambda }`$ (13)
$`8\pi T_4^1`$ $`=`$ $`e^\mu \left(\dot{\mu }^{}{\displaystyle \frac{\dot{\mu }\nu ^{}}{2}}\right)`$ (14)
$`8\pi T_1^4`$ $`=`$ $`e^\nu \left(\dot{\mu }^{}{\displaystyle \frac{\dot{\mu }\nu ^{}}{2}}\right)`$ (15)
where $`\mathrm{\Lambda }`$ is the cosmological constant, an over dot denotes the time derivative, and a prime denotes the derivative with respect to $`r`$.
Another quantity which is useful to consider is the extrinsic curvature. This concept arises in the Arnowitt et al. splitting of four dimensional spacetime into 3 dimensional space plus one dimensional time. The relationship between the metrics is given by
$`ds^2`$ $`=`$ $`{}_{}{}^{(4)}g_{ij}^{}dx^idx^j`$ (16)
$`=`$ $`{}_{}{}^{(3)}g_{\alpha \beta }^{}(dx^\alpha +N^\alpha dt)(dx^\beta +N^\beta dt)+N^2dt^2,`$ (17)
where the $`N^\alpha `$ are the three shift functions and $`N`$ is the lapse (of proper time) function. The intrinsic curvature is the analogue of $`R`$ \[eq. 7\] in three dimensions. The extrinsic curvature measures the fractional shrinkage and deformation as one advances in time from one space-like hyperplane to the next. The extrinsic curvature tensor is given as
$$K_{\alpha \beta }=\frac{1}{2N}\left[N_{\alpha |\beta }+N_{\beta |\alpha }\frac{g_{\alpha \beta }}{t}\right],$$
(18)
where the notation $`)_{|\alpha }`$ means the covariant derivative with respect to $`x^\alpha `$.
The reason for the interest in the extrinsic curvature in our case is that the necessary and sufficient junction conditions to join two metrics in 4 dimensional spacetime across a 3 dimensional hypersurface is that both $`g_{\alpha \beta }`$ and $`K_{\alpha \beta }`$ be continuous across the surface. In our case we are concerned with line elements of the class of eq. 8. The split is between $`r`$ instead of $`t`$ and the other three variables. In this case the shift functions are all zero, and the lag function is $`N=\mathrm{exp}(0.5\mu )`$. The non-zero elements of $`K`$ are
$`K_{22}`$ $`=`$ $`\frac{1}{2}(\mu ^{}r^2+2r)\mathrm{exp}(0.5\mu ),`$ (19)
$`K_{33}`$ $`=`$ $`\frac{1}{2}(\mu ^{}r^2+2r)\mathrm{sin}^2\theta \mathrm{exp}(0.5\mu ),`$ (20)
$`K_{44}`$ $`=`$ $`\frac{1}{2}\nu ^{}\mathrm{exp}(\nu 0.5\mu ).`$ (21)
Returning to the three-space, constant-time formalism, we give the equations of motion of a free test particle as seen by local co-moving observers along the test particle’s path. That is to say, a set of observers whose coordinates do not change with time. In other words, we need the equations for the geodesic curves in spacetime. For the class of line elements we are considering, it is simplest to start with a Lagrangian formulation. By eq. 2 we may write this formulation as,
$$s=Ldt,L=\frac{ds}{dt}=\left(g_{ij}\dot{x}^i\dot{x}^j\right)^{1/2}$$
(22)
The standard Euler-Lagrange equations for an extreme in path length (Here we seek a minimum distance between two fixed endpoints.) are
$$\frac{d}{dt}\left(\frac{L}{\dot{x}^\alpha }\right)=\frac{L}{x^\alpha }.$$
(23)
Thus, using the diagonal nature of the metric tensor, we obtain for the standard geodesic equations,
$$\left(\frac{ds}{dt}\right)\frac{d}{dt}\left[g_{\alpha i}\left(\frac{ds}{dt}\right)^1\dot{x}^i\right]=\frac{1}{2}g_{ij,\alpha }\frac{dx^i}{dt}\frac{dx^j}{dt},$$
(24)
where, of course $`dx^4/dt=1`$. From the line element we get
$$\left(\frac{ds}{dt}\right)=\left[e^\nu +g_{\alpha \beta }\dot{x}^\alpha \dot{x}^\beta \right]^{1/2}$$
(25)
These results display, for the class of line elements we are considering, the three, second-order, non-linear, coupled equations for the three coordinates $`x^\alpha `$ of a test particle as a function of time. We have assumed isosynchronous coordinates everywhere in the three-dimensional, spacelike hyper-surface. The square roots can removed from eq. 24 by rewriting it as
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{ds}{dt}}\right)^2{\displaystyle \frac{d}{dt}}\left[\left({\displaystyle \frac{ds}{dt}}\right)^2\right]\left[g_{\alpha i}\dot{x}^i\right]`$ (26)
$`+{\displaystyle \frac{d}{dt}}\left[g_{\alpha i}\dot{x}^i\right]={\displaystyle \frac{1}{2}}g_{ij,\alpha }{\displaystyle \frac{dx^i}{dt}}{\displaystyle \frac{dx^j}{dt}},`$ (27)
## III “Swiss Cheese Model”
As mentioned in the first section, there is a popular cosmological model which in the large has the Friedmann-Lemaître line element,
$`ds^2`$ $`=`$ $`{\displaystyle \frac{a(t)^2}{[1+(r/2R)^2]^2}}[(dx^1)^2+(dx^2)^2+(dx^3)^2]`$ (29)
$`+c^2(dx^4)^2`$
where in the notation of Sec. II,
$`e^\mu `$ $`=`$ $`{\displaystyle \frac{a(t)^2}{[1+(a(t)r/2a(t)R)^2]^2}}`$ (30)
$`e^\nu `$ $`=`$ $`c^2.`$ (31)
where $`c`$ is the velocity of light and $`a(t)R`$ is the radius of curvature of the model universe.
The non-zero elements of the stress-energy tensor associated with this line element are, by eq. 13,
$`8\pi T_1^1`$ $`=`$ $`8\pi T_2^2=8\pi T_3^3={\displaystyle \frac{1}{[a(t)R]^2}}+2{\displaystyle \frac{\ddot{a}}{ac^2}}+\left({\displaystyle \frac{\dot{a}}{ac}}\right)^2\mathrm{\Lambda }`$ (33)
$`=8\pi p_0,`$
$`8\pi T_4^4`$ $`=`$ $`{\displaystyle \frac{3}{[a(t)R]^2}}+3\left({\displaystyle \frac{\dot{a}}{ac}}\right)^2\mathrm{\Lambda }=8\pi \rho _{00}`$ (34)
The extrinsic curvature tensor, as given by eq. 20 for this metric has the non-zero elements,
$$K_{22}=\frac{a(t)r\left[1(r/2R)^2\right]}{[1+(r/2R)^2]^2},K_{33}=\mathrm{sin}^2\theta K_{22}.$$
(35)
In this cosmological model on scales small compared to that of the universe, as in, for example, the solar system, a quite different metric is used. It is Schwarzschild’s exterior solution. The line element for it is
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{m}{2r}}\right)^4[(dx^1)^2+(dx^2)^2+(dx^3)^2]`$ (37)
$`+c^2\left({\displaystyle \frac{1m/2r}{1+m/2r}}\right)^2(dx^4)^2,`$
where in the notation of Sec. II,
$$e^\mu =\left(1+\frac{m}{2r}\right)^4,e^\nu =c^2\left(\frac{1m/2r}{1+m/2r}\right)^2.$$
(38)
Here $`m`$ is an abbreviation for $`GM/c^2`$ where $`M`$ is the mass concentration. This metric is only valid outside the Schwarzschild radius $`r>>r_0=2GM/c^2`$. Inside this radius, Schwarzschild’s interior solution is required but we shall not be concerned with this aspect here.
There are no non-zero elements of the stress-energy tensor associated with this line element, unless the cosmological constant is different from zero. Thus it corresponds to zero pressure and zero density, except for a mass concentration in the center. If $`\mathrm{\Lambda }0`$ then,
$$T_1^1=T_2^2=T_3^3=T_4^4=\mathrm{\Lambda }/8\pi $$
(39)
This result would correspond to a uniform density and pressure through out space, rather than the empty space with $`\mathrm{\Lambda }=0`$.
A cross comparison of eq. 39 with eq. 33 shows that unless,
$$\frac{1}{[a(t)R]^2}+\left(\frac{\dot{a}}{ac}\right)^2=\frac{\ddot{a}}{ac^2}=0,$$
(40)
the stress-energy tensor is discontinuous at the boundary. The only solution of these equations, as $`R`$ is a constant, is $`\dot{a}`$ is a constant, i.e., a linear expansion factor, or a static flat universe.
Before we can us eq. 20 to compute extrinsic curvature for the Schwarzschild case, we must first reparameterize the metric so that the “Swiss cheese model” metric boundary is given by one of the coordinates equals a constant. It will be convenient to use spherical coordinates. If we follow the textbook approach, then we should start with a zero pressure cosmological model. By Birkoff’s theorem in this case we may hollow out a sphere and replace it by a mass concentration at the center. Since the density of the “cosmic fluid” is inversely proportional to $`a^3`$, the radius of the sphere is fixed in the Friedmann-Lemaître coordinates. The surface of the sphere is a three dimensional hypersurface parameterized by the time and the two angle variables of spherical coordinates. To proceed, we note that an observer sitting on the boundary is on a geodesic for $`r_{FL}=`$ a constant for all time. The relation between the Friedmann Lemaître time $`\overline{t}`$ and the Schwarzschild variables as seen by this observer is
$`(ds)^2`$ $`=`$ $`c^2(d\overline{t}_e)^2`$ (41)
$`=`$ $`\left[c^2\left({\displaystyle \frac{1{\displaystyle \frac{m}{2r_e}}}{1+{\displaystyle \frac{m}{2r_e}}}}\right)^2\left(1+{\displaystyle \frac{m}{2r_e}}\right)^4\left({\displaystyle \frac{dr_e}{dt_e}}\right)^2\right](dt_e)^2,`$ (42)
Since the observer is on a geodesic, we may deduce, by means of the time component of eq. 24, that he sees
$$c\left(\frac{1{\displaystyle \frac{m}{2r}}}{1+{\displaystyle \frac{m}{2r}}}\right)^2\frac{dt}{ds}=K^{1/2},$$
(43)
where $`K`$ is a constant of integration. Note is taken that, as the integration is over $`t`$, $`K`$ may depend on $`r`$, of course. By combining eq. 41 and eq. 43, we obtain,
$`{\displaystyle \frac{dr_e}{dt_e}}`$ $`=`$ $`{\displaystyle \frac{c\left(1{\displaystyle \frac{m}{2r_e}}\right)}{\left(1+{\displaystyle \frac{m}{2r_e}}\right)^3}}\left[1{\displaystyle \frac{1}{K}}\left({\displaystyle \frac{1{\displaystyle \frac{m}{2r_e}}}{1+{\displaystyle \frac{m}{2r_e}}}}\right)^2\right]^{1/2},`$ (44)
$`{\displaystyle \frac{dt_e}{d\overline{t}_e}}`$ $`=`$ $`K^{1/2}\left({\displaystyle \frac{1+{\displaystyle \frac{m}{2r_e}}}{1{\displaystyle \frac{m}{2r_e}}}}\right)^2,`$ (45)
$`{\displaystyle \frac{dr_e}{d\overline{t}_e}}`$ $`=`$ $`{\displaystyle \frac{cK^{1/2}}{\left(1{\displaystyle \frac{m^2}{4r_e^2}}\right)}}\left[1{\displaystyle \frac{1}{K}}\left({\displaystyle \frac{1{\displaystyle \frac{m}{2r_e}}}{1+{\displaystyle \frac{m}{2r_e}}}}\right)^2\right]^{1/2},`$ (46)
which gives the behavior of $`r_e,t_e`$ as a function of the Friedmann-Lemître time variable, as seen by the comoving observer sitting on the metric boundary in the Swiss-cheese model.
By equating the coefficients of the angular variables, we obtain,
$`\left(1+{\displaystyle \frac{m}{2r_e}}\right)^2r_e`$ $`=`$ $`{\displaystyle \frac{a(\overline{t}_e)\overline{r}_e}{1+\left({\displaystyle \frac{\overline{r}_e}{2R}}\right)^2}},`$ (47)
$`{\displaystyle \frac{\dot{a}(\overline{t}_e)}{a(\overline{t}_e)}}`$ $`=`$ $`{\displaystyle \frac{1}{r_e}}{\displaystyle \frac{dr_e}{dt_e}}\left({\displaystyle \frac{1{\displaystyle \frac{m}{2r_e}}}{1+{\displaystyle \frac{m}{2r_e}}}}\right){\displaystyle \frac{dt_e}{d\overline{t}_e}}`$ (48)
Note is taken that the boundary coordinate $`\overline{r}_e`$ in the Friedmann-Lemaître metric depends on the central mass concentration, and on the radius of curvature of the universe, i.e. on $`m`$ and $`R`$.
Next we combine eqs. 45 and 47 and comparing the result with the “cosmological equation.” This equation is the equation for the $`T_4^4`$ component of the stress energy tensor derived from the Friedmann-Lemaître line element. The result of the application of this textbook method is that the constant of integration is
$$K=\left[\frac{1\left({\displaystyle \frac{\overline{r}_e}{2R}}\right)^2}{1+\left({\displaystyle \frac{\overline{r}_e}{2R}}\right)^2}\right]^2,$$
(49)
which is a function of $`m`$ and $`R`$.
We are now in a position to introduce a reparameterization of the Schwarzschild metric. We choose the new parameters $`\widehat{t}=\overline{t}`$ and $`\widehat{r}`$. The second variable is chosen so as to keep the metric continuous at the interface, and to keep the metric diagonal, if possible. The reparameterized Schwarzschild metric is given by eq. 81. Eq. 41 insures that $`\widehat{g}_{44}=c^2`$. The equation of continuity, and the vanishing of the elements $`\widehat{g}_{14}=\widehat{g}_{41}`$ yield the conditions,
$`\left({\displaystyle \frac{dr_e}{d\widehat{r}_e}}\right)^2g_{11}+\left({\displaystyle \frac{dt_e}{d\widehat{r}_e}}\right)^2g_{44}`$ $`=`$ $`\overline{g}_{11},`$ (50)
$`\left(g_{11}{\displaystyle \frac{dr}{d\widehat{t}}}\right){\displaystyle \frac{dr}{d\widehat{r}}}+\left({\displaystyle \frac{dt}{d\widehat{t}}}g_{44}\right){\displaystyle \frac{dt}{d\widehat{r}}}`$ $`=`$ $`0`$ (51)
The solution of these equations is
$`{\displaystyle \frac{dr}{d\widehat{r}}}`$ $`=`$ $`{\displaystyle \frac{K^{1/2}a(\widehat{t})}{\left(1{\displaystyle \frac{m^2}{4r_e^2}}\right)\left[1+\left({\displaystyle \frac{\overline{r}_e}{2R}}\right)^2\right]}},`$ (52)
$`{\displaystyle \frac{dt}{d\widehat{r}}}`$ $`=`$ $`{\displaystyle \frac{K^{1/2}a(\widehat{t})\left(1+{\displaystyle \frac{m}{2r}}\right)^3}{c\left[1+\left({\displaystyle \frac{\overline{r}_e}{2R}}\right)^2\right]\left(1{\displaystyle \frac{m^2}{4r_e^2}}\right)\left(1{\displaystyle \frac{m}{2r}}\right)}}`$ (54)
$`\times \left[1{\displaystyle \frac{1}{K}}\left({\displaystyle \frac{1{\displaystyle \frac{m}{2r}}}{1+{\displaystyle \frac{m}{2r}}}}\right)^2\right]^{1/2},`$
$`\widehat{g}_{11}`$ $`=`$ $`{\displaystyle \frac{\left(1{\displaystyle \frac{m^2}{4r^2}}\right)^2a(\widehat{t})^2}{\left(1{\displaystyle \frac{m^2}{4r_e^2}}\right)^2\left[1+\left({\displaystyle \frac{r_e}{2R}}\right)^2\right]^2}}`$ (55)
Our construction insures the continuity of $`\widehat{g}_{11}=\overline{g}_{11}`$ at the metric interface. The continuity of the (22) and the (33) elements are insured by eq. 47. All the off-diagonal elements vanish.
We may now apply eq. 20 to obtain the extrinsic curvature. The result is just eq. 35 with $`(r,t,\theta )`$ replaced by $`(\widehat{r},\widehat{t},\widehat{\theta })`$. This result shows the continuity of the extrinsic curvature.
There is however, one important item to note. Consider the derivatives perpendicular to the boundary hypersurface, evaluated at the interface.
$`{\displaystyle \frac{d\mathrm{log}(\overline{g}_{11})}{d\overline{r}}}`$ $`=`$ $`{\displaystyle \frac{\overline{r}}{R^2+0.25\overline{r}^2}},`$ (56)
$`{\displaystyle \frac{d\mathrm{log}(\widehat{g}_{11})}{d\widehat{r}}}`$ $`=`$ $`{\displaystyle \frac{a(\widehat{t}){\displaystyle \frac{m}{r_e^3}}\left[1\left({\displaystyle \frac{r_e}{2R}}\right)^2\right]}{\left(1{\displaystyle \frac{m^2}{4r_e^2}}\right)^3\left[1+\left({\displaystyle \frac{r_e}{2R}}\right)^2\right]^2}}.`$ (57)
To leading order in a flat universe, eq. 57 becomes,
$$\frac{d\mathrm{log}(\overline{g}_{11})}{d\overline{r}}=0,\frac{d\mathrm{log}(\widehat{g}_{11})}{d\widehat{r}}=a(\widehat{t})\frac{m}{r_e^3}.$$
(58)
That such a discontinuity may occur is well known. I will explore some of the consequences of this discontinuity in the next section.
As was remarked above, it is well known that it is a necessary condition that the pressure equal zero in order for this matching to occur. As the pressure $`p=T_\alpha ^\alpha `$ for each $`\alpha `$ (not summed here), the zero-pressure condition means, by eq. 39 which gives the pressure for the Schwarzschild case in terms of the $`T_\alpha ^\alpha `$ elements that $`\mathrm{\Lambda }=0`$. Turning to the the Friedmann-Lemaître case, we see from eq. 33 that it must be that
$$\frac{1}{[a(t)R]^2}+2\frac{\ddot{a}}{ac^2}+\left(\frac{\dot{a}}{ac}\right)^2=0.$$
(59)
If use substitute the standard form $`a(t)=(t/t_0)^\psi `$ in eq. 59, then the only solutions are $`R=\mathrm{}`$ together with $`\psi =0`$ or $`2/3`$. The first solution is the trivial case of a static universe and is of no interest in the present discussion. In the second case, the stress-energy tensor element $`T_4^4>0`$, which means there is a mass discontinuity. This is expected by the structure of the approximation of scooping out a hollow sphere, and then having a uniform density outside.
## IV Dynamics in the “Swiss Cheese Model”
Following up on the discontinuity in the derivate of the metric, perpendicular to the metric boundary in the “Swiss Cheese Model,” we investigate the some of the dynamic properties of this model. It suffices for my purposes to consider the simpler case of a flat ($`R=\mathrm{}`$) universe. The equations of motion are as follows. For the Schwarzschild metric, we get from eq. 24, to leading order, the well known Newtonian equations
$$\frac{d^2\stackrel{}{\rho }}{dt^2}=\frac{GM}{\rho ^3}\stackrel{}{\rho }$$
(60)
The solutions to this equation are the familiar Newtonian conic sections.
For the Friedmann-Lemaître metric from eq. 24 we get a first intergal as
$$\dot{\stackrel{}{x}}=\frac{\stackrel{}{A}}{a^2(t)}$$
(61)
For my purposes, it is more convenient to use a variable more closely equal to the proper distances. Thus for $`\stackrel{}{\rho }=a(t)\stackrel{}{x}`$, eq. 61 becomes,
$$\frac{d^2\stackrel{}{\rho }}{dt^2}=\frac{\ddot{a}(t)}{a(t)}\stackrel{}{\rho }$$
(62)
In both metrics $`\rho `$ is within plotting accuracy for the proper distance in the examples I will consider.
The general solution of eq. 62 is
$$\stackrel{}{\rho }=\stackrel{}{A}t^\psi +\stackrel{}{B}t^{1\psi }=\stackrel{}{A}t^{2/3}+\stackrel{}{B}t^{1/3},$$
(63)
when the currently, theoretically favored value of $`\psi =\frac{2}{3}`$ is chosen. Let us take the example in rectangular coordinates where initially $`x=\lambda ,\dot{x}=0,y=0,\dot{y}=\lambda `$ at time $`t=t_0`$. Then eq. 63 becomes,
$`x`$ $`=`$ $`\lambda \left[\left({\displaystyle \frac{t}{t_0}}\right)^{2/3}+2\left({\displaystyle \frac{t}{t_0}}\right)^{1/3}\right],`$ (64)
$`y`$ $`=`$ $`3t_0\lambda \left[\left({\displaystyle \frac{t}{t_0}}\right)^{2/3}\left({\displaystyle \frac{t}{t_0}}\right)^{1/3}\right].`$ (65)
It is worth noting that for $`tt_0t_0`$ that the motion is almost exactly that of a straight line, as is to be expected in flat empty space. Specifically, direct calculation yields
$$\ddot{x}(t_0)=\frac{2\lambda }{9t_0^2},\ddot{y}(t_0)=0.$$
(66)
This apparent acceleration is of order $`H_0^2`$.
The solutions in eq. 63 are quadratic in the parameter $`t^{1/3}`$. Thus we can obtain the equation for the trajectory as a quadratic plus linear expression in $`\rho _x`$ and $`\rho _y`$. The form would be
$$\rho _y=c(a\rho _x+b\rho _y)^2+d(a\rho _x+b\rho _y)$$
(67)
This form is readily recognized as a parabola. Thus it is the case that a freely moving particle in an Friedmann-Lemaître expanding space always appears to be moving in a parabola. This effect is caused by the small $`\ddot{a}`$ term which appears in the equation 62 as a forcing term.
As an illustration of the behavior of the “Swiss cheese model” I have computed the following trajectories. I use as a unit of time the Hubble time, that is $`1/H_0`$, which is of the order of $`10^{18}`$ seconds. As a unit of distance I use $`\sqrt[3]{M_{}G/H_0^2}`$ which is about 20 million Astronomical units. $`M_{}`$ is the mass of the sun. I set $`t_0=1`$ to switch to our current units. The metric interface is at $`r=t^{2/3}`$ in these units for flat spacetime, as mentioned above. In order to follow a Friedmann-Lemaître trajectory the test particle’s distance from the origin must be larger at every time than that for the interface. Thus,
$$\lambda ^2\left[10t^{4/3}22t+13t^{2/3}\right]>t^{4/3},$$
(68)
The trajectory will intersect the interface if eq. 68 is an equality. By means of the quadratic formula, an intersection will occur if
$$t^{1/3}=\frac{11\pm \sqrt{9+13/\lambda ^2}}{101/\lambda ^2}.$$
(69)
It will be observed that for $`\lambda <\sqrt{13}/3`$ there are two real roots. If $`\lambda =1`$, then $`t^{1/3}=1,\mathrm{\hspace{0.33em}13}/9`$. On the other hand, if $`\lambda >\sqrt{13}/3`$, the roots are imaginary, so there are no intersections. If $`\lambda =\sqrt{13}/3`$ there is a double root at $`t^{1/3}=13/11`$. In this case the parabolic trajectory just grazes the metric interface.
In Fig. 1, I illustrate the two different trajectories when $`r_0=\sqrt{13}/3`$ and $`\dot{\varphi }_0=1.0`$. In the case where $`r_0`$ is just any arbitrary amount smaller, the expanding interface overtakes the test particle following its Friedmann-Lemaître parabolic trajectory and it must then follow the static Schwarzschild equations of motion. The Schwarzschild metric takes over at $`t=(13/11)^3`$ as explained above, and after that the trajectory is an ellipse with semimajor axis 7.6287 and the semiminor axis 3.986. These imply an eccentricity of 0.69068 and a semilatus rectum parameter of $`p=2.086`$. On the other hand, if $`r_0`$ is any arbitrary amount larger, it escapes the moving interface and continues to follow the parabolic trajectory. It is evident from Fig. 1 that future trajectories are, in some cases, discontinuous functions of the initial conditions for the “Swiss cheese model.”
Put another way, the Schwarzschild metric permits closed orbits and the Friedmann-Lemaître metric does not. In terms of the latter coordinates, one can choose initial conditions so that the parabolic trajectory just grazes the metric interface (fixed radial coordinate in this metric) and the speed is low enough that for infinitesimally different initial conditions the trajectory crosses the interface and is caught in a bound state, or alternatively misses the interface and proceeds on its parabolic trajectory. All these effects take place in supposedly empty space of the order of 20 million AU from a mass concentration of size $`M_{}`$, and are quite counter to one’s physical intuition that such discontinuities should not occur there.
## V An Acceptable Metric
I propose the following line element to represent a mass condensation in an expanding and curved universe. It is of the form 8 where
$`e^\mu `$ $`=`$ $`{\displaystyle \frac{a(t)^2}{[1+(a(t)r/2a(t)R)^2]^2}}\left(1+{\displaystyle \frac{m}{2a(t)r}}\right)^4,`$ (70)
$`e^\nu `$ $`=`$ $`c^2\left({\displaystyle \frac{1m/2a(t)r}{1+m/2a(t)r}}\right)^2,m{\displaystyle \frac{GM}{c^2}}\left[1+\left({\displaystyle \frac{r}{2R}}\right)^2\right]^{1/2}`$ (71)
where $`G`$ is Newton’s constant of gravitation. It is to be noted that this metric is an adaptation of the Schwarzschild metric in curved space. It is not claimed that this metric is unique. Certainly any coordinate transformation of this metric is equivalent. It does however have the property that in the limit where $`a(t)`$ is a constant, it reduces to the Schwarzschild metric in curved space. Also, when the central mass vanishes it reduces to the Friedmann-Lemaître metric. It is an infinitely differentiable solutions to the Einstein field equations which corresponded to values of the stress-energy tensor which are only of second order in $`R^1`$ and $`H_0`$, except at the central mass. This latter property is also true of the Friedmann-Lemaître metric in curved space, as may be seen in eq. 33. Thus we have, by this example whose properties are globally similar to the Friedmann-Lemaître metric, demonstrated that the Swiss cheese model is not require to match the expansion of the universe observed at large scales and the absence of any such expansion observable in the solar system. As shown in the previous section since the Swiss cheese model has a rather severe and very unphysical deficiency, we think it is time to begin the search for an acceptable metric.
The non-zero elements of the extrinsic curvature is given by (2.11) as
$`K_{22}`$ $`=`$ $`a(t)r{\displaystyle \frac{[1(r/2R)^2]}{[1+(r/2R)^2]^2}}[1(m/2a(t)r)^2]`$ (72)
$`K_{33}`$ $`=`$ $`\mathrm{sin}^2\theta K_{22}`$ (73)
$`K_{44}`$ $`=`$ $`{\displaystyle \frac{mc^2(1m/2a(t)r)}{a(t)^2r^2(1+m/2a(t)r)^5}}`$ (74)
In the limit that $`m0`$ these curvatures reduce to those of eq. 35 and in the limit $`R\mathrm{}`$ and $`\dot{a}=0`$ ($`a=1`$) they reduce to
$`K_{22}`$ $`=`$ $`r\left[1\left({\displaystyle \frac{m}{2r}}\right)^2\right],K_{33}=\mathrm{sin}^2\theta K_{22},`$ (75)
$`K_{44}`$ $`=`$ $`{\displaystyle \frac{mc^2(1m/2r)}{r^2(1+m/2r)^5}}`$ (76)
These curvatures agree with those of the unreparameterized Schwarzschild metric. The curvatures in eq. 73 differ only in that $`r`$ is replaced by $`a(t)r`$, $`R`$ by $`a(t)R`$, and there are corrections for the overall curvature of space. Thus, it reflects, as does the metric, the very same behavior, in terms of $`a(t)r`$ instead of $`r`$ as was found by Schwarzschild for his metric.
Both the metric and the extrinsic curvature are continuous outside the Schwarzschild radius, which is necessary for a metric to be acceptable.
We now turn to the stress-energy tensor. The current metric shares with both the Friedmann-Lemaître metric, eq. 29, and the Schwarzschild metric, eq. 37, the property that $`T_1^4=T_4^1=0`$, so there is no mass-density flux. This property follows by direct computation from eq. 13 and eq. 70. For the other non-zero components, we find
$`8\pi T_1^1`$ $`=`$ $`{\displaystyle \frac{1+(m/2a(t)r))^2}{a(t)^2R^2(1+m/2ar)^5(1m/2a(t)r)}}`$ (78)
$`+{\displaystyle \frac{2\ddot{a}a[1+m/2a(t)r)]+3\dot{a}^2[1m/2a(t)r]}{a(t)^2c^2(1m/2a(t)r)}}\mathrm{\Lambda }`$
$`8\pi T_2^2`$ $`=`$ $`8\pi T_3^3=8\pi T_1^1`$ (79)
$`8\pi T_4^4`$ $`=`$ $`{\displaystyle \frac{3}{a(t)^2R^2(1+m/2a(t)r)^5}}+{\displaystyle \frac{3\dot{a}^2}{c^2a^2}}\mathrm{\Lambda },`$ (80)
In the limit as $`m0`$ or $`r\mathrm{}`$ these tensor elements reduce to those of eq. 33, and also in the limit as $`R\mathrm{}`$ and $`\dot{a}0`$ they vanish as with the Schwarzschild metric, unless $`\mathrm{\Lambda }0`$, in which case the result (3.7) is obtained. The dominate terms in the diagonal elements (pressure and density) are a sum of terms of second order in $`1/R`$ and terms containing two time derivatives of the universal expansion factor $`a(t)`$. The metric given by eqs. 70 and 8 is the solution of Einstein’s field equations eq. 4 when the stress-energy tensor, eq. 79, is specified.
## VI Dynamics
In this section we investigate the equations of motion a test particle in our proposed metric as described by eq. 8 and eq. 70. Our treatment diverges from that of McVittie at this point as he directly substitutes $`\rho =a(t)r`$ into the line element rather than using the transformation equations
$$\overline{g}_{kl}=\frac{x^i}{\overline{x}^k}\frac{x^j}{\overline{x}^l}g_{ij},$$
(81)
The correct change of the line element for this change of variables is
$`ds^2`$ $`=`$ $`e^U\left[(d\rho )^2+\rho ^2(d\theta )^2+\rho ^2\mathrm{sin}^2\theta (d\varphi )^2\right]`$ (83)
$`+\left[e^Ve^U\left({\displaystyle \frac{\rho \dot{a}}{a}}\right)^2\right](d\tau )^2+2e^U\left({\displaystyle \frac{\rho \dot{a}}{a}}\right)d\rho d\tau ,`$
where
$`e^U`$ $`=`$ $`{\displaystyle \frac{(1+m/2\rho )^4}{[1+(\rho /2a(\tau )R)^2]^2}}`$ (84)
$`e^V`$ $`=`$ $`c^2\left({\displaystyle \frac{1m/2\rho }{1+m/2\rho }}\right)^2`$ (85)
Clearly these are non-synchronous coordinates as the coefficient of $`d\rho d\tau `$ is non-zero.
We will make this change of variables later after the equations of motion have been derived. The main difference is that McVittie omits the $`\dot{\mu }`$ terms from his subsequent equations. The equations of motion then become, using eq. 24, for $`\theta `$ and $`\varphi `$,
$`\left({\displaystyle \frac{ds}{dt}}\right){\displaystyle \frac{d}{dt}}\left[\left({\displaystyle \frac{ds}{dt}}\right)^1e^\mu r^2\dot{\theta }\right]`$ $`=`$ $`e^\mu r^2\mathrm{sin}\theta \mathrm{cos}\theta \dot{\varphi }^2`$ (86)
$`\left({\displaystyle \frac{ds}{dt}}\right){\displaystyle \frac{d}{dt}}\left[\left({\displaystyle \frac{ds}{dt}}\right)^1e^\mu r^2\mathrm{sin}^2\theta \dot{\varphi }\right]`$ $`=`$ $`0`$ (87)
One can see by inspection, the motion in the plane $`\theta =\pi /2`$ is a solution. The intergal of the second equation gives the result,
$$r^2\dot{\varphi }=Ae^\mu \left(\frac{ds}{dt}\right),$$
(88)
where $`A`$ is a constant of integration. This equation is the conservation of angular momentum in this coordinate system. To obtain the equation of motion for $`r`$, it is convenient to use eq. 27. We obtain,
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{ds}{dt}}\right)^2{\displaystyle \frac{d}{dt}}\left[\left({\displaystyle \frac{ds}{dt}}\right)^2\right]e^\mu \dot{r}{\displaystyle \frac{d}{dt}}\left[e^\mu \dot{r}\right]=`$ (89)
$`{\displaystyle \frac{1}{2}}[\mu ^{}e^\mu \dot{r}^2+(\mu ^{}r^2+2r)e^\mu \dot{\theta }^2`$ (90)
$`+(\mu ^{}r^2+2r)e^\mu \mathrm{sin}^2\theta \dot{\varphi }^2\nu ^{}e^\nu ],`$ (91)
where, from the line element,
$$\left(\frac{ds}{dt}\right)^2=e^\mu \left[\dot{r}^2+r^2\dot{\theta }^2+r^2\mathrm{sin}^2\theta \dot{\varphi }^2\right]+e^\nu .$$
(92)
At this point, we simplify to motion in the $`\theta =\pi /2`$ plane. Thus $`\dot{\theta }=\ddot{\theta }=0`$, and eq. 91 becomes,
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{ds}{dt}}\right)^2{\displaystyle \frac{d}{dt}}\left[\left({\displaystyle \frac{ds}{dt}}\right)^2\right]\dot{r}+e^\mu {\displaystyle \frac{d}{dt}}\left[e^\mu \dot{r}\right]=`$ (93)
$`{\displaystyle \frac{1}{2}}\left[\mu ^{}\dot{r}^2+(\mu ^{}r^2+2r)\dot{\varphi }^2\nu ^{}e^{\nu \mu }\right].`$ (94)
These formulas yield the equations of motion of a test particle in the reference frame which is at rest with respect to the coordinate system at the location of the test particle. We want to find the motion with respect to an observer at rest at $`r=0`$. To this end we introduce the change of variables,
$$r=\frac{\rho }{a(t)},\dot{r}=\frac{\dot{\rho }}{a(t)}\frac{\dot{a}(t)\rho }{a(t)^2}$$
(95)
This change of variables yields coordinates which are equal to the proper distances, as viewed from the origin, when the mass concentration is absent. With the substitution 95 we obtain for eq. 88 and eq. 92
$`\rho ^2\dot{\varphi }`$ $`=`$ $`Aa(t)^2e^\mu \left({\displaystyle \frac{ds}{dt}}\right),`$ (96)
$`\left({\displaystyle \frac{ds}{dt}}\right)^2`$ $`=`$ $`{\displaystyle \frac{e^\mu }{a(t)^2}}\left[\left(\dot{\rho }{\displaystyle \frac{\dot{a}(t)\rho }{a(t)}}\right)^2+\rho ^2\dot{\varphi }^2\right]+e^\nu .`$ (97)
Finally, eq. 94 becomes,
$`\ddot{\rho }2{\displaystyle \frac{\dot{a}}{a}}\dot{\rho }+\rho \left[2\left({\displaystyle \frac{\dot{a}}{a}}\right)^2{\displaystyle \frac{\ddot{a}}{a}}\right]+\dot{\mu }\left(\dot{\rho }{\displaystyle \frac{\dot{a}}{a}}\rho \right)`$ (98)
$`=`$ $`{\displaystyle \frac{1}{2}}\mu ^{}a\left({\displaystyle \frac{\dot{\rho }}{a}}{\displaystyle \frac{\dot{a}}{a^2}}\rho \right)^2+{\displaystyle \frac{1}{2}}\left(\mu ^{}{\displaystyle \frac{\rho }{a}}+2\right)\left[{\displaystyle \frac{A^2a^4e^{2\mu }}{\rho ^3}}\right]\left({\displaystyle \frac{ds}{dt}}\right)^2`$ (100)
$`{\displaystyle \frac{1}{2}}\nu ^{}ae^{\nu \mu }{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{ds}{dt}}\right)^2{\displaystyle \frac{d}{dt}}\left[\left({\displaystyle \frac{ds}{dt}}\right)^2\right]\left(\dot{\rho }{\displaystyle \frac{\dot{a}}{a}}\rho \right)`$
where we have used eq. 96 to eliminate the $`\dot{\varphi }`$ dependence. There is also a $`\dot{\varphi }`$ in $`ds/dt`$ but it too can be eliminated by the substitution of $`\dot{\varphi }`$ from eq. 96 in eq. 97. To assess the importance of the various terms it is helpful to note the following dimensionless quantities, in “planetary units,”
$`T_{}H_0`$ $``$ $`5\times 10^{11},\left({\displaystyle \frac{v_{}}{c}}\right)^21\times 10^8,`$ (101)
$`{\displaystyle \frac{GM_{}}{c^2R_{}}}`$ $``$ $`1\times 10^8,\left({\displaystyle \frac{R_{}}{R_{\mathrm{Hubble}}}}\right)^20.6\times 10^{30},`$ (102)
where $`T_{},v_{},R_{}`$ are the orbital period, velocity, and radius of the earth, and $`M_{}`$ is the mass of the sun. Some limiting cases are of interest. First, we take the flat-space, slow-speed limit, i.e., $`R=\mathrm{}`$, and
$$\left(\frac{ds}{dt}\right)^2c^2\left(\frac{1{\displaystyle \frac{GM}{2c^2\rho }}}{1+{\displaystyle \frac{GM}{2c^2\rho }}}\right)^2.$$
(103)
Thus eq. 98 reduces to
$`\ddot{\rho }{\displaystyle \frac{\ddot{a}}{a}}\rho \left(\dot{\rho }{\displaystyle \frac{\dot{a}}{a}}\rho \right){\displaystyle \frac{\dot{a}m(3m/\rho )}{a\rho \left[1(m/2\rho )^2\right]}}`$ (104)
$`=`$ $`{\displaystyle \frac{m(\dot{\rho }\dot{a}\rho /a)^2}{\rho ^2(1+m/2\rho )}}+{\displaystyle \frac{A^2c^2(1m/2\rho )^3}{\rho ^3(1+m/2\rho )^{11}}}{\displaystyle \frac{mc^2(1m/2\rho )}{\rho ^2(1+m/2\rho )^7}}`$ (105)
We may further reduce these equations by discarding terms which are proportional to $`c^2`$ for when the velocities are much less than the speed of light. By eq. 70 these are the terms proportional to m alone, but we retain, of course, the terms in $`mc^2`$. Eq. 104 reduces further to
$$\ddot{\rho }\frac{\ddot{a}}{a}\rho =\frac{A^2c^2}{\rho ^3}\frac{GM}{\rho ^2}=\frac{GM}{\rho ^2}\left(\frac{\rho _0}{\rho }1\right).$$
(106)
The coefficient of the discarded term on the left-hand side of eq. 104 is of the order of $`10^{19}`$ as it is the product of two first order corrections. The discarded (first) term on the right-hand side of eq. 104 is smaller by a factor of $`\dot{\rho }^2/c^2`$ than the last term. The other items discard are factors of $`GM/c^2\rho `$ smaller than the dominant terms.
It is to be noticed that in the limit of eq. 106, that it differs from Newton’s equation of gravitation only by a term proportional to $`\ddot{a}`$. Note is taken the Noerdlinger and Petrosian do take account of this term. The constant $`\rho _0`$ is just another form of the constant of integration $`A`$. The magnitude of the $`\ddot{a}`$ term equals that for the sun’s gravity at about 0.5 kiloparsecs. The effects on the scale of the solar systems are too small to be measured.
## VII Examples
The equations of motion in a flat, Friedmann-Lemaître expanding universe for a slowly moving test particle under no external forces are, by eq. 106 and the corresponding reduction of eq. 87,
$$\ddot{\rho }\rho \dot{\varphi }^2=\frac{\ddot{a}}{a}\rho =\frac{\psi (1\psi )}{t^2}\rho ,\frac{d}{dt}\left(\rho ^2\dot{\varphi }\right)=0,$$
(107)
for the standard form of the universal expansion factor for the universe, as describe at the end of section III. $`t`$ is the current age of the universe. One easily recognizes these equations to be just exactly Newton’s equations in a plane in spherical coordinates. If we change to rectangular coordinates, we get exactly eq. 62. The behavior of the solutions is discussed in detail in Section IV above.
Next we consider the case where we add gravitational effects to their leading order. For purely radial motion, the $`A`$ of eq. 96 is zero. Thus the equation of motion 106 becomes,
$$\ddot{\rho }=\frac{\ddot{a}(t)}{a(t)}\rho \frac{GM}{\rho ^2}$$
(108)
This equation differs from the Schwarzschild metric equation, 60 by the addition of a term in $`\ddot{a}`$. We will consider this behavior over time periods short compared to the age of the universe. If we multiply by $`d\rho `$ and integrate, we get,
$$\frac{1}{2}\dot{\rho }^2=\frac{\ddot{a}}{2a}\rho ^2+\frac{GM}{\rho }+E_0$$
(109)
where $`E_0`$ is the constant of integration. If $`\ddot{a}=0`$, then $`E_0=0`$ would correspond to a test particle which had zero velocity at $`\rho =\mathrm{}`$. I choose to examine this special case. Then eq. 109 becomes
$`{\displaystyle _{\rho _0}^\rho }{\displaystyle \frac{\xi ^{1/2}d\xi }{2\sqrt{GM+\frac{\ddot{a}}{2a}\xi ^3}}}`$ $`=`$ $`tt_0.`$ (110)
$`{\displaystyle \frac{2}{3}}\sqrt{{\displaystyle \frac{a}{\ddot{a}}}}\mathrm{sin}^1\left({\displaystyle \frac{\xi ^{3/2}}{\sqrt{2GMa/\ddot{a}}}}\right)|_{\rho _0}^\rho `$ $`=`$ $`tt_0,`$ (111)
by Pierce’s tables. Thus,
$`\rho ^{3/2}`$ $`=`$ $`\rho _0^{3/2}+\sqrt{{\displaystyle \frac{2GMa}{\ddot{a}}}}\left[\mathrm{sin}\left({\displaystyle \frac{3}{2}}(tt_0)\sqrt{{\displaystyle \frac{\ddot{a}}{a}}}\right)\right]`$ (112)
$``$ $`\rho _0^{3/2}+{\displaystyle \frac{3}{2}}\sqrt{2GM}\left[(tt_0)+{\displaystyle \frac{3\ddot{a}}{8a}}(tt_0)^3+\mathrm{}\right].`$ (113)
It is evident from this solution that the leading order corrections due to the expansion of the universe are of the order $`H_0^2(tt_0)^2`$ which is extremely small on the planetary time scale. The solution when $`(tt_0)=O(t_0)`$ would take account of the time dependence of $`\ddot{a}/a`$. The more general case, where $`E_00`$, can also be integrated in terms of elliptic functions of the first and third kinds.
Next we investigate bound circular motion. To do so, I set $`\ddot{\rho }=0`$ in eq. 106 which gives,
$$\frac{A^2c^2}{\rho ^3}=\frac{GM}{\rho ^2}+\frac{\ddot{a}}{a}\rho .$$
(114)
If I use eq. 96 to reintroduce $`\dot{\varphi }`$, and remember that the period $`T=2\pi /\dot{\varphi }`$, then I find,
$$\frac{4\pi ^2\rho ^3}{GMT^2}=1+\frac{\ddot{a}\rho ^3}{GMa},$$
(115)
which is Kepler’s law relating the square of the period to the cube of the radius, with a correction caused by the expansion of space. For the case $`a(t)t^{2/3}`$ the correction term becomes, $`\frac{1}{2}H_0^2\rho ^3/GM`$. If we use the solar mass, then
$$\frac{4\pi ^2\rho ^3}{GMT^2}=13.307\times 10^{23}h_{50}^2\rho ^3,$$
(116)
where $`\rho `$ is in astronomical units and $`h_{50}=1`$ when $`H_0=50`$ km per second per megaparsec.
These results are illustrated in Fig. 2. In flat space at time $`t_0`$, the interface between the static Schwarzschild metric and the non-static Friedmann-Lemaître metric is at distance unity from the origin, in the units of section IV. We display the Schwarzschild result when started with unit angular velocity ($`\dot{\varphi }=1`$). It is a circle of radius unity. Next we display the trajectory just outside the interface. It is, as expected, a parabolic curve. I have started it with unit angular velocity (again $`\dot{\varphi }=1`$) at a distance of 1.25 from the mass concentration. It is to be noted that after a unit (Hubble) time has passed, the two trajectories are significantly separated.
The reason for starting it further out is, as explained in section IV, that if I were to have started it at a distance between 1.0 and $`\sqrt{13}/31.2018504`$, the trajectory would have been over taken by the expanding spherical interface between the two metrics of the “Swiss cheese model.”
In addition I show the trajectory using the presently considered metric. I have again started with unit distance and unit angular velocity. The result here is that it converges fairly quickly to an elliptical trajectory. In this case I find, using the standard equations,
$$v^2=GM\left(\frac{2}{r}\frac{1}{𝒜}\right),\frac{1}{2}rv^2=A=\frac{1}{2}\sqrt{\frac{GM}{𝒜}},$$
(117)
where $`𝒜,B`$ are the major and minor semi-axes, and the areal velocity $`A`$ is a constant by the conservation of angular momentum. In this case I find $`𝒜1.0437`$ and $`1.0216`$. In as much as the semilatus rectum parameter $`p=1`$, the latus rectum itself is clearly defined by the intersection of this ellipse with the unit circle (Schwarzschild trajectory). The latus rectum is the line through the focus (origin in this case) which is perpendicular to the semimajor axis. The eccentricity is $`e0.20462`$.
###### Acknowledgements.
The author is pleased to acknowledge helpful conversations with S. Habib, P. O. Mazur, E, Motolla, and M. M. Nieto. |
warning/0003/quant-ph0003131.html | ar5iv | text | # Disagreement between correlations of quantum mechanics and stochastic electrodynamics in the damped parametric oscillator
## I Introduction
Local hidden variable theories have been extensively compared to quantum mechanics over the last seventy or so years . Most comparisons between the two have investigated whether or not quantum mechanics is equivalent to a local hidden variable theory. Much evidence indicates that it is not. Many results in quantum mechanics have been found that are incompatible with all local hidden variable theories . Most of these results have involved idealized undamped systems. However, all experimental systems encounter damping. Thus, it is interesting (and more realistic) to compare quantum mechanics and local hidden variable theories in damped systems . This paper compares quantum mechanics and one local hidden variable theory (SED) in such a system.
One of the earliest works comparing local hidden variable theories to quantum mechanics was Bell’s theorem . It demonstrates that quantum mechanics is incompatible with all local hidden variable theories at a statistical level. It does so by deriving an upper bound on a function of two particle correlations for all local hidden variable theories, which quantum mechanics exceeds. Extensions of it have been formulated for large angular momentum and particle number systems . These extensions demonstrate nonclassical behavior in a regime usually regarded as being purely classical. Greenberger, Horne and Zeilinger (GHZ) have also extended Bell’s work, differentiating quantum mechanics from all local hidden variable theories for single, as opposed to ensemble, measurements. The three particle GHZ theorem has an “all or nothing” quality and distinguishes between local hidden variable theories and quantum mechanics in a single experimental run, once three basic correlations are established.
Comparisons between quantum mechanics and local hidden variable theories have also been made using continuous variables (which are discretized in formulating the comparison), such as quadrature phase amplitudes and there is currently much interest in this area. For quadrature phase amplitude measurements, these comparisons can have detector efficiencies of in excess of 99% . They also tend to relate more strongly to Einstein, Podolsky and Rosen’s original EPR paradox than earlier discrete variable ones. Indeed, the EPR paradox has been experimentally demonstrated using quadrature phase amplitudes . Additionally, quantum teleportation has been achieved using quadrature phase amplitudes further demonstrating the utility of continuous variables.
One commonly used local hidden variable theory is stochastic electrodynamics (SED) . Some authors have proposed it as an alternative to quantum mechanics . Furthermore, a semiclassical approach equivalent to it is also commonly used in parametric oscillator calculations . SED consists of adding Gaussian white noise to classical electrodynamics. It is equivalent to truncating third order derivative terms in the quantum mechanical Moyal equation, a commonly used approximation . Such terms are often negligible and thus SED reproduces many results of quantum mechanics . However, it cannot violate Bell inequalities for quadrature phase amplitude measurements, and is thus distinct from quantum mechanics . Various authors have explicitly shown differences between SED and quantum mechanics . In particular, it has been shown that the two theories predict different transient third order correlations for the undamped nondegenerate parametric oscillator . It has also been shown that they predict different macroscopic quadrature phase amplitude correlations in the damped nondegenerate parametric oscillator in the steady state .
In general, differences between quantum mechanics and local hidden variable theories are reduced or eliminated by damping . Furthermore, damping is a significant element of many realistic systems. It is thus important to consider its effects on differences between quantum mechanics and local hidden variable theories such as SED. However, all but a few of the comparisons between the quantum mechanics and local hidden variable theories referenced above have involved undamped systems. They are thus idealized in this respect. In contrast, damping is included in the calculatons in this paper. It is included to consider a theoretical model which is as realistic as possible and also to determine the sensitivity of differences between quantum mechanics and SED to its presence.
This paper extends a previous comparison between quantum mechanics and SED in the nondegenerate parametric oscillator . In particular, it contrasts both intracavity and external moments of the two theories’ in the same system with damping included. Expressions from both theories are compared for the intracavity moment $`\mathrm{\Delta }X_1(\tau )\mathrm{\Delta }(X_2(\tau )\mathrm{\Delta }X_3(\tau )`$, where $`\mathrm{\Delta }X_i(\tau )=X_i(\tau )X_i(\tau )`$, for $`i=1,2,3`$, $`X_i(\tau )`$ is quadrature phase amplitude, the subscripts represent different radiation modes and $`\tau `$ is a scaled time variable. A comparison is also made for an analogous external moment. Both analytic iterative and numerical techniques are used to calculate moments. The results produced by these techniques show that the intracavity and external moments differ greatly between the two theories. In particular, the analytic method shows that the moments of quantum mechanics are cubic in the system’s nonlinear coupling constant to leading order whilst those of SED are linear. The two theories are compared over a range of nonlinear coupling constant, damping and average initial pump photon number values. The results of these comparisons show a number of qualitative trends. Most importantly, quantum mechanics and SED differ in the situations considered with the largest particle number and damping to nonlinear coupling ratios, although the differences are reduced in relative size.
Stochastic techniques are used to obtain results both for quantum mechanics and SED. The positive-P coherent state representation is used to calculate quantum mechanical predictions. It is particularly well suited to the calculation of quantum dynamics in damped quantum optical systems when nonclassical behavior is present. It is able to handle arbitrarily large photon numbers. It converges quickly (in the sense of sampling error) when systems’ dimensionless nonlinearities are relatively small, as is the case with nonlinear optical experiments. By contrast, the method used for SED calculations corresponds to commonly used approaches in quantum optics, where the field is treated as a semiclassical object surrounded by (classical) vacuum fluctuations. Both methods are used to generate analytic predictions and are also numerically simulated.
## II Quantum mechanics
This paper considers an idealized nondegenerate parametric oscillator, resonant at three frequencies $`\omega _1\mathrm{and}\omega _2`$ (signal and idler frequencies) and $`\omega _3=\omega _1+\omega _2`$ (pump frequency). It contains a nonlinear medium that couples the modes and converts higher energy pump photons into lower energy signal and idler ones. The system’s interaction Hamiltonian, including linear losses, is given by
$$\widehat{H}=i\mathrm{}G(\widehat{a}_1^{}\widehat{a}_2^{}\widehat{a}_3\widehat{a}_1\widehat{a}_2\widehat{a}_3^{})+\underset{i=1}{\overset{3}{}}\widehat{\mathrm{\Gamma }}_i\widehat{a}_i^{}+\widehat{\mathrm{\Gamma }}_i^{}\widehat{a}_i,$$
(1)
where $`\widehat{a}_i^{}\mathrm{and}\widehat{a}_i`$ are creation and annihilation operators for oscillator modes, $`\widehat{\mathrm{\Gamma }}_i^{}\mathrm{and}\widehat{\mathrm{\Gamma }}_i`$ are environment mode operators and G is a nonlinear interaction strength constant. Initially, the system has a coherent state in the pump mode, and vacuum states in the signal and idler modes.
A number of quasiprobability representations exist to describe quantum states, the most famous being the Glauber-Sudarshan representation. It is produced by decomposing quantum density operators using a diagonal coherent state basis. Thus,
$$\widehat{\rho }=𝑑\alpha ^2P(\alpha ,\alpha ^{})|\alpha \alpha |,$$
(2)
where $`\widehat{\rho }`$ is a density operator and $`P(\alpha ,\alpha ^{})`$ is the Glauber-Sudarshan representation. The Glauber-Sudarshan representation can be negative and is hence not a strict probability density function. A more recent representation is the positive-P representation which is an actual probability density function over an off-diagonal coherent state basis. It further differs from the Glauber-Sudarshan representation by using a phase space of doubled dimension. The positive-P variables $`\{\alpha _i,\alpha _i^+\}`$, where $`i`$ is a positive integer, are analogous to complex field amplitudes, with $`\alpha _i`$ and $`\alpha _i^+`$ describing a particular radiation mode. However, $`\{\alpha _i\}`$ and $`\{\alpha _i^+\}`$ are independent and hence $`\alpha _i(\alpha _i^+)^{}`$, though their averages are complex conjugate and thus $`\alpha _i=\alpha _i^+^{}`$. Variable averages are equal to normally ordered quantum averages once the substitutions $`\alpha _i\widehat{a}_i`$ and $`\alpha _i^+\widehat{a}_i^{}`$ are made. For example, $`\alpha _1\alpha _1^+=\widehat{a}_1^{}\widehat{a}_1_\rho `$, where $`\widehat{O}_\rho `$ denotes $`Tr(\widehat{\rho }\widehat{O})`$, as usual in quantum mechanics.
Stochastic equations of motion for positive-P variables for the damped nondegenerate parametric oscillator are, in terms of $`\tau `$ (time scaled by $`\mathrm{\Gamma }`$, a typical damping constant with units of inverse time),
$`{\displaystyle \frac{d\alpha _1}{d\tau }}`$ $`=`$ $`\gamma _1\alpha _1+g\alpha _2^+\alpha _3+\sqrt{g\alpha _3}\xi _1`$ (3)
$`{\displaystyle \frac{d\alpha _1^+}{d\tau }}`$ $`=`$ $`\gamma _1\alpha _1^++g\alpha _2\alpha _3^++\sqrt{g\alpha _3^+}\xi _1^+`$ (4)
$`{\displaystyle \frac{d\alpha _2}{d\tau }}`$ $`=`$ $`\gamma _2\alpha _2+g\alpha _1^+\alpha _3+\sqrt{g\alpha _3}\xi _2`$ (5)
$`{\displaystyle \frac{d\alpha _2^+}{d\tau }}`$ $`=`$ $`\gamma _2\alpha _2^++g\alpha _1\alpha _3^++\sqrt{g\alpha _3^+}\xi _2^+`$ (6)
$`{\displaystyle \frac{d\alpha _3}{d\tau }}`$ $`=`$ $`\gamma _3\alpha _3g\alpha _1\alpha _2`$ (7)
$`{\displaystyle \frac{d\alpha _3^+}{d\tau }}`$ $`=`$ $`\gamma _3\alpha _3^+g\alpha _1^+\alpha _2^+.`$ (8)
Here $`\xi _1,\xi _2,\xi _1^+\mathrm{and}\xi _2^+`$ are complex Gaussian white noises with the following correlations
$`\xi _i(\tau _1)\xi _j(\tau _2)`$ $`=`$ $`\delta _{3i,j}\delta (\tau _1\tau _2)`$ (9)
$`\xi _i^+(\tau _1)\xi _j^+(\tau _2)`$ $`=`$ $`\delta _{3i,j}\delta (\tau _1\tau _2)`$ (10)
$`\xi _i^+(\tau _1)\xi _j(\tau _2)`$ $`=`$ $`0,`$ (11)
where $`i,j=1,2`$. In Eq. (3), $`\gamma _i=\mathrm{\Gamma }_i/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }_i`$ is a damping constant for mode $`i`$ with units of inverse time, $`g=G/\mathrm{\Gamma }`$ and $`\tau =\mathrm{\Gamma }t`$. It is assumed that G, $`\mathrm{\Gamma }_i`$ and $`\mathrm{\Gamma }`$ are real. Initial conditions are $`\alpha _1(0)=0,\alpha _2(0)=0\mathrm{and}\alpha _3(0)=ϵ`$. It is noted that Eq. (3) is only valid when boundary terms in phase space can be neglected. These are asymptotically small in the limit of short times or large damping ratios .
Eq. 3 is solved using an analytic iterative method. This method treats damping terms exactly, and noise and nonlinear terms iteratively. It involves, firstly, rewriting the equations forming Eq. (3) as $`\dot{\alpha }_i=\gamma _i\alpha _i+f_i(\{\alpha _j,\alpha _j^+\},\tau )`$ or $`\dot{\alpha }_i^+=\gamma _i\alpha _i^++f_i^+(\{\alpha _j,\alpha _j^+\},\tau )`$, where $`i,j=1,2,3`$. Successively higher order approximations for $`\{\alpha _i(\tau ),\alpha _i^+(\tau )\}`$ are then found using increasingly better approximations for $`f_i`$ and $`f_i^+`$. Thus, $`(m+1)^{th}`$ order terms are given by
$`\alpha _i^{(m+1)}(\tau )`$ $`=`$ $`\alpha _i^{(0)}(\tau )+{\displaystyle _{\tau _1=0}^{\tau _1=\tau }}𝑑\tau _1\mathrm{exp}[\gamma _i(\tau _1\tau )]f_i(\{\alpha _j^{(m)},\alpha _j^{+(m)}\},\tau _1)`$ (12)
$`\alpha _i^{+(m+1)}(\tau )`$ $`=`$ $`\alpha _i^{+(0)}(\tau )+{\displaystyle _{\tau _1=0}^{\tau _1=\tau }}𝑑\tau _1\mathrm{exp}[\gamma _i(\tau _1\tau )]f_i^+(\{\alpha _j^{(m)},\alpha _j^{+(m)}\},\tau _1),`$ (13)
where $`\alpha _k^{(0)}(\tau )=\alpha _k(\tau =0)\mathrm{exp}(\gamma _k\tau )`$ and $`\alpha _k^{(0)+}(\tau )=\alpha _k^{(0)}(\tau )^{}`$ where $`k=1,2,3`$. For example,
$$\alpha _1^{(m+1)}(\tau )=\alpha _1^{(0)}(\tau )+_{\tau _1=0}^{\tau _1=\tau }𝑑\tau _1\mathrm{exp}[\gamma _1(\tau _1\tau )]\left(g\alpha _2^{+(m)}(\tau _1)\alpha _3^{(m)}(\tau _1)+\sqrt{g\alpha _3^{(m)}(\tau _1)}\xi _1(\tau _1)\right)$$
(14)
and first order approximations are
$`\alpha _i^{(1)}(\tau )`$ $`=`$ $`{\displaystyle _{\tau _i=0}^{\tau _i=\tau }}𝑑\tau _i\mathrm{exp}[\gamma _i(\tau _i\tau )]\sqrt{gϵ}\mathrm{exp}({\displaystyle \frac{\gamma _i\tau _i}{2}})\xi _i(\tau _1)`$ (15)
$`\alpha _3^{(1)}(\tau )`$ $`=`$ $`ϵ\mathrm{exp}(\gamma _3\tau )`$ (16)
$`\alpha _i^{+(1)}(\tau )`$ $`=`$ $`{\displaystyle _{\tau _i=0}^{\tau _i=\tau }}𝑑\tau _i\mathrm{exp}[\gamma _i(\tau _i\tau )]\sqrt{gϵ^{}}\mathrm{exp}({\displaystyle \frac{\gamma _i\tau _i}{2}})\xi _i^+(\tau _1)`$ (17)
$`\alpha _3^{+(1)}(\tau )`$ $`=`$ $`ϵ^{}\mathrm{exp}(\gamma _3\tau ),`$ (18)
where $`i=1,2`$.
## III Stochastic diagrams
The iterative method of the previous section can be used, in conjunction with stochastic diagrams, to readily produce analytic approximations for the intracavity moments of quantum mechanics considered in this paper. Stochastic diagrams are schematic representations of the combinatoric parts of an iterative process. They clearly lay out all terms produced by different orders of iteration. Fundamental stochastic diagrams appear as one of three classes. Those associated with initial conditions appear as straight lines, those with noise terms as straight lines with a cross at their end and those with nonlinear terms as straight lines containing a fork, as shown in Fig. 1(a)-(c). Higher order iterative terms are represented by stochastic diagrams using either combinations of the three basic classes. For example, one of the iterative terms in $`\alpha _1^{(2)}(\tau )`$ is
$$_{\tau _1=0}^{\tau _1=\tau }𝑑\tau _1\mathrm{exp}[\gamma _1(\tau _1\tau )]g\alpha _3^{(0)}(\tau _1)_{\tau _2=0}^{\tau _2=\tau _1}𝑑\tau _2\mathrm{exp}[\gamma _2(\tau _2\tau _1)]\sqrt{gϵ^{}}\mathrm{exp}(\frac{\gamma _2\tau _2}{2})\xi _2^+(\tau _2).$$
It combines all three basic classes and is represented by the stochastic diagram in Fig. 1(d). All iterative terms can be represented by stochastic diagrams.
Stochastic diagrams can also be used to determine the orders of iterative terms. In particular, they can be used to determine the orders of such terms in the system’s nonlinear coupling constant g. This paper focuses on the order of terms in this constant. For quantum mechanics, initial value iterative terms are $`O(g^0)`$, noise iterative terms $`O(g^{1/2})`$ and nonlinear iterative terms $`O(g)`$. Hence, lines in stochastic diagrams count as order zero, crosses as order 1/2 and vertices as order 1. A term’s order is simply found by considering its stochastic diagram and adding a half to its order for every cross and one for every vertex. For example, the term represented in Fig. 1(d) has one vertex and one cross and thus is $`O(g^{\frac{3}{2}})`$. A notation that denotes the order in g of a term by a superscript \[n\] is used in this section.
Stochastic diagrams are now used to determine the intracavity moments of quantum mechanics considered in this paper. Consider all eight moments of the form $`\mathrm{\Delta }𝒜_1(\tau )\mathrm{\Delta }𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$, where $`\mathrm{\Delta }𝒜_\mathrm{i}(\tau )=𝒜_i(\tau )𝒜_i(\tau )`$ and $`𝒜_i(\tau )`$ is either $`\widehat{a}_i`$ or $`\widehat{a}_i^{}`$. These are equal to the positive-P variable moments which replace $`\widehat{a}_i`$ and $`\widehat{a}_i^{}`$ by $`\alpha _i(\tau )`$ and $`\alpha _i^+(\tau )`$ respectively. Now, consider the equations that constitute Eq. (3). Their forms do not change when they are expressed in terms of $`\alpha _i(\tau ),\alpha _i^+(\tau ),\alpha _3(\tau )`$ and $`\alpha _3^+(\tau )`$, where $`i=1,2`$. From this, it follows that $`𝒜_i(\tau )=𝒜_i(\tau )`$ and hence $`\alpha _i(\tau )=\alpha _i^+(\tau )=0`$, where again $`i=1,2`$. Thus, $`\mathrm{\Delta }𝒜_1(\tau )\mathrm{\Delta }𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$ can be simplified to $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$.
An approximate expression for $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$ is now obtained using the iterative method in Section II (and stochastic diagrams). This method can be used to produce power series expressions in $`g`$ for the positive-P variables. These expressions can then be used to generate power series expressions in $`g`$ for the moments of the form $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$. As $`g1`$ in realistic systems, these power series expressions can be approximated by their lowest order nonzero terms.
Figs 2 (a) and (b) show the stochastic diagrams required to determine the moments of the form $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$. Naively, it might be thought that the lowest order nonzero terms from $`𝒜_1(\tau ),𝒜_2(\tau )\mathrm{and}\mathrm{\Delta }𝒜_3(\tau )`$ simply need to be multiplied together and the average of the subsequent product determined to calculate the lowest order nonzero term in $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$. This is not always true. Sometimes, $`𝒜_1(\tau ),𝒜_2(\tau )\mathrm{and}\mathrm{\Delta }𝒜_3(\tau )`$ are not necessarily zero and yet $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$ is zero. For example, the lowest order nonzero terms for the positive-P variables in $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3^+(\tau )`$ are
$$\alpha _i^{[1/2]}(\tau )=_{\tau _i=0}^{\tau _i=\tau }𝑑\tau _i\mathrm{exp}[\gamma _i(\tau _i\tau )]\sqrt{gϵ}\mathrm{exp}(\frac{\gamma _i\tau _i}{2})\xi _i(\tau _i),$$
(19)
where $`i=1,2`$ and
$`\mathrm{\Delta }\alpha _3^{+[2]}(\tau )`$ $`=`$ $`{\displaystyle _{\tau _3=0}^{\tau _3=\tau }}𝑑\tau _3\mathrm{exp}[\gamma _3(\tau _3\tau )]g{\displaystyle _{\tau _4=0}^{\tau _4=\tau _3}}𝑑\tau _4\mathrm{exp}[\gamma _1(\tau _4\tau _3)]\sqrt{gϵ^{}}\mathrm{exp}({\displaystyle \frac{\gamma _1\tau _4}{2}})\xi _1^+(\tau _4)`$ (23)
$`{\displaystyle _{\tau _5=0}^{\tau _5=\tau _3}}𝑑\tau _5\mathrm{exp}[\gamma _2(\tau _5\tau _3)]\sqrt{gϵ^{}}\mathrm{exp}({\displaystyle \frac{\gamma _2\tau _5}{2}})\xi _2^+(\tau _5)`$
$`{\displaystyle _{\tau _3=0}^{\tau _3=\tau }}d\tau _3\mathrm{exp}[\gamma _3(\tau _3\tau )]g{\displaystyle _{\tau _4=0}^{\tau _4=\tau _3}}d\tau _4\mathrm{exp}[\gamma _1(\tau _4\tau _3)]\sqrt{gϵ^{}}\mathrm{exp}({\displaystyle \frac{\gamma _1\tau _4}{2}})\xi _1^+(\tau _4)`$
$`{\displaystyle _{\tau _5=0}^{\tau _5=\tau _3}}d\tau _5\mathrm{exp}[\gamma _2(\tau _5\tau _3)]\sqrt{gϵ^{}}\mathrm{exp}({\displaystyle \frac{\gamma _2\tau _5}{2}})\xi _2^+(\tau _5).`$
However, the average of their product is zero as
$`\alpha _1^{[1/2]}(\tau )\alpha _2^{[1/2]}(\tau )\mathrm{\Delta }\alpha _3^{+[2]}(\tau )`$ (24)
$`=`$ $`g^3ϵϵ^{}{\displaystyle _{\tau _1=0}^{\tau _1=\tau }}{\displaystyle _{\tau _2=0}^{\tau _2=\tau }}{\displaystyle _{\tau _3=0}^{\tau _3=\tau }}{\displaystyle _{\tau _4=0}^{\tau _4=\tau _3}}{\displaystyle _{\tau _5=0}^{\tau _5=\tau _3}}𝑑\tau _1𝑑\tau _2𝑑\tau _3𝑑\tau _4𝑑\tau _5`$ (27)
$`\mathrm{exp}[\gamma _1(\tau _1\tau )]\mathrm{exp}({\displaystyle \frac{\gamma _1\tau _1}{2}})\mathrm{exp}[\gamma _2(\tau _2\tau )]\mathrm{exp}({\displaystyle \frac{\gamma _2\tau _2}{2}})\mathrm{exp}[\gamma _3(\tau _3\tau )]\mathrm{exp}[\gamma _1(\tau _4\tau _3)]\mathrm{exp}({\displaystyle \frac{\gamma _1\tau _4}{2}})`$
$`\mathrm{exp}[\gamma _5(\tau _5\tau _3)]\mathrm{exp}({\displaystyle \frac{\gamma _2\tau _5}{2}})(\xi _1(\tau _1)\xi _2(\tau _2)\xi _1^+(\tau _4)\xi _2^+(\tau _5)\xi _1(\tau _1)\xi _2(\tau _2)\xi _1^+(\tau _4)\xi _2^+(\tau _5))`$
and
$$\xi _1(\tau _1)\xi _2(\tau _2)\xi _1^+(\tau _4)\xi _2^+(\tau _5)=\xi _1(\tau _1)\xi _2(\tau _2)\xi _1^+(\tau _4)\xi _2^+(\tau _5).$$
(28)
Thus, the two noise terms cancel each other and the right hand side of Eq. (28) is zero. Taking such a consideration into account, the moments of the form $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$ are determined by carefully considering the lowest order nonzero terms of their constituent positive-P variables and then finding the average of these variables’ products.
Consider Fig. (2)(a), which contains the lowest order stochastic diagrams for $`\alpha _i(\tau )`$ and $`\alpha _i^+(\tau )`$, where $`i=1,2`$. The first diagram in it represents the initial value terms $`\alpha _i^{(0)}(\tau )`$ and $`\alpha _i^{(0)+}(\tau )`$, which are zero and do not contribute to any moments. The second represents terms containing $`\xi _1,\xi _1^+,\xi _2`$ or $`\xi _2^+`$, which are not necessarily zero and thus may contribute to moments. Fig. 2(b) contains the lowest order stochastic diagrams for $`\alpha _3(\tau )`$ and $`\alpha _3^+(\tau )`$. In it, all terms represented by stochastic diagrams containing initial value lines are zero except for the $`O(g^0)`$ ones. This is so as these terms contain either $`\alpha _i^{(0)}(\tau )\mathrm{or}\alpha _i^{+(0)}(\tau )`$, where $`i=1,2`$, which are both zero. In addition, all $`O(g^0)`$ terms represented by stochastic diagrams in Fig. 2(b) are canceled out by other $`O(g^0)`$ terms. This occurs because $`\mathrm{\Delta }𝒜_3(\tau )`$ appears in the moments considered. Its two components, $`𝒜_3(\tau )`$ and $`𝒜_3(\tau )`$, contain the same $`O(g^0)`$ term and hence their $`O(g^0)`$ terms cancel each other. It follows that the only remaining stochastic diagram in Fig. 2(b), which represents the $`O(g^2)`$ term containing two noise components, denotes the lowest order term in $`\mathrm{\Delta }𝒜_3(\tau )`$ that is not necessarily zero.
The lowest order nonzero terms determined above are now used to calculate $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$. The lowest order contribution to $`\mathrm{\Delta }\alpha _3(\tau )`$ that is not necessarily zero $`\mathrm{\Delta }\alpha _{3\mathrm{lowest}}(\tau )`$ is
$`\mathrm{\Delta }\alpha _{3\mathrm{lowest}}(\tau )`$ $`=`$ $`{\displaystyle _{\tau _3=0}^{\tau _3=\tau }}𝑑\tau _3\mathrm{exp}[\gamma _3(\tau _3\tau )]g{\displaystyle _{\tau _4=0}^{\tau _4=\tau _3}}𝑑\tau _4\mathrm{exp}[\gamma _1(\tau _4\tau _3)]\sqrt{gϵ}\mathrm{exp}({\displaystyle \frac{\gamma _1\tau _4}{2}})\xi _1(\tau _4)`$ (32)
$`{\displaystyle _{\tau _5=0}^{\tau _5=\tau _3}}𝑑\tau _5\mathrm{exp}[\gamma _2(\tau _5\tau _3)]\sqrt{gϵ}\mathrm{exp}({\displaystyle \frac{\gamma _2\tau _5}{2}})\xi _2(\tau _5)`$
$`{\displaystyle _{\tau _3=0}^{\tau _3=\tau }}d\tau _3\mathrm{exp}[\gamma _3(\tau _3\tau )]g{\displaystyle _{\tau _4=0}^{\tau _4=\tau _3}}d\tau _4\mathrm{exp}[\gamma _1(\tau _4\tau _3)]\sqrt{gϵ}\mathrm{exp}({\displaystyle \frac{\gamma _1\tau _4}{2}})\xi _1(\tau _4)`$
$`{\displaystyle _{\tau _5=0}^{\tau _5=\tau _3}}d\tau _5\mathrm{exp}[\gamma _2(\tau _5\tau _3)]\sqrt{gϵ}\mathrm{exp}({\displaystyle \frac{\gamma _2\tau _5}{2}})\xi _2(\tau _5).`$
When $`\gamma _1=\gamma _2=\gamma _3=\gamma `$, the average of the product of $`\mathrm{\Delta }\alpha _{3\mathrm{lowest}}(\tau )`$ and the lowest order nonzero terms in $`\alpha _1(\tau )`$ and $`\alpha _2(\tau )`$ is approximately equal to $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$ when $`g1`$ and thus
$`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$ $``$ $`{\displaystyle \frac{ϵ^2g^3\mathrm{exp}[3\gamma \tau ]}{\gamma ^2}}\times \left[{\displaystyle \frac{\mathrm{exp}[\gamma \tau ]}{\gamma }}2\tau {\displaystyle \frac{\mathrm{exp}[\gamma \tau ]}{\gamma }}\right].`$ (33)
As daggered positive-P variables are complex conjugate to undaggered ones on average, $`\alpha _1^+(\tau )\alpha _2^+(\tau )\mathrm{\Delta }\alpha _3^+(\tau )=\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )^{}`$. The other six moments of the form $`𝒜_1(\tau )𝒜_2(\tau )\mathrm{\Delta }𝒜_3(\tau )`$ are zero to $`O(g^3)`$. As does $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3^+(\tau )`$, they all have two $`O(g^3)`$ terms that cancel each other. To explain such behaviour in general, the following argument is given. These other six moments can be rewritten as
$$𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau ),$$
(34)
where it is understood that the moments in which $`𝒜_1(\tau )=\alpha _1(\tau ),𝒜_2(\tau )=\alpha _2(\tau ),𝒜_3(\tau )=\alpha _3(\tau )`$ and $`𝒜_1(\tau )=\alpha _1^+(\tau ),𝒜_2(\tau )=\alpha _2^+(\tau ),𝒜_3(\tau )=\alpha _3^+(\tau )`$ are excluded. All $`O(g^3)`$ terms in the six moments of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ under consideration contain noises in one of the following three forms, $`\xi _1(\tau _a)\xi _2(\tau _b)\xi _1^{}(\tau _c)\xi _2^{}(\tau _d)`$, $`\xi _i(\tau _a)\xi _{3i}^{}(\tau _b)\xi _i(\tau _c)\xi _{3i}(\tau _d)`$ and $`\xi _i^{}(\tau _a)\xi _{3i}(\tau _b)\xi _i^{}(\tau _c)\xi _{3i}^{}(\tau _d)`$, where $`i=1,2`$ and time arguments are dummy variables. All $`O(g^3)`$ terms in the six moments of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ under consideration contain the same noises as their corresponding $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ term. However, in these terms of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ four noise averages from corresponding terms of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ are split into the product of two averages of two noises. For example, $`\alpha _1(\tau )\alpha _2(\tau )\alpha _3^+(\tau )`$ contains noises in the form $`\xi _1(\tau _a)\xi _2(\tau _b)\xi _1^+(\tau _c)\xi _2^+(\tau _d)`$ whilst $`\alpha _1(\tau )\alpha _2(\tau )\alpha _3^+(\tau )`$ contains them in the form $`\xi _1(\tau _a)\xi _2(\tau _b)\xi _1^+(\tau _c)\xi _2^+(\tau _d)`$. Using the formula
$`\xi _1(\tau _a)\xi _2(\tau _b)\xi _3(\tau _c)\xi _4(\tau _d)`$ $`=`$ $`\xi _1(\tau _a)\xi _2(\tau _b)\xi _3(\tau _c)\xi _4(\tau _d)`$ (37)
$`+\xi _1(\tau _a)\xi _3(\tau _c)\xi _2(\tau _b)\xi _4(\tau _d)`$
$`+\xi _1(\tau _a)\xi _4(\tau _d)\xi _2(\tau _b)\xi _3(\tau _c)`$
it can be shown that noise expressions in moments of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ under consideration factorize. In particular, they reduce to the noise expression in the six corresponding term of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$. It follows that cancellation occurs between the $`O(g^3)`$ terms in corresponding moments of the form $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ and $`𝒜_1(\tau )𝒜_2(\tau )𝒜_3(\tau )`$ under consideration as the two terms are identical. Consequently, all six moments under consideration are $`O(g^4)`$. They are also typically much smaller than the two $`O(g^3)`$ moments, $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$ and $`\alpha _1^+(\tau )\alpha _2^+(\tau )\mathrm{\Delta }\alpha _3^+(\tau )`$, as $`g1`$ for realistic systems. To be precise, as all moments are complex quantities, the magnitudes of $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$ and $`\alpha _1^+(\tau )\alpha _2^+(\tau )\mathrm{\Delta }\alpha _3^+(\tau )`$ are much larger than the magnitudes of the other six moments.
The above results are now more closely related to experiments by considering quadrature phase amplitudes $`X_{i,\theta _i}(\tau )`$. In particular, calculations are performed to determine the in principle experimentally observable third order quadrature phase amplitude moment $`M(\tau )`$, where $`M(\tau )=\mathrm{\Delta }X_{1,\theta _1}(\tau )\mathrm{\Delta }X_{2,\theta _2}(\tau )\mathrm{\Delta }X_{3,\theta _3}(\tau )`$, according to quantum mechanics and SED. In quantum mechanics, quadrature phase amplitudes are expressed in terms of creation and annihilation operators by the equation
$$\widehat{X}_{i,\theta _i}=\frac{\widehat{a}_i\mathrm{exp}(i\theta _i)+\widehat{a}_i^{}\mathrm{exp}(i\theta _i)}{2}.$$
(38)
Using Eq. (38) and operator-positive-P variable correspondences $`\widehat{M}(\tau )_{QM}`$, the value of $`M(\tau )`$ for quantum mechanics can be expressed as
$$\widehat{M}(\tau )_{QM}=\frac{1}{8}\underset{i=1}{\overset{3}{}}\mathrm{\Delta }\alpha _i(\tau )e^{i\theta _i}+\mathrm{\Delta }\alpha _i^+(\tau )e^{i\theta _i}.$$
(39)
Upon expanding the right hand side of Eq. (39), the two lowest order terms in g, $`\mathrm{\Delta }\alpha _1(\tau )\mathrm{\Delta }\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$ and $`\mathrm{\Delta }\alpha _1^+(\tau )\mathrm{\Delta }\alpha _2^+(\tau )\mathrm{\Delta }\alpha _3^+(\tau )`$, usually dominate. When they do
$$\widehat{M}(\tau )_{QM}\frac{1}{4}[\mathrm{cos}\mathrm{\Theta }\mathrm{Re}\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )\mathrm{sin}\mathrm{\Theta }\mathrm{Im}\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )],$$
(40)
where $`\mathrm{\Theta }=\theta _1+\theta _2+\theta _3`$. However, when $`\mathrm{cos}\mathrm{\Theta }=0`$ and $`\mathrm{Im}(\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau ))=O(\tau ^4)`$ or when $`\mathrm{sin}\mathrm{\Theta }=0`$ and $`\mathrm{Re}(\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau ))=O(\tau ^4)`$ Eq. (40) is not necessarily true. Such situations can be avoided though because $`\mathrm{\Theta }`$ and $`ϵ`$ are controllable parameters. They are ignored in present considerations. When $`ϵ`$ is real, the $`O(g^3)`$ term in $`\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )`$ is also real and so
$$\widehat{M}(\tau )_{QM}\frac{1}{4}\mathrm{cos}\mathrm{\Theta }\alpha _1(\tau )\alpha _2(\tau )\mathrm{\Delta }\alpha _3(\tau )\frac{ϵ^2g^3\mathrm{cos}\mathrm{\Theta }\mathrm{exp}(3\gamma \tau )}{4\gamma ^2}\left[\frac{\mathrm{exp}(\gamma \tau )}{\gamma }2\tau \frac{\mathrm{exp}(\gamma \tau )}{\gamma }\right].$$
(41)
Thus, Eq. (41) shows that $`\widehat{M}(\tau )_{QM}`$ is cubic in g, within the domain considered, as shown in Fig. 3.
## IV comparison of quantum mechanics and stochastic electrodynamics
This section compares the predictions of quantum mechanics and SED for the intracavity moment $`M(\tau )`$. SED is a semiclassical theory which adds Gaussian white noise to classical electrodynamics. It describes electromagnetic field modes by complex field amplitudes $`\beta `$. For the nondegenerate parametric oscillator, the set of such amplitudes $`\{\beta _1,\beta _2,\beta _3\}`$ evolves via the equations
$`{\displaystyle \frac{\beta _1}{\tau }}`$ $`=`$ $`\gamma _1\beta _1+g\beta _2^{}\beta _3+\sqrt{\gamma _1}\xi _1`$ (42)
$`{\displaystyle \frac{\beta _2}{\tau }}`$ $`=`$ $`\gamma _2\beta _2+g\beta _1^{}\beta _3+\sqrt{\gamma _2}\xi _2`$ (43)
$`{\displaystyle \frac{\beta _3}{\tau }}`$ $`=`$ $`\gamma _3\beta _3g\beta _1\beta _2+\sqrt{\gamma _3}\xi _3,`$ (44)
where the same time variable as in the quantum case is used and the $`\xi ^{}s`$ are independent complex Gaussian white noises with the following correlations
$$\xi _i(\tau _1)\xi _j^{}(\tau _2)=\delta _{ij}\delta (\tau _1\tau _2),$$
(45)
where $`i,j=1,2,3`$. The field amplitudes $`\beta _1,\beta _2\mathrm{and}\beta _3`$ initially have Gaussian fluctuations in their real and imaginary parts of variance 1/4. The only nonzero correlations present in these fluctuations are thus
$$\mathrm{\Delta }\beta _i(0)\mathrm{\Delta }\beta _i^{}(0)=\frac{1}{2},$$
(46)
where $`i=1,2,3`$. Initial conditions are $`\beta _1(0)=\beta _2(0)=0\mathrm{and}\beta _3(0)=ϵ`$.
The SED prediction for the intracavity moment $`M(\tau )`$ is $`M(\tau )_{SED}`$, which is given by the equation
$$M(\tau )_{SED}=\mathrm{\Delta }X_{1,\theta _1}(\tau )\mathrm{\Delta }X_{2,\theta _2}(\tau )\mathrm{\Delta }X_{3\theta _3}(\tau ),$$
(47)
where $`X_{i,\theta _i}(\tau )=\left(\beta _i(\tau )e^{i\theta _i}+\beta _i^{}(\tau )e^{i\theta _i}\right)/2`$. It is calculated using a similar iterative method to the one in Section II, except that noise terms are now treated exactly instead of iteratively. Zeroth order approximations for this iterative method are thus
$`\beta _i^{(0)}(\tau )`$ $`=`$ $`\beta _i(0)\mathrm{exp}(\gamma _i\tau )+{\displaystyle _{\tau _{i=0}}^{\tau _i=\tau }}𝑑\tau _i\mathrm{exp}[\gamma _i(\tau _i\tau )]\sqrt{\gamma _i}\xi _i`$ (48)
$`\beta _3^{(0)}(\tau )`$ $`=`$ $`\beta _3(0)\mathrm{exp}(\gamma _3\tau ),`$ (49)
where $`i=1,2`$. Higher order $`(m+1)^{th}`$ order approximations are
$`\beta _i^{(m+1)}(\tau )`$ $`=`$ $`\beta _i^{(0)}(\tau )+{\displaystyle _{\tau _i=0}^{\tau _i=\tau }}𝑑\tau _i\mathrm{exp}[\gamma _i(\tau _i\tau )]g\beta _{3i}^{(m)}(\tau _i)\beta _3^{(m)}(\tau _i)`$ (50)
$`\beta _3^{(m+1)}(\tau )`$ $`=`$ $`\beta _3^{(0)}(\tau ){\displaystyle _{\tau _3=0}^{\tau _3=\tau }}𝑑\tau _3\mathrm{exp}[\gamma _3(\tau _3\tau )]g\beta _1^{(m)}(\tau _3)\beta _2^{(m)}(\tau _3),`$ (51)
where $`i=1,2`$.
The lowest order nonzero term in g of $`M(\tau )_{SED}`$ is now found using the same method as for the lowest order nonzero term of $`\widehat{M}(\tau )_{QM}`$. Consider the moments of the form $`\mathrm{\Delta }_1(\tau )\mathrm{\Delta }_2(\tau )\mathrm{\Delta }_3(\tau )`$, where $`_n(\tau )`$ is either $`\beta _n(\tau )`$ or $`\beta _n^{}(\tau )`$. The stochastic diagrams required to determine the order of the lowest order nonzero terms of these moments are shown in Fig. 4. Note that noise terms are now $`O(g^0)`$, instead of $`O(g^{\frac{1}{2}})`$ as for quantum mechanics. Using the stochastic diagrams in Fig. 4 it is found that, when $`\gamma _1=\gamma _2=\gamma _3=\gamma `$ and $`g,\tau _f1`$,
$$\mathrm{\Delta }\beta _1(\tau )\mathrm{\Delta }\beta _2(\tau )\mathrm{\Delta }\beta _3^{}(\tau )\frac{g}{12\gamma }[1\mathrm{exp}(3\gamma \tau )]\mathrm{\Delta }\beta _1^{}(\tau )\mathrm{\Delta }\beta _2^{}(\tau )\mathrm{\Delta }\beta _3(\tau )^{}.$$
(52)
The other six moments of the form $`\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3`$ are all $`O(g^2)`$. Thus, $`\mathrm{\Delta }\beta _1(\tau )\mathrm{\Delta }\beta _2(\tau )\mathrm{\Delta }\beta _3^{}(\tau )`$ and $`\mathrm{\Delta }\beta _1^{}(\tau )\mathrm{\Delta }\beta _2^{}(\tau )\mathrm{\Delta }\beta _3(\tau )`$ dominate these other six moments when $`g1`$ and hence
$$M(\tau )_{SED}\frac{1}{4}\mathrm{cos}\mathrm{\Phi }\mathrm{\Delta }\beta _1(\tau )\mathrm{\Delta }\beta _2(\tau )\mathrm{\Delta }\beta _3^{}(\tau )\frac{g\mathrm{cos}\mathrm{\Phi }}{4\gamma }[1\mathrm{exp}(1\gamma \tau )],$$
(53)
where $`\mathrm{\Phi }=\theta _1+\theta _2\theta _3`$, when $`\mathrm{\Phi }0`$. Eq. (53) shows that $`M(\tau )_{SED}`$ is linear in g, as shown in Fig. 3. This is in contrast to the cubic behaviour of $`\widehat{M}(\tau )_{QM}`$. Thus, quantum mechanics and SED predict greatly different values for $`M(\tau )`$ when $`g1`$.
Consideration now is given to the effect of damping strength on the size of the difference between $`M(\tau )_{SED}`$ and $`\widehat{M}(\tau )_{QM}`$. Fig. 5 (a) shows $`\widehat{M}_{QM}`$ and $`M_{SED}`$ as functions of $`\gamma `$ for $`g=0.1,\tau =1`$ and $`ϵ=1`$. It indicates that the difference between them is somewhat sensitive to $`\gamma `$, decreasing exponentially with increasing $`\gamma `$ and quickly approaching zero. However, for the shorter time $`\tau =0.1`$, Fig. 5(b) shows that this difference is not as sensitive to damping. It, approximately, only decreases linearly with increasing $`\gamma `$.
SED and the positive-P representation treat fluctuations very differently, as is evident by comparing noise terms in Eqs (3) and (42). This difference in treatment underlies the differences between the two theories’ results. Firstly, noise terms in the positive-P representation are nonlinear and are scaled by either $`\sqrt{g\alpha _3}`$ or $`\sqrt{g\alpha _3^+}`$, whilst those in SED are linear and are scaled by $`\sqrt{\gamma _i}`$. Secondly, noise terms possess different correlations in the two cases. Thirdly, in quantum mechanics no energy fluctuations occur in the vacuum state, whilst in SED $`\{\beta _i\}`$ fluctuates, as does the total energy. Assuming quantum mechanics is true, in SED fluctuations in the vacuum lead to an overestimate of $`M(\tau )`$ for small g.
## V numerical results
The analytic results for $`M(\tau )_{SED}`$ and $`\widehat{M}(\tau )_{QM}`$ in Sections III and IV only include lowest order nonzero terms. This leaves the sums of all higher order terms as neglected and these may be significant. For this reason, the validity of the analytic approximations are checked by comparison with highly accurate numerical simulation results.
Numerical simulation methods for stochastic differential equations (SDE’s) are both somewhat complex and not widely known. Thus, explanations are given for the numerical technique used to solve the SDE’s in Eqs (3) and (42). Normal ODE techniques such as the Runge-Kutta method cannot be used to solve SDE’s as they contain discontinuous source terms. Instead, a semi-implicit numerical method is employed. Only its application to Eq. (3) is explained as its application Eq. (42) is similar. Each of the equations in Eq. (3) can be rewritten as
$$\frac{x_i}{\tau }=A_i(𝐱)+\underset{j}{}B_{ij}(𝐱)\zeta _j(\tau ),$$
(54)
where $`x_i`$ is either $`\alpha _i`$ or $`\alpha _i^+`$, for $`i=1,2,3`$, $`𝐱`$ is a vector whose components are $`\{\alpha _i,\alpha _i^+\}`$, $`A_i`$ is the function of $`𝐱`$ formed by the damping and nonlinear terms in the evolution equation for $`x_i`$ and $`b_{ij}`$ is a matrix whose elements are coefficients of the noise terms $`\{\zeta _j\}`$ where $`\zeta _j`$ is either $`\xi _j`$ or $`\xi _j^{}`$, for $`j=1,2`$. The semi-implicit method used determines $`\overline{𝐱}^{(n)}`$, an approximation to $`𝐱`$ at the midpoint of the interval $`(\tau _n,\tau _{n+1})`$. This approximation is found using iteration such that the $`p^{th}`$ order approximation to a component of $`\overline{𝐱}^{(n)}`$ $`\overline{x_i}^{(n)[p]}`$ is given by the equation
$$\overline{x}_i^{(n)[p]}=x_i^{(n)}+\frac{1}{2}\left[\mathrm{\Delta }\tau A_i(\overline{𝐱}^{(n)[p1]})+\underset{j}{}B_{ij}(\overline{𝐱}^{(n)[p1]})\mathrm{\Delta }W_j(\overline{\tau }_n)\right],$$
(55)
where $`x_i^{(n)}`$ is the value of $`x_i`$ at time $`\tau _n`$ , $`\mathrm{\Delta }\tau =\tau _{n1}\tau _n`$, $`\mathrm{\Delta }W_j(\overline{\tau }_n)=\zeta _j^{(n)}(\overline{\tau }_n)\mathrm{\Delta }\tau `$ and $`\overline{\tau }_n`$ is the midpoint of the interval $`(\tau _{n1},\tau _n)`$. The zeroth order approximation to $`\overline{x}_i^{(n)}`$ is given by the equation $`\overline{x}_i^{(n)[0]}=x_i^{(n)}`$. The approximation to $`\overline{𝐱}^{(n)}`$ calculated is then used to generate $`\mathrm{\Delta }x_i^{(n)}`$, an approximation to the change in $`x_i`$ over the interval $`(\tau _n,\tau _{n+1})`$. This is done by solving the equation
$$\mathrm{\Delta }x_i^{(n)}=A_i(\overline{𝐱}^{(n)})\mathrm{\Delta }\tau _n+\underset{j}{}B_{ij}(\overline{𝐱}^{(n)})\mathrm{\Delta }W_j(\overline{\tau }_n).$$
(56)
Repeated use of Eq. (56) determines $`x_i^{(n)}`$ for successively later and later times and thus solves Eq. (54). Two of the most important parameters used in the numerical simulations are the step size and the number of stochastic paths that are averaged over. The former is always 0.0025 and the latter is $`O(10^6)`$ for most simulations. However, large sampling errors necessitated averaging over $`O(10^7)`$ paths for $`g=0.1`$ SED simulations.
Results from the numerical and analytic simulations of $`M(\tau )_{SED}`$ and $`\widehat{M}(\tau )_{QM}`$ over a range of $`g`$ and $`N`$ values, where $`N`$ is the average initial number of pump photons ($`N=|ϵ^2|`$), are shown in Figs 6-8. In all cases $`\theta _1=\theta _2=\theta _3=0`$ and relative numerical errors are small. All $`g=0.1`$ analytic results are in agreement with their numerical counterparts. However, $`g=1`$ analytic results for N=1 and N=10 are not. This disagreement is explained by noting that the analytic results are only necessarily valid when $`g1`$.
A number of qualitative trends can be seen in Figs 6-8. In Figs 6 and 7 (N=1 and N=10) the results of SED and quantum mechanics are so distinct that they have different signs, with those of quantum mechanics being negative and those of SED being positive. This trend only holds for short times ($`\tau <0.07`$) in Figs 8(a) and (b) (N=100). For longer times, SED and quantum mechanics predict the same sign. This trait is consistent with the fact that Figs 8(a) and (b) show results for the largest number of photons in the pump mode. SED and quantum mechanics are at their most classical level for this case and thus might be expected to differ the least. For constant N Figs 6-8 also show that as $`g`$ is decreased the results of quantum mechanics and SED become more similar. This occurs because lower $`g`$ values are associated with larger damping to nonlinear coupling ratios and therefore move SED and quantum mechanics closer to the classical domain.
## VI External moments
Thus far, only intracavity fields have been considered. However, it is the external fields that leak out of a cavity that are observed. In realistic systems, intracavity photons are transmitted through imperfect mirrors into the external environment where they are detected. Thus, an external field analogue of $`M(\tau )`$, $`M^{(E)}(\tau _s,\tau _f)`$, where $`\tau _s`$ and $`\tau _f`$ are initial and final measurement times, is calculated according to quantum mechanics and SED to consider what is actually observed in the laboratory.
The first step in calculating the external moment $`\widehat{M}^{(E)}(\tau )_{QM}`$ for quantum mechanics is defining the external quadrature phase amplitudes constituting it. This is done within the the context of homodyne detection as quadrature phase amplitudes are commonly measured using it. A schematic diagram for balanced homodyne detection is shown in Fig. 9. An external signal field flux $`\widehat{\mathrm{\Phi }}_{iOUT}`$, where $`i=1,2,3`$, and a local oscillator field flux $`E_i`$ are incident on a 50-50 beam splitter BS. An external local oscillator phase variable is represented by $`\overline{\theta }_i`$. The two field fluxes combine and are detected by two photodiodes $`D_{+i}`$ and $`D_i`$. The detected photocurrents are then converted to amplified electrical currents whose difference is found. An external quadrature phase amplitude for quantum mechanics $`\widehat{X}_{i,\overline{\theta }_i}^{(E)}`$ is defined as this difference yielding, when $`E_i`$ is real,
$$\widehat{X}_{i,\overline{\theta }_i}^{(E)}(\tau )=\frac{eA\eta _iE_i(\widehat{\mathrm{\Phi }}_{iOUT}(\tau )e^{i\overline{\theta }_i}+\widehat{\mathrm{\Phi }}_{iOUT}^{}{}_{}{}^{}(\tau )e^{i\overline{\theta }_i})}{2},$$
(57)
where e is the magnitude of the charge of an electron, A is an amplification factor and $`\eta _i`$ is a detector efficiency factor for both detectors associated with external field modes denoted by $`i`$. In realistic experiments detection occurs over a finite period of time and thus
$$_{\tau =\tau _s}^{\tau =\tau _f}\frac{d\tau }{\mathrm{\Gamma }}\widehat{X}_{i,\overline{\theta }_i}^{(E)}(\tau )$$
(58)
corresponds to what is observed. Only the $`\overline{\theta }_i=0`$ case is considered. Thus, an external moment analogue of the intracavity moment $`\widehat{M}(\tau )_{QM}`$ can be defined as
$`\widehat{M}^{(E)}(\tau _s,\tau _f)_{QM}`$ $`=`$ $`\mathrm{\Gamma }^3{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle _{\tau _i=\tau _s}^{\tau _i=\tau _f}}𝑑\tau _i\mathrm{\Delta }\widehat{X}_{i,\overline{\theta }_i}^{(E)}(\tau _i)`$ (59)
$`=`$ $`\mathrm{\Gamma }^3{\displaystyle _{\tau _1=\tau _s}^{\tau _1=\tau _f}}{\displaystyle _{\tau _2=\tau _s}^{\tau _2=\tau _f}}{\displaystyle _{\tau _3=\tau _s}^{\tau _3=\tau _f}}{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{\Delta }\widehat{X}_{i,\overline{\theta }_i}^{(E)}(\tau _i).`$ (60)
To calculate $`\widehat{M}^{(E)}(\tau _s,\tau _f)_{QM}`$, the relation between the unknown external output fields that define it and known intracavity fields needs to be ascertained. Gardiner and Collett have formulated an input-output theory which relates the two via the equation
$$\widehat{\mathrm{\Phi }}_{iOUT}(\tau )=\sqrt{2\mathrm{\Gamma }}\widehat{a}_i(\tau )+\widehat{\mathrm{\Phi }}_{iIN}(\tau ),$$
(61)
where $`\widehat{\mathrm{\Phi }}_{iIN}(\tau )`$ is the input field flux associated with intracavity mode $`i`$. All input fields are assumed to be in vacuum states. This allows the use of Eq. (5.3) from , which can be expressed as, in this paper’s notation,
$`\widehat{\mathrm{\Phi }}_{iOUT}^{}(\tau _1)\widehat{\mathrm{\Phi }}_{iOUT}^{}(\tau _2)\mathrm{}\widehat{\mathrm{\Phi }}_{iOUT}^{}(\tau _n)\widehat{\mathrm{\Phi }}_{iOUT}(\tau _{n+1}^{})\mathrm{}\widehat{\mathrm{\Phi }}_{iOUT}(\tau _m^{})`$ (62)
$`=(2\mathrm{\Gamma })^{m/2}\stackrel{~}{T}[\widehat{a}_i^{}(\tau _1)\widehat{a}_i^{}(\tau _2)\mathrm{}\widehat{a}_i^{}(\tau _n)]T[\widehat{a}_i(\tau _{n+1}^{})\mathrm{}\widehat{a}_i(\tau _m^{})]`$ , (63)
where $`\stackrel{~}{T}`$ and $`T`$ are time anti-ordering and time ordering operators respectively. Using Eq. (57), the integrand of Eq. (59) can be expressed in terms of $`\widehat{\mathrm{\Phi }}_{iOUT}`$ and $`\widehat{\mathrm{\Phi }}_{iOUT}^{}`$. It can then be expressed in terms of particular $`\widehat{a}_i(\tau _i)`$ and $`\widehat{a}_{i}^{}{}_{}{}^{}(\tau _i)`$ averages using Eq. (62). In turn, these averages are equivalent to the eight positive-P averages of the form $`\mathrm{\Delta }𝒜_1(\tau _1)\mathrm{\Delta }𝒜_2(\tau _2)\mathrm{\Delta }𝒜_3(\tau _3)`$, where $`𝒜_i`$ is either $`\alpha _i`$ or $`\alpha _i^+`$. As was determined in Section III, two of these averages, $`\mathrm{\Delta }\alpha _1(\tau _1)\mathrm{\Delta }\alpha _2(\tau _2)\mathrm{\Delta }\alpha _3(\tau _3)`$ and $`\mathrm{\Delta }\alpha _1^+(\tau _1)\mathrm{\Delta }\alpha _2^+(\tau _2)\mathrm{\Delta }\alpha _3^+(\tau _3)`$, are of lower order in $`g`$ than the others and hence dominate when $`g1`$. Thus, the external field moment of quantum mechanics $`\widehat{M}^{(E)}(\tau _s,\tau _f)_{QM}`$ can be expressed as, when $`\overline{\theta }_i=0`$, where $`i=1,2,3`$,
$`\widehat{M}^{(E)}(\tau _s,\tau _f)_{QM}`$ $``$ $`{\displaystyle \frac{\sqrt{2}(eA\eta E)^3\mathrm{\Gamma }^{3/2}}{4}}{\displaystyle _{\tau _1=\tau _s}^{\tau _1=\tau _f}}{\displaystyle _{\tau _2=\tau _s}^{\tau _2=\tau _f}}{\displaystyle _{\tau _3=\tau _s}^{\tau _3=\tau _f}}𝑑\tau _1𝑑\tau _2𝑑\tau _3`$ (65)
$`\mathrm{\Delta }\alpha _1(\tau _1)\mathrm{\Delta }\alpha _2(\tau _2)\mathrm{\Delta }\alpha _3(\tau _3)+\mathrm{\Delta }\alpha _1^+(\tau _1)\mathrm{\Delta }\alpha _2^+(\tau _2)\mathrm{\Delta }\alpha _3^+(\tau _3),`$
where $`\eta =\eta _i`$ and $`E=E_i`$, where $`i=1,2,3`$. To simplify the algebra only the $`\tau _s=0`$ case is investigated, so that only the moment $`\widehat{M}^{(E)}(\tau _f)_{QM}(\widehat{M}^{(E)}(0,\tau _f)_{QM})`$ is considered. External fields are only considered for small times ($`\tau _f<<1`$), and so, to a given order in $`g`$, $`\widehat{M}^{(E)}(\tau _f)_{QM}`$’s lowest nonzero order term in $`\tau _f`$ dominates. Hence $`\widehat{M}^{(E)}(\tau _f)_{QM}`$ can be approximated by its lowest nonzero order term in both $`g`$ and $`\tau _f`$. Thus,
$$\widehat{M}^{(E)}(\tau _f)_{QM}\frac{\sqrt{2}}{48}g^3ϵ^2(eA\eta E)^3\mathrm{\Gamma }^{3/2}\tau _f^6.$$
(66)
The SED external moment $`M^{(E)}(\tau _f)_{SED}`$ is now calculated. It is given by the same expression as $`\widehat{M}(\tau _f)_{QM}`$, the right hand side of Eq. (59) (when $`\tau _s=0`$), except that the quadrature phase amplitude operator $`\widehat{X}_{i\overline{\theta }_i}^{(E)}(\tau )`$, is replaced by its SED c-number analogue. This external SED c-number quadrature phase amplitude is defined as, when $`E_i`$ is real,
$$X_{i,\overline{\theta }_i}^{(E)}(\tau )=\frac{eA\eta _iE_i(\beta _{iOUT}(\tau )e^{i\overline{\theta }_i}+\beta _{iOUT}^{}(\tau )e^{i\overline{\theta }_i})}{2},$$
(67)
where $`\beta _{iOUT}(\tau )`$ is the output field flux associated with the intracavity field denoted by $`i`$. In analogy with Eq. (61), it is assumed that the SED input-output relation is
$$\beta _{iOUT}(\tau )=\sqrt{2\mathrm{\Gamma }}\beta _i(\tau )+\beta _{iIN}(\tau ),$$
(68)
where $`\beta _{iIN}(\tau )`$ is the input field flux for the intracavity mode $`i`$. When all input fields are in vacuum states, as is the case, $`\beta _{iIN}(\tau )`$ is a Gaussian white noise with a self correlation characterized by
$$\beta _{iIN}(\tau _i)\beta _{iIN}^{}(\tau _i^{})=\frac{\delta (\tau _i\tau _i^{})}{2\mathrm{\Gamma }}.$$
(69)
A calculation analogous to the quantum mechanical one earlier in this section can be performed using Eqs (67) and (68) to obtain an expression for $`M^{(E)}(\tau _f)_{SED}`$ in terms of particular intracavity averages. When lowest order nonzero approximations to these averages are considered, the following result is obtained when $`\overline{\theta }_i`$ for $`i=1,2,3`$, and $`g,\tau _f1`$,
$$M^{(E)}(\tau _f)_{SED}\frac{\sqrt{2}}{16}g\tau _f^4(eA\eta E)^3\mathrm{\Gamma }^{3/2}.$$
(70)
Upon comparing Eq. (70) to result of quantum mechanics in Eq. (66), it is seen that the leading order term in $`g`$ in Eq. (66) is $`O(g^3)`$ whilst in Eq. (70) it is $`O(g)`$. Hence, as was the case for the intracavity moment, quantum mechanics and SED predict significantly different results for the observable external field moment $`M^{(E)}(\tau _f)`$.
## VII Signal to noise ratio
In actual experiments, only finite samples of results are obtained, as opposed to infinite ones. Hence, in practice the population means considered thus far are estimated from sample means. These sample means fluctuate from sample to sample and thus have signal to noise ratios, which are now determined for small times ($`\tau _f<<1`$). This paper focuses on differences between quantum mechanics and SED. Thus, a calculation is performed of the signal to noise ratio of the difference between the two theories’ external sample moments. First, the noise of the external sample moment in quantum mechanics is determined. It is then assumed that the noise of the external sample moment of SED is the same. Noise results are combined with the external moment results of Section VI to produce $`S(\tau _f)`$, the signal to noise ratio of the difference between the two theories’ external sample moments. This quantity $`S(\tau _f)`$ is given by
$`S(\tau _f)`$ $`=`$ $`{\displaystyle \frac{|m^{(E)}(\tau _f)_{SED}\widehat{m}^{(E)}(\tau _f)_{QM}|}{\sqrt{s^2(m^{(E)}(\tau _f)_{SED})+s^2(\widehat{m}^{(E)}(\tau _f)_{QM})}}}`$ (71)
$`=`$ $`{\displaystyle \frac{\sqrt{n1}|M^{(E)}(\tau _f)_{SED}\widehat{M}^{(E)}(\tau _f)_{QM}|}{\sqrt{2}\sigma (\widehat{M}^{(E)}(\tau _f))}},`$ (72)
where $`n`$ is the number of observations in the sample considered, $`m^{(E)}(\tau _f)_{SED}`$ and $`\widehat{m}^{(E)}(\tau _f)_{QM}`$ are sample averages of $`M^{(E)}(\tau _f)`$ according to SED and quantum mechanics respectively, and $`s^2(A)`$ denotes the sample variance of A. The only significant unknown quantity on the right hand side of Eq. (71) is $`\sigma (\widehat{M}^{(E)}(\tau _f))`$, which is now determined. Expressing $`\sigma (\widehat{M}^{(E)})`$ explicitly yields
$$\sigma (\widehat{M}^{(E)}(\tau _f))=\sqrt{\widehat{M}^{(E)}(\tau _f)^2_{QM}\widehat{M}^{(E)}(\tau _f)_{QM}^2}.$$
(73)
The moment $`\widehat{M}^{(E)}(\tau _f)_{QM}`$ was determined in Section VI and so $`\widehat{M}^{(E)}(\tau _f)^2_{QM}`$ is now calculated. In the calculation that follows only the $`\overline{\theta }_i=0`$, where $`i=1,2,3`$, and $`g,\tau _f1`$ cases are considered.
The moment $`\widehat{M}^{(E)}(\tau _f)^2_{QM}`$ can be expressed in terms of external quadrature phase operators as
$`\widehat{M}^{(E)}(\tau _f)^2_{QM}`$ $`=`$ $`[\mathrm{\Gamma }^3{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle _{\tau _i=0}^{\tau _i=\tau _f}}𝑑\tau _i\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)]^2`$ (74)
$`=`$ $`\mathrm{\Gamma }^6{\displaystyle _{\tau _1=0}^{\tau _1=\tau _f}}{\displaystyle _{\tau _1^{}=0}^{\tau _1^{}=\tau _f}}{\displaystyle _{\tau _2=0}^{\tau _2=\tau _f}}{\displaystyle _{\tau _2^{}=0}^{\tau _2^{}=\tau _f}}{\displaystyle _{\tau _3=0}^{\tau _3=\tau _f}}{\displaystyle _{\tau _3^{}=0}^{\tau _3^{}=\tau _f}}𝑑\tau _1𝑑\tau _1^{}𝑑\tau _2𝑑\tau _2^{}𝑑\tau _3𝑑\tau _3^{}{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{}),`$ (75)
where $`\widehat{X}_i^{(E)}(\tau _i)=\widehat{X}_{i\overline{\theta }_i=0}^{(E)}(\tau _i).`$ The integrand of Eq. (74), which is denoted by K, can be expressed as
$$K=\underset{i=1}{\overset{3}{}}\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})+f(\tau _i,\tau _i^{}),$$
(76)
where $`f(\tau _i,\tau _i^{})`$ is a function which includes terms resulting from coupling between modes. These coupling terms vanish when $`g=0`$ and thus are at least O(g). It follows that $`K`$ can be be re-expressed as
$$K=\underset{i=1}{\overset{3}{}}\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})+O(g).$$
(77)
The moment $`\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})`$ is now calculated using a normally ordered approach that has been previously employed to solve similar problems ,. This method expresses $`\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})`$ in terms of normally ordered photocurrent averages and then determines these averages. It first defines $`\widehat{X}_i^{(E)}(\tau _a)`$, where $`\tau _a`$ is any $`\tau `$ variable, as the difference between the amplified electrical currents, $`\widehat{X}_{+i}^{(E)}(\tau _a)`$ and $`\widehat{X}_i^{(E)}(\tau _a)`$, produced by the photocurrents detected at the detectors $`D_{+i}`$ and $`D_i`$ in Fig. 9 in Section VI. Using this definition ($`\widehat{X}_i^{(E)}(\tau _a)=\widehat{X}_{+i}^{(E)}(\tau _a)\widehat{X}_i^{(E)}(\tau _a)`$), $`\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})`$ can be expressed as
$$\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})=(\mathrm{\Delta }\widehat{X}_{+i}^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i))(\mathrm{\Delta }\widehat{X}_{+i}^{(E)}(\tau _i^{})\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})).$$
(78)
Upon expansion, the right hand side of Eq. (78) contains two types of terms, those of the form $`\widehat{X}_{Ci}^{(E)}(\tau _a)`$, where C is either + or -, and those of the form $`\widehat{X}_{Ci}^{(E)}(\tau _i)\widehat{X}_{Di}^{(E)}(\tau _i^{})`$, where D is either + or -. Terms of the form $`\widehat{X}_{Ci}^{(E)}(\tau _a)`$ are given by the equation
$$\widehat{X}_{Ci}^{(E)}(\tau _a)=_{\mathrm{}}^+\mathrm{}\frac{ds_1}{\mathrm{\Gamma }}G_{Ci}^{(1)}(s_1)J^0(\tau _as_1),$$
(79)
where $`G_{Ci}^{(1)}(s_1)`$ is a first order Glauber correlation function and $`J^0(\tau _as_1)`$ is an electrical current pulse produced by a single photodetection event. In following previous work , square electrical current pulses of the form
$`J^0\left(ab\right)=\{\begin{array}{cc}Ae\mathrm{\Gamma }/\tau _d\hfill & bab+\tau _d\hfill \\ 0\hfill & a<b\mathrm{and}a>b+\tau _d\hfill \end{array}`$ (82)
are considered in the limit of $`\tau _d0`$, which is taken at some appropriate later stage of the calculation. The Glauber correlation function $`G_{Ci}^{(1)}(s_1)`$ can be expressed as a power series in $`g`$ and $`s_1`$ and thus as $`_{m,n=0}^\mathrm{},\mathrm{}c_{mn}g^ms_1^n`$. Due to the form of $`J^0(\tau _as_1)`$, when $`\tau _a<<1`$, as is being assumed, only photodetection events at small times $`s_1`$, contribute to $`\widehat{X}_{Ci}^{(E)}(\tau _a)`$. This fact, coupled with the knowledge that only the $`g<<1`$ case is considered, means that the $`n=m=0`$ term in the power series for $`G_{Ci}^{(1)}(s_1)`$ dominates when $`J^0(\tau _as_1)`$ is nonzero. Hence, upon calculating this dominant term by expressing $`\widehat{\mathrm{\Phi }}_i`$ and $`\widehat{\mathrm{\Phi }}_i^{}`$ in terms of intracavity field operators, in the limit of large local oscillator amplitude,
$$G_{Ci}^{(1)}(s_1)\frac{\eta _{Ci}}{2}E_i^2,$$
(83)
where $`\eta _{Ci}`$ is a detector efficiency factor for the photodetector $`D_{Ci}`$. It follows that
$$\widehat{X}_{Ci}^{(E)}(\tau _a)\frac{\eta _{Ci}E_i^2Ae}{2}.$$
(84)
Terms of the form $`\widehat{X}_{Ci}^{(E)}(\tau _i)\widehat{X}_{Di}^{(E)}(\tau _i^{})`$ in Eq. (78) can be expressed as
$`\widehat{X}_{Ci}^{(E)}(\tau _i)\widehat{X}_{Di}^{(E)}(\tau _i^{})`$ $`=`$ $`\delta _{CD}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{ds_1}{\mathrm{\Gamma }}}G_{Ci}^{(1)}(s_1)J^{(0)}(\tau _is_1)J^{(0)}(\tau _i^{}s_1)`$ (86)
$`+{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{ds_1ds_2}{\mathrm{\Gamma }^2}}G_{C,Di}^{(2)}(s_1,s_2)J^{(0)}(\tau _is_1)J^{(0)}(\tau _i^{}s_2),`$
where $`G_{C,Di}^{(2)}(s_1,s_2)`$ is a second order Glauber correlation function and $`\delta _{C,D}`$ is one when C and D are the same and zero otherwise. In the limit of $`\tau _d0`$,
$$_{\mathrm{}}^+\mathrm{}\frac{ds_1}{\mathrm{\Gamma }}G_{Ci}^{(1)}(s_1)J^{(0)}(\tau _is_1)J^{(0)}(\tau _i^{}s_1)\frac{(eAE_i)^2\eta _{Ci}\delta (\tau _i\tau _i^{})\mathrm{\Gamma }}{2}$$
(87)
to leading nonzero order in $`g,\tau _i`$ and $`\tau _i^{}`$. It is of equal order in $`g`$ and lower order in $`\tau _i`$ and $`\tau _i^{}`$ than the second term in Eq. (86) and hence is much larger than this second term when it is nonzero as the $`\tau _i,\tau _i^{}<<1`$ case is being considered. Thus
$$\widehat{X}_{Ci}^{(E)}(\tau _i)\widehat{X}_{Di}^{(E)}(\tau _i^{})\frac{\delta _{CD}(AeE_i)^2\eta _{Ci}\delta (\tau _i\tau _i^{})\mathrm{\Gamma }}{2}.$$
(88)
From Eqs (84) and (88) it can be seen that the single integral terms in $`\widehat{X}_{+i}^{(E)}(\tau _i)\widehat{X}_{+i}^{(E)}(\tau _i^{})`$ and $`\widehat{X}_i^{(E)}(\tau _i)\widehat{X}_i^{(E)}(\tau _i^{})`$ are of the same order in $`g`$ and lower order in $`\tau _i`$ and $`\tau _i^{}`$ than any other terms contributing to $`\widehat{X}_i^{(E)}(\tau _i)\widehat{X}_i^{(E)}(\tau _i^{})`$ and hence dominate. It follows that
$$\widehat{X}_i^{(E)}(\tau _i)\widehat{X}_i^{(E)}(\tau _i^{})(AeE)^2\eta _i\delta (\tau _i\tau _i^{})\mathrm{\Gamma },$$
(89)
where $`\eta _i=\eta _{Ci}=\eta _{Di}`$ and $`E=E_i`$, where $`i=1,2,3`$. As right hand side of Eq. (89) is $`O(g^0)`$, $`_{i=1}^3\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i)\mathrm{\Delta }\widehat{X}_i^{(E)}(\tau _i^{})`$ is also $`O(g^0)`$ and hence from Eq. (77),
$$K((AeE)^2\mathrm{\Gamma })^3\underset{i=1}{\overset{3}{}}\eta _i\delta (\tau _i\tau _i^{}).$$
(90)
Substituting this approximation for K into Eq. (74) yields
$`\widehat{M}^{(E)}(\tau _f)^2`$ $``$ $`\mathrm{\Gamma }^6((eAE)^2\eta \mathrm{\Gamma })^3{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle _{\tau _i=0}^{\tau _i=\tau _f}}{\displaystyle _{\tau _i^{}}^{\tau _i^{}=\tau _f}}𝑑\tau _i𝑑\tau _i^{}\delta (\tau _i\tau _i^{})`$ (91)
$`=`$ $`\left[{\displaystyle \frac{(eAE)^2\eta \tau _f}{\mathrm{\Gamma }}}\right]^3,`$ (92)
where $`\eta =\eta _i`$, where $`i=1,2,3`$. Thus
$$\sigma (\widehat{M}^{(E)}(\tau _f))\left[\frac{(eAE)^2\eta \tau _f}{\mathrm{\Gamma }}\right]^{3/2}.$$
(93)
Hence, the signal to noise ratio of the difference between the external sample moments of quantum mechanics and SED is
$$S(\tau _f)\frac{\sqrt{n1}\eta ^{3/2}g\tau _f^{5/2}}{16}.$$
(94)
## VIII Realistic systems
Realistic parameter values are now considered to determine if the theoretical difference between SED and quantum mechanics could be observed experimentally. In particular, the signal to noise ratio of the difference between the sample moments of quantum mechanics and SED $`S`$ is calculated using realistic parameter values for nondegenerate parametric oscillators containing the commonly used crystals, silver gallium selinide ($`\mathrm{AgGaSe}_2`$) and potassium titanyl phosphate (KTP). The non-linear interaction strength G for parametric down conversion is given by
$$Gd_{eff}\sqrt{\frac{2\mathrm{}\omega _1\omega _2\omega _3}{ϵ_0V}}\frac{l}{L},$$
(95)
where V is the cavity volume, l is the crystal length and L the cavity length. Cavity and crystal length values of 10cm are chosen. The cavity volume V is given by the formula $`V=\pi \mathrm{\Omega }^2L`$, where $`\mathrm{\Omega }`$ is the spot size. This volume is minimized in order to maximize G and thus the external difference between quantum mechanics and SED. It is assumed that the damping constant used to scale time $`\mathrm{\Gamma }`$ equals the unscaled damping constant for each mode $`\mathrm{\Gamma }_i`$ ($`\mathrm{\Gamma }=\mathrm{\Gamma }_i`$). This common damping constant $`\mathrm{\Gamma }`$ is calculated from the formula $`\mathrm{\Gamma }=T\times \frac{c}{2L}`$, where c is the speed of light and T is a mirror transmission coefficient. A $`T`$ value of $`T=0.01`$ is used. Using the above information, Table I shows realistic parameter values for $`d_{eff}`$, $`V`$, $`G`$, pump, signal and idler wavelengths, and resulting $`g`$ and $`\mathrm{\Gamma }`$ values. Results for $`\widehat{M}^{(E)}_{QM}`$ and $`M^{(E)}_{SED}`$ are obtained using Eqs (66) and (70) for when $`\eta =1`$, $`\tau _f=0.1`$, $`E=10^9s^{1/2}`$, $`A=1/e`$ and $`ϵ=10^3`$. These are displayed in Table II, which shows that the external results of quantum mechanics and SED differ greatly. Due to local oscillator amplification, they are also macroscopically distinct with respect to photon number, even though the initial number of intracavity photons is small on average. Another appealing feature of the difference between the two theories is that detector efficiencies approach one as photodiodes as opposed to photomultipliers are used for detection. Thus, no fair sampling assumptions need to be made.
The question remains of whether or not the population difference between SED and quantum mechanics could be reliably observed in a finite sample of results. To answer it, $`S`$ is now considered. Figs 10 and 11 show graphs of $`S`$ versus sample size $`n`$ for $`\mathrm{AgGaSe}_2`$ and KTP for the same parameter values as used in the last paragraph. These show reasonable $`S`$ values and indicate that large sample sizes must be obtained to produce a signal to noise ratio of one, the smallest signal to noise ratio required to clearly observe the signal. In particular, sample sizes of $`1.8\times 10^{13}`$ (KTP) and $`3.7\times 10^{11}`$ ($`\mathrm{AgGaSe}_2`$) need to be obtained to generate a signal to noise ratio of one. An individual observation takes a time of the order $`t=\tau _f/\mathrm{\Gamma }=6.7\times 10^9s`$ and so, assuming minimal time delay between measurements, $`1.8\times 10^{13}`$ observations would take about 33 hours and $`3.7\times 10^{11}`$ observations about 41 minutes. It is conceivable that measurements could be taken over both times. Furthermore, as the signal to noise ratio scales as $`\frac{1}{g}`$, higher $`g`$ materials would enable the difference to be observed even more readily.
## IX discussion
It has been shown that there exist a significant, potentially experimentally observable, difference between quantum mechanics and SED. Due to local oscillator amplification, this difference can involve macroscopically distinct external fields for the two theories. Thus, it can be considered macroscopic if it is legitimate to include the local oscillators as part of the system and not as external measuring apparatuses. The difference is also potentially experimentally observable, as a realistic system and state are considered and is present at realistic parameter values. The system is practical as parametric oscillators and balanced homodyne detection are widely used, and damping is included. The state is realistic as the initial intracavity coherent state can be approximated well by a laser. It follows that the difference can be seen as providing the basis for an experimentally achievable macroscopic test of quantum mechanics against one local hidden variable theory (SED). Such a test is significant as all experimental tests of quantum mechanics against local hidden variable theories to date have been microscopic. It is true that many macroscopic tests have been proposed, but most of them consider highly idealized states or systems that are not currently able to be experimentally implemented. In particular, many of them do not consider damping, even though it is known to rapidly destroy the correlations of quantum mechanics present in Schroedinger cat and other entangled states. The calculations in this paper do include damping and show that the difference between SED and quantum mechanics is not overly sensitive to it. Most importantly, it remains for realistic damping values. The test proposed in this paper can be seen as being in the novel and largely unexplored domain of macroscopic experimental tests of quantum mechanics.
Even if the local oscillators are not included as part of the system investigated, the external difference between quantum mechanics and SED is still at least mesoscopic as average initial pump photon numbers up to $`10^6`$ are considered. From this perspective, the difference is still distinct from many earlier microscopic ones known to exist between quantum mechanics and all local hidden variable theories. It is also, perhaps, more surprising than some of them as it occurs in a larger particle number system.
Two noteworthy features of the external difference between quantum mechanics and SED are that it involves continuous variables and high efficiency detection. That it involves continuous variables is significant because most previous differences between quantum mechanics and local hidden variable theories have involved discrete ones. Furthermore, it is, perhaps, more surprising that a difference between quantum mechanics and a local hidden variable theory can be found for continuous variables as continuous variables are more closely related to classical ones (which are all continuous) than discrete ones. Low detector efficiency forms the basis of a significant loophole in most tests between quantum mechanics and SED to date . The use of photodiodes for detection in the scheme discussed means that such a loophole is avoided.
The calculations in Section VIII show it is difficult to observe the external difference between quantum mechanics and SED. This is mainly a result of small experimental nonlinearities. They cause few signal and idler photons to be created and thus the experimental signal is weak relative to its noise. For small enough measurement samples, SED results cannot be clearly distinguished from those of quantum mechanics. This fact is consistent with the knowledge that SED reproduces many features of quantum mechanics. However, it is a distinct theory and does differ from quantum mechanics in particular cases, as this paper has shown.
The external difference between quantum mechanics and SED would be easier to observe if larger nonlinear coupling constants were used. These could be achieved by using organic nonlinear crystals such as N-(4-nitrophenyl)-L-prolinol (NPP) . However, phase matching would be difficult with such crystals. In addition, they are typically only transparent within a small frequency range. Alternatively, higher nonlinearities could be achieved by using Josephson-parametric amplifiers , which can have even larger nonlinearities than organic nonlinear crystals. Another possibility, in the area of atom optics, is to to utilize BEC nonlinear effects, in which atom-molecule coupling is induced through photon-associaton .
To conclude, this paper compared particular moments of quantum mechanics to those of SED for the nondegenerate parametric oscillator. Both internal and external moments were considered and an analytic iterative technique showed them both to be cubic in the system’s nonlinear coupling constant for quantum mechanics and linear for SED. Numerical simulations were performed to check the approximate intracavity analytic result and were in agreement with them when the system’s nonlinear coupling constant was much less than one. Realistic parameter values were considered and it was shown that the external sample difference between SED and quantum mechanics had a small signal to noise ratio in typical parametric oscillators. The presence of intense local oscillators means that the results could be seen as providing the basis for a macroscopic experimental test of quantum mechanics against SED.
## acknowledgements
DTP would like to thank Professor Gerard Milburn, Dr Howard Wiseman, Dr Karen Kheruntsyan, Professor Brian Orr, Michael Gagen and Cynthia Freeman for their assistance with the paper. He would also like to thank The University of Queensland for its financial support. WJM acknowledges support from the Australian Research Council. |
warning/0003/gr-qc0003029.html | ar5iv | text | # A New Way to Make Waves
## I Introduction
I will describe a new way to “make waves”. Wave phenomena is treated mathematically by hyperbolic equations. Einstein’s theory of general relativity is the supreme example. It treats gravity as a distortion in the geometry of space-time. Gravitational waves are ripples in the geometry which travel through space changing the shape and size of objects in their path. At this historic time, Einstein’s theory is providing the basis for a new way to observe the distant universe. A worldwide network of gravitational wave observatories is under construction. It is designed to detect the gravitational waves produced throughout the cosmos by collisions between black holes.
The new way to make waves which I will describe was developed to simulate the production of gravitational waves by using a computer to solve Einstein’s equations. But the approach applies to any wave phenomena. In order to get the fundamental ideas across, I present them in terms of very simple systems. Only near the end will I give in to the urge to tell you about the application to black holes. However, even in simple examples, the ideas are easiest to understand in terms of space-time language and pictures.
The simplest hyperbolic equation is the 1-dimensional advection equation
$$(_t+c_x)\mathrm{\Phi }=0.$$
(1)
It has the general solution $`\mathrm{\Phi }(t,x)=F(xct)`$, which describes a wave traveling with velocity $`c`$ with initial waveform $`F(x)`$ at $`t=0`$. The curves $`x=x_0+ct`$ are called the characteristics of the equation. They are the paths in space-time (worldlines) along which you must move to “ride the wave” in the sense that $`\mathrm{\Phi }=const`$ along a characteristic worldline. Since $`F`$ is arbitrary, we can choose $`F=0`$ for $`x<0`$ and $`F=1`$ for $`x>0`$ to obtain a wave with a shock front along the characteristic $`x=ct`$. In modern parlance, this shock wave carries one bit of information. It illustrates the general way information is propagated by waves along characteristics in any hyperbolic system.
The standard way to solve a hyperbolic equation such as Eq. (1) is by means of the Cauchy problem. One poses initial Cauchy data $`\mathrm{\Phi }(t=0,x)`$ in some spatial region $`\mathrm{\Omega }`$ and then evolves the solution in time. The characteristics causally relate the initial data to a unique evolution throughout some domain of dependence $`D(\mathrm{\Omega })`$ in space-time. As a result, shock fronts, which represent a sudden signal, can only occur across characteristics. Hyperbolic equations lead to ordinary differential equations along the characteristics which govern the propagation of shock discontinuities. That makes it important for the purpose of numerical simulation to enforce the propagation along characteristics as extensively as possible. In complicated cases, where the wave velocity has functional dependence $`c(\mathrm{\Phi },x)`$ and an analytic solution is not possible, there still exists a well established theory of the general properties of hyperbolic systems based upon characteristics .
The method of characteristics is a computational algorithm for evolving Cauchy data based upon this theory. As applied to the advection equation (1), the characteristic velocities are calculated at a given discretized time $`t_n=n\mathrm{\Delta }t`$ and the algorithm used to determine the field at time $`t_{n+1}`$ by requiring that $`\mathrm{\Phi }`$ be constant along the characteristic curve from a space-time point $`P`$ at time level $`n`$ to point $`Q`$ at time level $`n+1`$, i.e.
$$\mathrm{\Phi }_Q=\mathrm{\Phi }_P,$$
(2)
where $`x_Q=x_P+c\mathrm{\Delta }t`$. When the wave velocity is not constant in either space or time, the points $`P`$ and $`Q`$ along the characteristics do not in general lie exactly on spatial grid points, so that interpolations are necessary. But, because there are no sudden changes along characteristics (only across them), such interpolations give excellent accuracy by riding the wave like a surfer.
The method of characteristics for the Cauchy problem extends to the generalization of Eq. (1) to any symmetric hyperbolic system of first differential order equations in a multi-dimensional space with many evolution variables and with source terms for the creation of waves . However, there is another classification of hyperbolic systems, which is more familiar to physicists, based upon second order differential equations whose principal part has the form
$$g^{\alpha \beta }_\alpha _\beta \mathrm{\Phi }=0,$$
(3)
where $`\alpha `$ is a space-time index, i.e. $`x^\alpha =(t,x,y,z)`$ for 4-dimensional space-time, and a sum over the values of the repeated indices is implied. Equation (3) is classified as hyperbolic if the symmetric matrix $`g^{\alpha \beta }`$ has one negative eigenvalue and its remaining eigenvalues are positive. Displacements $`x^\alpha x^\alpha +dx^\alpha `$ along characteristic worldlines satisfy
$$g_{\alpha \beta }dx^\alpha dx^\beta =0,$$
(4)
where $`g_{\alpha \beta }`$, called the metric tensor, is the inverse matrix to $`g^{\alpha \beta }`$. The simplest example is the wave equation in one spatial dimension,
$$[\frac{1}{c^2}_t^2_x^2]\mathrm{\Phi }=0,$$
(5)
with the displacement along the characteristics satisfying $`c^2dt^2dx^2=0`$. By rewriting Eq. (5) in the form
$$(\frac{1}{c}_t_x)(\frac{1}{c}_t+_x)\mathrm{\Phi }=0,$$
(6)
a comparison with Eq. (1) shows that the general solution is $`\mathrm{\Phi }(t,x)=F(xct)+G(x+ct)`$. There are two characteristics $`x=x_0\pm ct`$ through each spatial point $`x_0`$. Information is propagated along these characteristics which now criss-cross the space-time. The conventional method of characteristics could be applied to Eq. (6) by the standard technique of reducing it to a coupled system of first order differential equations. But it is simpler to integrate Eq.(6) over a parallelogram $`\mathrm{\Sigma }`$ in space time whose sides are characteristics meeting at the corners $`P`$, $`Q`$, $`R`$ and $`S`$. This gives the 2-dimensional version of Eq. (2)
$$\mathrm{\Phi }_Q\mathrm{\Phi }_P+\mathrm{\Phi }_R\mathrm{\Phi }_S=0,$$
(7)
as illustrated in Fig. 1 for the case $`c=1`$ where the characteristics $`x=x_0\pm t`$ have $`45^o`$ slope in space-time. By using Eq. (7), the method of characteristics can be implemented as a computational algorithm for the Cauchy evolution of Eq. (5) by positioning the point $`Q`$ at time level $`n+1`$, the points $`P`$ and $`S`$ at time level $`n`$ and the point $`R`$ at time level $`n1`$. On a $`(t,x)`$ computational grid satisfying $`\mathrm{\Delta }x=c\mathrm{\Delta }t`$, the characteristics pass through diagonal grid points. The Cauchy data, for the wave equation, which consists of the initial values $`\mathrm{\Phi }(t=0,x)`$ and $`_t\mathrm{\Phi }(t=0,x)`$, is used to initialize an iterative evolution scheme at time levels $`n=0`$ and $`n=1`$. An exceptional feature is that the one-dimensional wave equation with constant wave velocity can be evolved without error by using Eq. (7) as a finite difference equation on such a grid. For a variable velocity, the points $`P`$, $`Q`$, $`R`$ and $`S`$ cannot be placed exactly on the grid and interpolations are necessary.
For the wave equation in two or more spatial dimensions, the manner in which characteristics determine domains of dependence and lead to propagation equations is qualitatively the same. The major difference is that an infinite number of characteristics now pass through each point. For the 3-dimensional wave equation,
$$[\frac{1}{c^2}_t^2_x^2_y^2_z^2]\mathrm{\Phi }=0$$
(8)
the characteristics which pass through the point $`(x_0,y_0,z_0)`$ at time $`t_0`$ are straight lines which trace out an expanding spherical wavefront of radius $`c(tt_0)`$ at a given time $`t`$; or, from a space-time point of view, these characteristics trace out the 3-dimensional characteristic cone
$$(xx_0)^2+(yy_0)^2+(zz_0)^2c^2(tt_0)^2=0.$$
(9)
In electrodynamics or relativity, the characteristics are light rays and this is called the light cone. I will often use that language here. The analogue in hydrodynamics is the Mach cone.
The future (past) light cone consists of the radially outward (inward) characteristics parameterized by $`t>t_0`$ ($`t<t_0`$). There is a 2-parameter set of characteristics through each point corresponding to the sphere of angular directions $`(\theta ,\varphi )`$ at that point. This leads to some arbitrariness in formulating an evolution algorithm for Cauchy data based upon the method of characteristics. There are an infinite number of characteristics and associated propagation equations which can be used to evolve Cauchy data from time $`t_0`$ to time $`t_0+\mathrm{\Delta }t`$.
For a practical numerical scheme, it is thus necessary either to average these propagation equations appropriately over the sphere of characteristic directions at or select out some finite number of characteristics whose propagation equations comprise a nonredundant set. The latter approach has been successfully carried out by Butler . For the case of plane flow of an inviscid fluid (a problem in two spatial dimensions), he formulates an algorithm based upon four “preferred” characteristics. In this problem, the geometry is further complicated because the characteristics are dynamically dependent upon the fluid variables, in contrast to the essentially time independent characteristic cones of Eq. (9). In the numerical scheme, the characteristics must themselves be determined by some finite difference approximation.
That summarizes the Cauchy problem and its solution by the method of characteristics.
## II Characteristic evolution
I now present a procedure which avoids the arbitrariness and awkwardness of the method of characteristics in a 3-dimensional Cauchy evolution. It is based upon a characteristic initial value approach rather than a Cauchy approach. In order to understand the distinction, it is essential to view space-time as 4-dimensional, with initial data given on a three-dimensional hypersurface. The Cauchy approach is based upon evolution of initial data given on a spatial hypersurface in space-time, i.e. points of space at the same instant of time $`t=t_0`$, to data at a later time $`t=t_0+\mathrm{\Delta }t`$. In the characteristic approach, data is initially posed on an outgoing characteristic cone (light cone), which defines a hypersurface at constant retarded time $`u=u_0`$. Characteristic evolution then proceeds iteratively to characteristic cones $`u_n=u_0+n\mathrm{\Delta }u`$. The numerical grid is intrinsically based upon these outgoing characteristic hypersurfaces. The algorithm for evolving from retarded time $`u_n`$ to $`u_{n+1}`$ is based upon characteristics which are uniquely and intrinsically picked out by the geometry of the retarded time hypersurfaces.
This characteristic evolution procedure radically differs from Cauchy evolution. It uses concepts developed in the 1960’s for studies of general relativity . These were prompted by the inability of the major mathematical tools such as Green’s functions and Fourier analysis to overcome the difficulties posed by nonlinearity of the equations and the ambiguity of the general coordinate freedom in the theory. Its success later motivated a more extensive mathematical treatment of the characteristic initial value problem . These new techniques were especially designed for investigation of gravitational waves. A computational algorithm based on this approach has been the most successful in simulating the production of gravitational waves from black holes.
The approach has two novel ingredients:
* the use of characteristic hypersurfaces to formulate a characteristic initial value problem (CIVP) and
* the use of compactification methods to describe on a finite grid the waves propagating to infinity.
Characteristic hypersurfaces are 3-dimensional sets traced out in space-time by characteristic worldlines or, equivalently, by their wavefronts. They provide a natural coordinate system to describe waves . It is very fruitful to use an initial value scheme which describes time evolution by means of the retarded time coordinate defined by characteristic hypersurfaces as a substitute for the familiar Cauchy scheme based upon constant time hypersurfaces. As will be shown, this approach uses a completely different form of the mathematical equations and the free initial data.
Compactification methods provide a rigorous description of the radiation field of a source observed in the asymptotic limit of going to infinity along a characteristic worldline.. The key idea is to introduce a new coordinate which ranges over values from 0 to 1 as the actual distance from the source ranges from 0 to $`\mathrm{}`$. The hyperbolic equations are rewritten in terms of these new coordinates. Asymptotic behavior at the “points at infinity” can then be studied in terms of the new coordinate which ranges over finite values. In this way, the concept of the radiation zone as an asymptotic limit at infinity is given rigorous meaning. Characteristic hypersurfaces are important here since waves travel to infinity along characteristic worldlines, not along Cauchy hypersurfaces of constant time. Even for field equations as complicated as those of general relativity, this procedure provides a finite geometrical description of waves travelling to infinity. The limit points at infinity form a boundary to the compactified space-time, which I will refer to as radiative infinity. At a given retarded time, this boundary has the topology of a sphere, representing a sphere of observers at infinity.
It should be emphasized that radiative infinity differs drastically from spatial infinity (the limit of going to infinity holding time constant). Early considerations of compactifying infinite space for computational purposes were discarded because they were based upon spatial infinity . An attempt to cover infinite space this way by a finite grid at constant time fails in a hyperbolic problem because there is necessarily an infinite physical distance between a grid point at infinity and its neighbors. This makes it impossible on a spatial grid at fixed time to resolve radiation with finite wave length which propagates to infinity. In contrast, for a grid constructed on a characteristic hypersurface, the grid points ride the wave without noticing its finite wavelength in the approach to radiative infinity.
In terms of characteristic coordinates $`u=ctx`$ and $`v=ct+x`$, Eq. (6) becomes
$$_u_v\mathrm{\Phi }=0,$$
(10)
Here $`\mathrm{\Psi }=_v\mathrm{\Phi }`$ satisfies the propagation equation
$$_u\mathrm{\Psi }=0$$
(11)
along the characteristics in the $`u`$-direction. This is the essence of how the use of characteristic coordinates simplifies the treatment of waves.
These ideas provide the physical basis for a new computational algorithm . I will illustrate how it applies to a nonlinear version of the wave equation in 3 spatial dimensions. However, this evolution algorithm can be taken over intact to other hyperbolic physical systems, including electromagnetic fields, the Yang-Mills gauge fields of elementary particle physics, as well as general relativity, because of the common mathematical structure of these theories as second differential order hyperbolic equations. Although it has not been explored how this approach might be implemented in the case of a first differential order hyperbolic system, such as hydrodynamics, the general ideas should be applicable.
## III Characteristic initial value problem for nonlinear waves
As a simple illustration of the characteristic initial value problem, consider the nonlinear scalar wave equation (SWE) in 3-spatial dimensions, which we write in spherical polar coordinates $`(t,r,\theta ,\varphi )`$ as
$`[{\displaystyle \frac{1}{c^2}}_t^2_x^2_y^2_z^2]\mathrm{\Phi }={\displaystyle \frac{1}{c^2}}_t^2\mathrm{\Phi }{\displaystyle \frac{1}{r}}_r^2(r\mathrm{\Phi })+{\displaystyle \frac{L^2\mathrm{\Phi }}{r^2}}=S(\mathrm{\Phi })`$ (12)
$`r=\sqrt{(}x^2+y^2+z^2),`$ (13)
where $`c`$ is the wave velocity, $`L^2`$ denotes the standard angular momentum operator
$$L^2\mathrm{\Phi }=\frac{_\theta (\mathrm{sin}\theta _\theta \mathrm{\Phi })}{\mathrm{sin}\theta }\frac{_\varphi ^2\mathrm{\Phi }}{\mathrm{sin}^2\theta }.$$
(14)
and $`S(\mathrm{\Phi })`$ represents a nonlinear source term. Rather than using ordinary time $`t`$, characteristic evolution uses the “retarded time” coordinate $`u=ctr`$. The outgoing radial characteristics are the curves of constant $`u`$, $`\theta `$ and $`\varphi `$. These are the curves in the $`r`$-direction, holding $`u=const`$, shown in the space-time Fig. 1. In $`(u,r,\theta ,\varphi )`$ coordinates, the SWE (12) takes the form
$$2_u_rg=_r^2g\frac{L^2g}{r^2}+rS.$$
(15)
where $`g=r\mathrm{\Phi }`$.
The striking feature about Eq.(15) is that it is only first differential order in retarded time $`u`$, unlike the more familiar form of the SWE (12) which is of second differential order in time $`t`$. When data is given on a characteristic initial hypersurface $`u=u_0`$, we need only specify the initial value of the field $`\mathrm{\Phi }`$, and then use Eq. (15) to calculate its retarded time derivative $`_u\mathrm{\Phi }`$ in order to evolve the initial data. This is in contrast to the conventional Cauchy scheme, where both $`\mathrm{\Phi }`$ and $`_t\mathrm{\Phi }`$ must be supplied at initial time $`t=t_0`$ and Eq. (12) is the used to compute the second time derivative $`_t^2\mathrm{\Phi }`$ in order to evolve the data.
In a computational implementation of the CIVP, rather than finite differencing Eq. (15) directly, it is advantageous to first convert it into an integral equation which is subsequently discretized. Using both the outgoing characteristic coordinate $`u`$ and the ingoing characteristic coordinate $`v=ct+r`$, the SWE (15) takes the form
$$4_u_vg=\frac{L^2g}{r^2}+rS(\frac{g}{r})$$
(16)
where $`g=r\mathrm{\Phi }`$. In the $`u`$-$`v`$ plane formed by fixing the angular coordinates $`(\theta ,\varphi )`$, we construct a parallelogram $`\mathrm{\Sigma }`$ made up of incoming and outgoing radial characteristics which intersect at vertices $`P,Q,R,S`$ as depicted in Fig. 1. By integrating Eq. (16) over the area $`\mathrm{\Sigma }`$ bounded by these vertices, we may establish the identity
$$g_Q=g_P+g_Sg_R+\frac{1}{2}_\mathrm{\Sigma }𝑑u𝑑r\left[\frac{L^2g}{r^2}+rS(\frac{g}{r})\right].$$
(17)
This simple identity is the 3-dimensional analogue of Eq. (7), adapted to include a source term. It is the starting point for an evolution algorithm which incorporates the essential role that characteristics play in the SWE.
In order to study the far field wave behavior, we transform this equation to the new radial coordinate
$$x=r/(1+r),\mathrm{\hspace{0.17em}\hspace{0.17em}0}x1.$$
(18)
This serves to map an infinite radial domain into a finite coordinate region, and assigns infinitely distant radial points to the edge of the coordinate patch (radiative infinity) at $`x=1`$, where the radiation signal can be identified.
### A Numerical Algorithm
To develop a discrete evolution algorithm, we work on the lattice of points
$`u_n`$ $`=`$ $`n\mathrm{\Delta }u`$ (19)
$`x_i`$ $`=`$ $`i\mathrm{\Delta }x`$ (20)
$`\zeta _{j,k}`$ $`=`$ $`(j+ik)\mathrm{\Delta }\varphi `$ (21)
where a complex stereographic coordinate $`\zeta `$ is used to cover the the sphere in two patches (North and South) in order to avoid the polar singularities of the spherical coordinates $`(\theta ,\varphi )`$. We denote the field at these sites by
$$g_{ijk}^n=g(u_n,x_i,\zeta _{j,k}).$$
(22)
(We will generally suppress the angular indices $`j`$ and $`k`$.)
With respect to the $`(u,x)`$ coordinate grid, it is not possible to place the corners P, Q, R and S at grid points since the slope of the characteristics in the compactified $`x`$ coordinate depends upon location. As a consequence, the field $`g`$ at these points must be interpolated from neighboring grid points. The essential feature of the placement of the parallelogram on the grid is that the sides formed by the ingoing characteristics intersect adjacent $`u`$-hypersurfaces at equal but opposite $`x`$-displacements from the neighboring grid points, as illustrated in Fig. 2. The field values at the vertices of the parallelogram are obtained by quadratic interpolation. Cancellations between the interpolation errors at the four vertices yields the accuracy
$$g_Qg_Pg_S+g_R=G_QG_PG_S+G_R+O((\mathrm{\Delta }x)^3\mathrm{\Delta }u),$$
(23)
where $`G`$ represents the exact analytic solution.
The integral in Eq. (17) can be evaluated by treating the integrand as a constant over the parallelogram, with value at the center. The radial coordinate of the point at the center is $`r_c=(r_P+r_S)/2`$. To compute the nonlinear term, the value of $`g`$ at $`r_c`$ is taken as the average $`g_c=(g_P+g_S)/2`$, with $`g_P`$ and $`g_S`$ evaluated from second-order linear interpolations over adjacent points on the grid. The angular derivatives in Eq. (17) are replaced with standard second-order-accurate finite difference approximations. $`L^2g`$ is calculated on the grid points, and the same interpolation procedure is used to obtain the value of $`L^2g_c`$. The integral term is then approximated by
$`{\displaystyle _\mathrm{\Sigma }}[L^2g+r^3S(g/r)]𝑑u𝑑r/r^2=[L^2g_c+r_c^3S(g_c/r_c)]{\displaystyle _\mathrm{\Sigma }}𝑑u𝑑r/r^2`$ (24)
$`=2\mathrm{log}\left({\displaystyle \frac{r_Qr_R}{r_Pr_S}}\right)\left[L^2g_c+r_c^3S({\displaystyle \frac{g_c}{r_c}})\right],`$ (25)
where the integrand is accurate to second order in $`\mathrm{\Delta }x`$ and $`\mathrm{\Delta }\varphi `$. The resulting finite difference equation
$`g_i^{n+1}(x_Q+x_Px_{i2}x_{i1})=`$ (30)
$`2g_{i1}^{n+1}(x_Q+x_Px_{i2}x_i)g_{i2}^{n+1}(x_Q+x_Px_{i1}x_i)`$
$`+\{g_{i+1}^n(x_S+x_Rx_{i1}x_i)2g_i^n(x_S+x_Rx_{i1}x_{i+1})`$
$`+g_{i1}^n(x_S+x_Rx_ix_{i+1})\}{\displaystyle \frac{(x_Sx_R)}{(x_Qx_P)}}`$
$`+{\displaystyle _\mathrm{\Sigma }}{\displaystyle \frac{drdu}{r^2}}\left[L^2(g_c)+r_c^3S({\displaystyle \frac{g_c}{r_c}})\right]{\displaystyle \frac{(\mathrm{\Delta }x)^2}{(x_Qx_P)}}.`$
relates values of $`g_i^{n+1}`$ with values at neighboring grid points which are either earlier in retarded time ($`g_{i1}^n,g_i^n,g_{i+1}^n`$), or else contemporary but located at smaller radius ($`g_{i2}^{n+1},g_{i1}^{n+1}`$). Consequently, it is possible to move through the grid, computing $`g_i^{n+1}`$ explicitly by an orderly march. This is achieved by starting at the origin at time $`u_{n+1}`$. Field values of $`g=r\mathrm{\Phi }`$ vanish there. Step outward to the next radial point, using Eq. (30) for all angular sites on the grid, and iterate this march out to radiative infinity thus updating the characteristic cone at $`u_{n+1}`$ and completing one retarded time step. This march is then iterated in retarded time.
The algorithm steps $`g`$ radially outward one cell with a local error of fourth order in grid size. This leads to second order global accuracy which is confirmed by convergence tests using known analytic solutions . A complete specification of the algorithm would require a description of how the startup procedure at the origin is handled and how stereographic coordinates are used to compute angular derivatives in a smooth way. Details of the use of North and South stereographic grids are given in Ref. . Physical behavior of nonlinear waves is treated in Ref. . Construction of an exact nonspherical solution for an $`S=\mathrm{\Phi }^3`$ self-interaction allows calibration of the algorithm in the nonlinear case where physical singularities form. The predicted second order accuracy is confirmed right up to the formation of the singularity. Other choices of nonlinear potential allow simulation of solitons.
The Courant-Friedrichs-Lewy (CFL) condition that the numerical domain of dependence contain the physical domain of dependence (determined by the characteristics) is a necessary condition for convergence of a finite difference algorithm. For a grid point at $`(u,r,\theta )`$, there are three critical grid points, at $`(u\mathrm{\Delta }u,r+\mathrm{\Delta }r,\theta )`$ and $`(u\mathrm{\Delta }u,r\mathrm{\Delta }r,\theta \pm \mathrm{\Delta }\theta )`$, which must lie inside its past characteristic cone. These gives rise to the inequalities $`\mathrm{\Delta }u<2\mathrm{\Delta }r`$ and $`\mathrm{\Delta }u<\mathrm{\Delta }r+(\mathrm{\Delta }r^2+r^2\mathrm{\Delta }\theta ^2)^{1/2}`$. At large $`r`$, the second inequality becomes $`\mathrm{\Delta }u<r\mathrm{\Delta }\theta `$ and the Courant limit on the time step is essentially the same as for a Cauchy evolution algorithm. However, near the vertex of the cone, the second inequality gives a stricter condition
$$\mathrm{\Delta }u<K\mathrm{\Delta }r\mathrm{\Delta }\theta ^2,$$
(31)
where the value of $`K`$ depends upon the start up procedure at the vertex. For the scalar wave equation, these stability limits were confirmed by numerical experiments and it was found that $`K4`$.
Local von Neumann stability analysis leads to no constraints on the algorithm. This may seem surprising because no analogue of a CFL condition on the time step arises. It can be understood in the following vein. The local structure of the code is implicit, since it involves 3 points at the upper time level. Implicit algorithms do not necessarily lead to a CFL condition. However, the algorithm is explicit in the way that the evolution starts up as an outward radial march from the origin. It is this startup procedure that introduces a CFL condition.
Operating within the CFL limit, the algorithm gives a stable, globally second order accurate evolution on a compactified grid . (In some nonlinear applications, artificial dissipation is necessary for stability .) Radiative infinity behaves as a perfectly transmitting boundary so that no radiation is reflected back into the system. Numerical evolution satisfies a conservation law relating the loss of energy to the radiation flux at infinity.
That summarizes the characteristic initial value problem and its implementation as a new computational approach to simulate waves.
## IV Cauchy-characteristic matching
Characteristic evolution has many advantages over Cauchy evolution. Its one disadvantage is caused by the existence of either (i) caustics where neighboring characteristics focus or (ii), a milder version of this, cross-over points where two distinct characteristics collide. The vertex of the characteristic cone is a highly symmetric example of a point caustic where a complete sphere of characteristics focus. I have already discussed how a point focus gives rise to a strong limitation imposed by the CFL condition.
Cauchy-characteristic matching (CCM) is a way to avoid such limitations by combining the strong points of characteristic and Cauchy evolution in formulating a global evolution. Here I illustrate the application of CCM to the nonlinear wave equation (12). This problem requires boundary conditions at infinity which ensure that the total energy and the energy loss due to radiation are both finite. In a 3-dimensional problem, these are the conditions responsible for the proper $`1/r`$ asymptotic decay of the radiation fields. However, for practical purposes, in the computational treatment of such a system by the Cauchy problem, an outer boundary is artificially established at some large but finite distance. Some condition is then imposed upon this boundary in an attempt to approximate the proper asymptotic behavior at infinity. Such an artificial boundary condition (ABC) typically causes partial reflection of the outgoing wave back into the system , which contaminates the accuracy of the evolution and the radiated signal. Furthermore, nonlinear wave equations often display backscattering so that it may not be correct to try to entirely eliminate incoming radiation from the numerical solution. The errors introduced by ABC’s are of an analytic origin, essentially independent of the computational discretization. In general, a systematic reduction of the error can only be achieved by simultaneously refining the grid and moving the computational boundary to a larger radius, which is computationally very expensive for three-dimensional simulations. CCM provides a global solution which does not introduce error at the analytic level.
For linear wave problems, a variety of ABC’s have been proposed. For recent reviews, see Ref’s. . During the last two decades, local ABC’s in differential form have been extensively employed by several authors with varying success. Some local ABC’s have been derived for the linear wave equation by considering the asymptotic behavior of outgoing solutions ; this approach may be regarded as a generalization of the Sommerfeld outgoing radiation condition. Although such ABC’s are relatively simple to implement and have a low computational cost, their final accuracy is often limited because their simplifying assumptions are rarely met in practice . Systematic improvement of the accuracy of local ABC’s can only be achieved by moving the computational boundary to a larger radius.
The disadvantages of local ABC’s have led to implementation of nonlocal ABC’s based on integral representations of the infinite domain problem . Even for problems where the Green’s function is known and easily computed, such approaches were initially dismissed as impractical ; however, the rapid increase in computer power has made it possible to implement nonlocal ABC’s for the linear wave equation even in 3 space dimensions . For a linear problem, this can yield numerical solutions which converge to the exact infinite domain problem as the grid is refined, keeping the artificial boundary at a fixed distance. However, due to nonlocality, the computational cost per time step usually grows at a higher power of grid size ($`O(N^4)`$ per time step in a 3-dimensional problem with $`O(N^3)`$ spatial grid points) than in a local approach , which is demanding even for today’s supercomputers. Further, the applicability of current nonlocal ABC’s is restricted to problems where nonlinearity may be neglected near the grid boundary .
To my knowledge, only a few works have been devoted to the development of ABC’s for strongly nonlinear problems . In practice, nonlinear problems are often treated by linearizing the governing equations in the far field , using either local or nonlocal linear ABC’s . Besides introducing an approximation at the analytical level, this procedure requires that the artificial boundary be placed sufficiently far from the strong-field region, which sharply increases the computational cost in multidimensional simulations. There seems to be no currently available ABC which is able to produce numerical solutions which converge (as the discretization is refined) to the infinite domain exact solution of a strongly nonlinear 3-dimensional wave problem, keeping the artificial boundary at a fixed location.
For such nonlinear problems, CCM produces an accurate solution out to radiative infinity with effort $`O(N^3)`$ per time-step. CCM increases the total computational cost only by a factor $`2`$ with respect to a pure Cauchy algorithm with a local ABC. The use of numerical methods based upon matching a characteristic initial-value formulation and a Cauchy formulation can effectively remove the above difficulties associated with a finite computational boundary. There is no need to truncate space-time at a finite distance from the sources, since compactification of the radial coordinate makes it possible to cover the whole space-time with a finite grid. In this way, the true radiation zone signal may be computed. Although the characteristic formulation has stability limitations in interior region where the characteristic hypersurfaces can develop caustics, it proves to be both accurate and computationally efficient in the treatment of the exterior, caustic-free region.
CCM is a new approach to global numerical evolution which is free of error at the analytic level. The characteristic algorithm provides the outer boundary condition for the interior Cauchy evolution, while the Cauchy algorithm supplies the inner boundary condition for the characteristic evolution. Since CCM consists of discretizing an exact analytic treatment of the radiation from source to radiative infinity, it generates numerical solutions which converge to the exact analytic solution of the radiating system even in the presence of strong nonlinearity. Thus, any desired accuracy can be achieved by refining the grid, without moving the matching boundary. In practice, the method performs extremely well even at moderate resolutions.
### A CCM for nonlinear waves
As an illustration of the computational implementation of CCM , consider again the nonlinear 3-dimensional wave equation (12). For simplicity, set the velocity $`c=1`$. In the standard computational implementation of the Cauchy problem for (12), initial data $`\mathrm{\Phi }(t_0,x,y,z)`$ and $`_t\mathrm{\Phi }(t_0,x,y,z)`$ are assigned and evolved in a bounded spatial region, with some ABC imposed at the computational boundary. In a characteristic initial-value formulation the SWE (12) is reexpressed in the form of Eq. (15) in terms of $`g=r\mathrm{\Phi }`$, using standard spherical coordinates and a retarded time coordinate $`u=tr`$. The initial data $`g(u_0,r,\theta ,\varphi )`$, on an initial outgoing characteristic cone $`u=u_0`$ is then evolved globally out to radiative infinity.
In CCM, (12) is solved in an interior region $`rR_m`$ using a Cauchy algorithm, while a characteristic algorithm solves the retarded coordinate version (15) for $`rR_m`$. The matching procedures ensure that, in the continuum limit, $`\mathrm{\Phi }`$ and its gradient are continuous across the interface $`r=R_m`$. This is a requirement for any consistent matching algorithm, since a discontinuity in the field or its gradient could act as a spurious boundary source, contaminating both the interior and exterior evolutions.
For technical simplicity, I illustrate the details of the method here for spherically symmetric waves but it has been successfully implemented in fully nonlinear 3-dimensional problems without symmetry. In a 3-dimensional problem without symmetry, the characteristic evolution is carried out on an exterior spherical grid, while the Cauchy evolution uses a Cartesian grid covering the interior spherical region. Although a spherical grid could also be used in the interior, a Cartesian grid avoids the necessity of cumbersome numerical procedures to handle the singularity of spherical coordinates at the origin. A Cartesian discretization in the interior and a spherical discretization in the exterior are the coordinates natural to the geometries of the two regions. However, this makes the treatment of the interface somewhat involved; in particular, guaranteeing the stability of the matching algorithm requires careful attention to the details of the inter-grid matching. Nevertheless, there is a reasonably broad range of discretization parameters for which CCM is stable .
With the substitution $`G=r\varphi `$ and the use of spherical coordinates, the spherically symmetric version of the SWE (12) reduces to the 1-dimensional wave equation
$$_{tt}G=_{rr}G+rS.$$
(32)
The initial Cauchy data is $`G(t_0,r)`$ and $`_tG(t_0,r)`$ in the region $`0rR_m`$. Together with the regularity condition $`G(t,0)=0`$, these data determine a unique solution in the domain of dependence $`D_1^{}`$ indicated in Fig. 3. The outer boundary of the domain of dependence is the ingoing radial characteristic $`C_1^{}`$ described by $`r=R_mt+t_0`$. The solution cannot be constructed throughout the complete interior region $`rR_m`$ without additional information, which can be furnished by giving the value of $`G`$ on the outgoing characteristic $`C_{0^+}`$ described by $`r=R_m+tt_0`$ (see Fig. 3). In terms of the coordinates $`u=tr`$ and $`r`$, $`C_{0^+}`$ is described by $`u=t_0R_m`$ In these coordinates, expressing $`g(u,r)=G(u+r,r)`$, the spherically symmetric version of the wave equation (15) is
$$2_{ur}g=_{rr}g+rS.$$
(33)
A unique solution of (33) is determined by characteristic initial data consisting of the value of $`g`$ on the initial outgoing characteristic $`C_{0^+}`$ and on the ingoing characteristic $`C_1^{}`$. These data determine the solution uniquely throughout the future of $`C_1^{}`$ and $`C_{0^+}`$, i.e. the region $`D_{1^+}`$ in Fig. 3.
The matching scheme proceeds as shown in Fig. 3. First, initial Cauchy data are evolved from $`t_0`$ to $`t_1`$ throughout the region $`D_1^{}`$, which is in its domain of dependence. Next, this induces characteristic data on $`C_1^{}`$ which combined with the initial characteristic data on $`C_{0^+}`$ allows a characteristic evolution throughout the region $`D_{1^+}`$, bounded in the future by the characteristic $`C_{1^+}`$. The solution determined from this initial stage induces Cauchy data at time $`t_1`$ in the region $`rR_m`$, inside the matching boundary. This process can then be iterated to carry out the entire future evolution of the system.
CCM is a discretized version of this scheme in which the criss-cross pattern of characteristics inside the radius $`R_m`$ is at the scale of a grid spacing. The discretized evolution algorithm consists of the following steps (see Fig. 4):
Step 1. Cauchy evolution. The interior integration scheme is implemented on a uniform spatial grid $`r_i=i\mathrm{\Delta }r`$ ($`0iM`$) with outer radius $`R_B=M\mathrm{\Delta }r`$. We discretize (32) using the standard second-order finite difference scheme
$$\frac{G_i^{n+1}2G_i^n+G_i^{n1}}{(\mathrm{\Delta }t)^2}=\frac{G_{i+1}^n2G_i^n+G_{i1}^n}{(\mathrm{\Delta }r)^2}+r_iS_i^n,$$
(34)
where $`G_i^n=G(t_n,r_i)`$, $`S_i^n=S(t_n,r_i)`$, and $`t_n=t_0+n\mathrm{\Delta }t`$. The interior evolution is initialized by evaluating $`G_i^0`$ and $`G_i^1`$ $`(0iM)`$ to second order accuracy from the Cauchy initial data. In the $`n`$th time step, (34) is used to compute $`G_i^{n+1},1iM`$ in terms of field values $`G_i^{n1}`$ and $`G_i^n`$. The regularity of $`\mathrm{\Phi }`$ at $`r=0`$ implies that $`G_0^n=0`$ for all $`n`$. The boundary values $`G_{M+1}^n`$ which are required by (34) are supplied by the matching procedure (step 3).
Step 2. Characteristic evolution. The characteristic algorithm is implemented on a uniform grid based on the dimensionless compactified radial coordinate
$$\eta =1\frac{1}{1+r/R_m},\frac{1}{2}\eta 1$$
(35)
so that points at radiative infinity (corresponding to $`\eta =1`$) are included in the grid. In order to include one of the technical problems in matching an interior Cartesian grid to an exterior spherical grid, let there be a small gap between the outer radius $`R_B`$ of the Cauchy grid and the matching radius $`R_m`$ (which is also the inner radius of the characteristic grid). In 3-dimensional Cartesian-spherical matching the outermost Cauchy grid points and the innermost characteristic grid points are necessarily distinct. This is represented here by a gap $`R_mR_B=\kappa \mathrm{\Delta }r`$, where $`\kappa 0`$ is an arbitrary parameter. The characteristic grid consists of the uniformly spaced points $`\eta _\alpha =\frac{1}{2}+\alpha \mathrm{\Delta }\eta `$ ($`0\alpha N_\eta `$), where $`\mathrm{\Delta }\eta =(2N_\eta )^1`$. The retarded time levels $`u=u_n`$ for the characteristic evolution are chosen so that $`u=u_n`$ intersects the time level $`t=t_n`$ of the Cauchy evolution at the matching radius; therefore, $`u_n=t_nR_m`$ and $`\mathrm{\Delta }u=\mathrm{\Delta }t`$. We denote by $`g_\alpha ^n`$ the value of $`g`$ at $`\eta =\eta _\alpha `$, $`u=u_n`$. The initial characteristic data consist of $`g_\alpha ^0`$, $`0\alpha N_\eta `$.
In the $`n`$th iteration of the evolution, we compute the field values at the grid points with $`u=u_n`$ using the values of $`g_\alpha ^{n1}`$, which are known either from initialization or from the previous iteration. As already described, this is done by the characteristic marching algorithm based on the integral identity (17).
At the inner boundary of the characteristic grid ($`\alpha =1`$), the previous scheme must be slightly modified, since $`g_{\alpha 2}^n`$ is not defined. For this initial step, $`PQRS`$ is chosen so that $`\eta _P=\eta _0`$, $`\eta _Q=\eta _1`$, and $`g_R`$, $`g_S`$ are approximated by quadratic interpolation in terms of $`g_0^{n1}`$, $`g_1^{n1}`$, $`g_2^{n1}`$, which have already been computed. Besides these field values, the final evaluation of $`g_1^n`$ still requires the value of $`g_0^n`$, which is supplied by the matching procedure (step 3).
Step 3. Matching. Numerous schemes are possible in the case of spherically symmetric matching. Here I describe one which is stable for a wide range of gap sizes, $`0\kappa 2`$ and works in the more complicated situation with an interior Cartesian grid and an exterior spherical grid.
The required boundary values $`G_{M+1}^n`$ and $`g_0^n`$ are computed by radial interpolations at constant $`t`$, using the field values at points $`A`$, $`B`$, $`E`$, and $`F`$ in Fig. 4. The first two of these field values are already known at the $`n`$th step, while the last two can be obtained by cubic radial interpolations along the previously evolved characteristics $`u=u_{n1}`$ and $`u=u_{n2}`$, respectively. At the initial step, point $`F`$ lies on the characteristic $`u=u_1=\mathrm{\Delta }tR_m`$, which is not evolved by the algorithm; this field value is supplied along with the initial data. Once the Cauchy and characteristic boundary values are computed, a new iteration may be performed starting from Step 1 above.
Since all the interpolations employed in the matching step have fourth order error, the matching algorithm has the same second order global accuracy exhibited by the separate Cauchy and characteristic algorithms, as confirmed by numerical tests.
In summary, CCM has been implemented and tested for nonlinear waves in 3-dimensional space. No special assumption is made about the waves crossing the computational interface $`r=R_m`$ and nonlinear effects in the exterior characteristic domain are automatically taken into account. In numerical experiments CCM converged to the exact solution (with the matching boundary fixed at an arbitrary position) in highly nonlinear problems. For comparison, nonlocal ABC’s yielded convergent results only in linear problems. In terms of both computational cost and accuracy, CCM is a very effective way to solve the 3-dimensional wave equations, with or without a nonlinear term. It is possible to achieve convergence with an ABC by refining the grid and simultaneously enlarging the radius of the outer boundary. However, this is very expensive computationally, especially for small target error in the determination of the radiated signal in a 3-dimensional problems . Because CCM is convergent under grid-refinement alone, for small target error its performance is significantly better than any available alternative. In strongly nonlinear problems. CCM appears to be the only available method which is able to produce numerical solutions which converge to the exact solution with a computational interface located at an arbitrary fixed position.
## V Application to black holes
By definition, a black hole traces out a world tube in space-time which is the boundary of what is visible to an outside observer. As a result, this worldtube is necessarily a characteristic hypersurface traced out in space-time by light rays (the characteristics of general relativity).
A mechanical analogue of a black hole can in principle be constructed by letting a reservoir of water empty through a funnel into a lower reservoir. By arranging the flow velocity in the funnel to exceed the velocity of sound in the upper reservoir, an acoustic version of a black hole results . Sound waves can travel from the upper reservoir to the lower but not in the reverse direction.
The simplest example of a black hole arises in the spherically symmetric collapse of a star whose energy has been depleted to the extent that internal pressure cannot withstand gravity. A point sized black hole first forms at the center of the collapsing star, which then expands into a sphere growing with the velocity of light, thus tracing out a light cone in space-time. Normally a light cone keeps expanding forever. What makes a black hole light cone unique is that gravity halts this expansion and eventually the black hole just hovers in equilibrium at a fixed size. Each point on the spherical black hole still moves along a light ray, but the sphere of light rays is in a delicate gravitational balance between growing and shrinking.
This spherically symmetric black hole was discovered analytically by Schwarzschild as a simple solution to Einstein’s equations. Unfortunately, Schwarzschild’s black hole does not emit gravitational waves (just as a spherically symmetric charge distribution does not emit electromagnetic waves).
The inspiral and merger of a binary system of black holes is a powerful source of gravitational waves, which lie in the frequency band detectable by the new gravity wave observatories. The lack of symmetry necessitates a computational treatment to determine the waveform of the radiated signal. The binary black holes trace out a characteristic hypersurface in space-time so that their simulation by characteristic evolution is a natural approach. Using a characteristic evolution algorithm similar to that which I have described here, except now applied to Einstein’s equation, we have found fascinating results. The binary black holes, rather than initially forming as a point caustic in the Schwarzschild case, form as a cross-over surface where light rays collide. This cross-over surface is itself bounded by a ring of caustics which in a sense mark the merger of the two individual black holes into a single black hole.
The individual black holes form as spheres but, as illustrated in Fig. 5, as the holes approach, their mutual gravitational tidal distortion produces sharp pincers just prior to merger. At merger, the pincers join to form a single temporarily toroidal black hole, as illustrated in Fig. 6. The inner hole of the torus subsequently closes up to produce first a single peanut shaped black hole and finally a spherical black hole. Details of this merger can be viewed at the web site http://artemis.phyast.pitt.edu/animations. We are now using characteristic evolution to compute the gravitational wave signal emitted in such black hole collisions in the anticipation that they will be detected by the new gravity wave observatories.
## VI Acknowledgements
I thank the Pittsburgh Supercomputing Center and NPACI for making computing time available for this research and the National Science Foundation for research support under grant NSF PHY 9510895. |
warning/0003/gr-qc0003114.html | ar5iv | text | # Dynamic Cosmic Strings II: Numerical evolution of excited Cosmic Strings
## I Introduction
This is the second paper in a series devoted to the study of dynamic cosmic strings. In the previous paper , henceforth referred to as paper I, we derived the equations of motion for a time dependent cylindrically symmetric cosmic string coupled to a gravitational field with two degrees of freedom. The treatment involved using a Geroch decomposition to reformulate the problem in terms of matter fields on a $`2+1`$ dimensional spacetime and two geometrical variables $`\nu `$ and $`\tau `$ which describe the gravitational degrees of freedom. Unlike the original cylindrical metric the reduced $`2+1`$ spacetime is asymptotically flat and this allows us to implement a numerical treatment of the field equations which includes null infinity as part of the grid. In paper I we presented a Cauchy-characteristic matching (CCM) code that reproduced the results of several vacuum solutions with both one and two degrees of freedom but did not perform entirely satisfactorily in the presence of matter. We have therefore developed a second implicit, purely characteristic code. The numerical details, the testing and the convergence analysis of this implicit code will be described in this paper. We also give a detailed analysis of the interaction between the string and a pulse of gravitational radiation.
After briefly describing the variables used for the coupled system in section II the details of the numerical scheme are presented in section III. We begin by describing the relaxation method we use for the much simpler problem of a static cosmic string in Minkowski spacetime and demonstrate how this approach leads naturally to the implicit scheme used to solve the general case of a dynamical cosmic string coupled to gravity. The testing of the code is described in section IV. This involves comparing it with exact solutions, checking that static cosmic string initial data remain static when they are evolved, and undertaking a time dependent convergence analysis of the code. We are able to show that the code is accurate, stable and shows clear second order convergence. In section V we analyse the interaction between an initially static cosmic string and a Weber–Wheeler type pulse of gravitational radiation. We show that this interaction causes the string fields $`S`$ and $`P`$ to oscillate and examine how this oscillation depends upon both the width and amplitude of the pulse and also upon the constants $`e`$, $`\lambda `$ and $`\eta `$. The key result is that the frequencies are essentially independent of the nature of the Weber–Wheeler wave but are proportional to the masses $`m_S`$ and $`m_P`$ associated with the scalar field $`S`$ and the vector field $`P`$. The string continues to ring after the gravitational pulse has largely radiated away, but the oscillations slowly decay and the variables eventually return to their static values. It is only possible to observe this effect because of the long term stability of the code. Finally in section VI we discuss our results and outline future work.
## II The variables and field equations
We begin with a brief summary of the variables used to describe the spacetime and the string. In cylindrical polar coordinates $`(t,\rho ,\varphi ,z)`$ we may write the line element for a cylindrically symmetric spacetime with two degrees of freedom as a modified version of that given by Jordan, Ehlers, Kundt and Kompaneets ,
$$ds^2=e^{2(\gamma \psi )}(dt^2d\rho ^2)\rho ^2e^{2\psi }d\varphi ^2e^{2(\psi +\mu )}(\omega d\varphi +dz)^2,$$
(1)
where $`\psi `$, $`\omega `$, $`\mu `$ and $`\gamma `$ are functions of $`t`$, $`\rho `$. As shown in paper I, in order to work with an asymptotically flat spacetime it is convenient to make a Geroch decomposition in which the 4-dimensional metric is replaced by a 3-dimensional one and two auxiliary scalar fields. These scalar fields are the norm of the axial Killing vector $`\nu `$ and the Geroch potential $`\tau `$ and are related to $`\psi `$ and $`\omega `$ by equations (15) and (16) of paper I, while the conformal 3-dimensional line element is given by
$$d\stackrel{~}{\sigma }^2=e^{2(\gamma +\mu )}(dt^2d\rho ^2)\rho ^2e^{2\mu }d\varphi ^2.$$
(2)
The cosmic string is described by a complex scalar field $`\mathrm{\Phi }`$ coupled to a U(1) gauge field $`A_\mu `$, which in cylindrical symmetry can be written as
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}S(t,\rho )e^{i\varphi },`$ (3)
$`A_\mu `$ $`=`$ $`{\displaystyle \frac{1}{e}}[P(t,\rho )1]_\mu \varphi .`$ (4)
It is also helpful to introduce rescaled coupling constants and variables given by
$`X`$ $`=`$ $`{\displaystyle \frac{S}{\eta }},`$ (5)
$`\alpha `$ $`=`$ $`{\displaystyle \frac{e^2}{\lambda }},`$ (6)
$`r`$ $`=`$ $`\sqrt{\lambda }\eta \rho ,`$ (7)
$`\stackrel{~}{t}`$ $`=`$ $`\sqrt{\lambda }\eta t.`$ (8)
Here $`\eta `$ is the vacuum expectation value of the scalar field while $`\alpha `$ represents the relative strength of the coupling between the scalar and vector field given by $`e`$, compared to the self-coupling of the scalar field given by $`\lambda `$. Critical coupling, for which the masses of the scalar and vector fields are equal, is given by $`\alpha =8`$ .
The field equations for a cosmic string coupled to gravity were derived in paper I in terms of Cauchy coordinates \[Paper I (26)–(33)\] and also in terms of compactified characteristic coordinates \[Paper I (43)–(49)\]. However the fully characteristic numerical scheme also requires characteristic equations in the inner region. In terms of the retarded time $`u=\stackrel{~}{t}r`$ and the radius $`r`$ the field equations are given by
$`\mathrm{}\nu `$ $`=`$ $`\nu _{,r}\mu _{,r}+{\displaystyle \frac{\tau _{,r}^2\nu _{,r}^2}{\nu }}\nu _{,u}\mu _{,r}\nu _{,r}\mu _{,u}+2{\displaystyle \frac{\nu _{,u}\nu _{,r}\tau _{,u}\tau _{,r}}{\nu }}+8\pi \eta ^2\left[2e^{2(\gamma +\mu )}(X^21)^2+e^{2\mu }\nu ^2{\displaystyle \frac{2P_{,u}P_{,r}P_{,r}^2}{\alpha r^2}}\right],`$ (9)
$`\mathrm{}\tau `$ $`=`$ $`\tau _{,r}\mu _{,r}2{\displaystyle \frac{\tau _{,r}\nu _{,r}}{\nu }}\tau _{,u}\mu _{,r}\tau _{,r}\mu _{,u}+2{\displaystyle \frac{\tau _{,r}\nu _{,u}+\tau _{,u}\nu _{,r}}{\nu }},`$ (10)
$`\mathrm{}\mu `$ $`=`$ $`\mu _{,r}^2+{\displaystyle \frac{\mu _{,r}}{r}}{\displaystyle \frac{\mu _{,u}}{r}}2\mu _{,u}\mu _{,r}+8\pi \eta ^2\left[2{\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}(X^21)^2+e^{2\gamma }{\displaystyle \frac{X^2P^2}{r^2}}\right],`$ (11)
$`0`$ $`=`$ $`2\gamma _{,r}+2r\gamma _{,r}\mu _{,r}r\mu _{,rr}+r\mu _{,r}^2{\displaystyle \frac{r}{2\nu ^2}}(\tau _{,r}^2+\nu _{,r}^2)8\pi \eta ^2\left[rX_{,r}^2+{\displaystyle \frac{1}{\alpha }}e^{2\mu }\nu {\displaystyle \frac{P_{,r}^2}{r}}\right],`$ (12)
$`\mathrm{}P`$ $`=`$ $`2{\displaystyle \frac{P_{,r}}{r}}+2{\displaystyle \frac{P_{,u}}{r}}P_{,r}\mu _{,r}+P_{,r}{\displaystyle \frac{\nu _{,r}}{\nu }}+P_{,r}\mu _{,u}+P_{,u}\mu _{,r}{\displaystyle \frac{P_{,r}\nu _{,u}+P_{,u}\nu _{,r}}{\nu }}\alpha {\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}PX^2,`$ (13)
$`\mathrm{}X`$ $`=`$ $`X_{,r}\mu _{,r}X_{,u}\mu _{,r}X_{,r}\mu _{,u}4{\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}X(X^21)e^{2\gamma }{\displaystyle \frac{XP^2}{r^2}},`$ (14)
where $`\mathrm{}`$ represents the flat-space d’Alembert operator
$$\mathrm{}=2\frac{^2}{ur}\frac{^2}{r^2}\frac{1}{r}\left(\frac{}{r}\frac{}{u}\right).$$
(15)
Note that there are two further Einstein equations which are not used in the numerical scheme as they are a consequence of the above equations and their derivatives. These equations are only used to provide a check on the numerical accuracy of the code.
## III Numerical methods
In order to solve the above field equations we have developed two independent codes. The first is based on the Cauchy-characteristic matching code of Dubal et al. and d’Inverno et al. . This code performs well in the absence of matter and has been used in paper I to study several cylindrically symmetric vacuum solutions. In paper I we have also given details of the convergence analysis and described the modifications with respect to that lead to long term stability with both gravitational degrees of freedom present. However the CCM code performed less satisfactorily in the evolution of the cosmic string. This is due to the existence of unphysical solutions to the evolution equations (9)–(14) which diverge exponentially as $`r\mathrm{}`$. Controlling the time evolution near null infinity by means of a bump function enabled us to select the physical solutions with regular behavior at $`I^+`$, but the bump function itself introduced noise which eventually gave rise to instabilities. We therefore implemented a second implicit, purely characteristic, code which allowed us to directly apply boundary conditions at the origin as well as null infinity and thus suppress diverging solutions. It is interesting that this problem is already present in the calculation of the static cosmic string in Minkowski spacetime. We will, therefore, first describe the numerical scheme used in the static Minkowskian case where the equations are fairly simple. We then present the modifications necessary for the static and dynamic case coupled to the gravitational field.
### A The static cosmic string in Minkowski spacetime
In (9)–(14) we set the metric variables to their Minkowskian values and all time derivatives to zero to obtain the equations for the static cosmic string in Minkowski spacetime (cf. )
$`r{\displaystyle \frac{d}{dr}}\left(r^1{\displaystyle \frac{dP}{dr}}\right)`$ $`=`$ $`\alpha X^2P,`$ (16)
$`r{\displaystyle \frac{d}{dr}}\left(r{\displaystyle \frac{dX}{dr}}\right)`$ $`=`$ $`X\left[P^2+4r^2(X^21)\right].`$ (17)
The boundary conditions are
$`P(0)=1,`$ $`\underset{r\mathrm{}}{lim}P(r)=0,`$ (18)
$`X(0)=0,`$ $`\underset{r\mathrm{}}{lim}X(r)=1.`$ (19)
In order to cover the whole spacetime with a finite coordinate range, we divide the computational domain into two regions. In the inner region ($`0r1`$) we use the coordinate $`r`$, while in the outer region we introduce the compactified radius
$`y`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{r}}},`$ (20)
which covers the range $`1y0`$ corresponding to the region $`r1`$. It is also useful to combine $`r`$ and $`y`$ into a single radial variable $`w`$ defined by
$`w=\{r\mathrm{for}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}r132/\sqrt{r}\mathrm{for}r>1\text{,}`$ (21)
so that the region $`0w<3`$ corresponds to $`0r<\mathrm{}`$ and infinity is mapped to $`w=3`$ .
In terms of the variable $`y`$, (16) and (17) take the form
$`y{\displaystyle \frac{d}{dy}}\left(y^5{\displaystyle \frac{dP}{dy}}\right)`$ $`=`$ $`4\alpha X^2P,`$ (22)
$`y{\displaystyle \frac{d}{dy}}\left(y{\displaystyle \frac{dX}{dy}}\right)`$ $`=`$ $`4X\left[P^2+4{\displaystyle \frac{(X^21)}{y^4}}\right].`$ (23)
The number of grid points in each region may differ, but each half-grid is uniform. Thus we use a total of $`N:=n_1+n_2`$ grid points where the points labelled $`n_1`$ and $`n_1+1`$ both correspond to the position $`r=1=y`$. The points $`n_1`$, $`n_1+1`$ form the interface between the two regions (see Figure 1). One point will contain the variables in terms of $`r`$, the other in terms of $`y`$. With the computational grid covering the whole spacetime, we now face a two point boundary value problem. Due to the existence of unphysical solutions diverging at $`y=0`$ shooting methods turned out to be unsuitable for solving this problem. On the other hand numerical relaxation, as described in for example, allows us to directly control the behavior of $`P`$ and $`X`$ at infinity. The form of equations (16), (17) suggests that in order to write them as a first order system we should introduce the auxiliary variables $`Q=r^1P_{,r}`$ and $`R=rX_{,r}`$. The equations may then be written in the form
$`P_{,r}`$ $`=`$ $`rQ,`$ (24)
$`X_{,r}`$ $`=`$ $`{\displaystyle \frac{R}{r}},`$ (25)
$`Q_{,r}`$ $`=`$ $`\alpha {\displaystyle \frac{PX^2}{r}},`$ (26)
$`R_{,r}`$ $`=`$ $`X\left[{\displaystyle \frac{P^2}{r}}+4r(X^21)\right].`$ (27)
The corresponding equations in the outer region are given by
$`P_{,y}`$ $`=`$ $`2{\displaystyle \frac{Q}{y^5}},`$ (28)
$`X_{,y}`$ $`=`$ $`2{\displaystyle \frac{R}{y}},`$ (29)
$`Q_{,y}`$ $`=`$ $`2\alpha {\displaystyle \frac{X^2P}{y}},`$ (30)
$`R_{,y}`$ $`=`$ $`2X\left({\displaystyle \frac{P^2}{y}}+4{\displaystyle \frac{X^21}{y^5}}\right).`$ (31)
Standard second order centered finite differencing results in $`4(N2)`$ non-linear algebraic equations which are supplemented by the 4 boundary conditions (19) and 4 interface relations
$`P_{n_1+1}`$ $`=`$ $`P_{n_1},`$ (32)
$`X_{n_1+1}`$ $`=`$ $`X_{n_1},`$ (33)
$`Q_{n_1+1}`$ $`=`$ $`Q_{n_1},`$ (34)
$`R_{n_1+1}`$ $`=`$ $`R_{n_1}.`$ (35)
We then start with piecewise linear initial guesses for $`P`$ and $`X`$ (and the corresponding derivatives $`Q`$ and $`R`$) and solve the $`4N`$ algebraic equations iteratively with a Newton–Raphson method. Results for various choices of the string parameter $`\alpha `$ are shown in Figure 6 of paper I.
### B The static cosmic string coupled to gravity
From the numerical point of view, the problem of solving for a static cosmic string coupled to gravity through the Einstein equations is virtually identical to that of a static string in Minkowski spacetime. The only difference is the much higher degree of complexity of the equations due to the appearance of the functions $`\nu `$, $`\tau `$, $`\mu `$ and $`\gamma `$ as extra variables. We do not present the equations here, since they may be derived from the fully dynamic case by setting all time derivatives to zero in the relevant equations. The solution is again obtained using the relaxation method described in the previous section. As our initial guess for the metric variables we use Minkowskian values, and for the string variables $`X`$ and $`P`$ we use the previously calculated values for a Minkowskian string with the same string parameters. The results are shown in Figure 7 of paper I.
### C The dynamic code
In the dynamic case all variables $`\nu `$, $`\tau `$, $`\mu `$, $`\gamma `$, $`P`$ and $`X`$ are functions of $`u,r`$ and we have to solve the system (9)–(14) of partial differential equations (PDEs). In order to control the behavior of the solution at infinity, we need a generalisation for PDEs of the relaxation scheme applied to ordinary differential equations (ODEs). In view of the characteristic feature of the relaxation scheme, namely the simultaneous calculation of new function values at all grid points, this generalisation leads directly to implicit evolution schemes as used for hyperbolic or parabolic PDEs. Therefore, the dynamic code is based on the implicit, second order in space and time Crank–Nicholson scheme (see for example). For each solution step we consider two spatial slices of the grid-type of Figure 1, labelled $`n`$ and $`n+1`$. We then apply centered finite differencing according to
$`f`$ $`=`$ $`{\displaystyle \frac{f_{k+1}^{n+1}+f_k^{n+1}+f_{k+1}^n+f_k^n}{4}},`$ (36)
$`f_{,r}`$ $`=`$ $`{\displaystyle \frac{f_{k+1}^{n+1}f_k^{n+1}+f_{k+1}^nf_k^n}{2\mathrm{\Delta }r}},`$ (37)
$`f_{,u}`$ $`=`$ $`{\displaystyle \frac{f_{k+1}^{n+1}+f_k^{n+1}f_{k+1}^nf_k^n}{2\mathrm{\Delta }u}},`$ (38)
where $`f`$ represents any of our variables. Assuming that all functions are known on slice $`n`$ we arrive at a large set of algebraic equations for the $`f_i^{n+1}`$, similar to the static case, which needs to be supplemented by boundary conditions and interface relations analogous to (32)–(35). Again the system of algebraic equations is solved iteratively with the Newton–Raphson method. The initial guess for the data on the new slice $`n+1`$ is taken from the previous slice and convergence is typically achieved within three iterations. For this purpose we rewrite the dynamic equations (9)–(14) as a first order system. The equations for the variables $`\nu `$, $`\tau `$ and $`X`$ involve radial derivatives which may be written in terms of the second order operator $`\frac{}{r}(r\frac{}{r})`$ so we introduce the corresponding variables $`N=r\nu _{,r}`$, $`T=r\tau _{,r}`$ and $`R=rX_{,r}`$ \[cf. paper I equations (110)–(115)\]. The equation for $`\mu `$ on the other hand involves $`\frac{}{r}(r^2\frac{}{r})`$, while that for $`P`$ involves $`\frac{}{r}(r^1\frac{}{r})`$. We therefore introduce the corresponding variables $`M=r^2\mu _{,r}`$ and $`Q=r^1P_{,r}`$. Finally the equation for $`\gamma `$ involves only one $`r`$ derivative and may be written in first order form without having to introduce any further quantities. In terms of these variables (9)–(14) become:
$`\nu _{,r}`$ $`=`$ $`{\displaystyle \frac{N}{r}},`$ (39)
$`\tau _{,r}`$ $`=`$ $`{\displaystyle \frac{T}{r}},`$ (40)
$`\mu _{,r}`$ $`=`$ $`{\displaystyle \frac{M}{r^2}},`$ (41)
$`P_{,r}`$ $`=`$ $`rQ,`$ (42)
$`X_{,r}`$ $`=`$ $`{\displaystyle \frac{R}{r}},`$ (43)
$`2N_{,u}`$ $`=`$ $`N_{,r}+{\displaystyle \frac{T^2N^2}{r\nu }}+{\displaystyle \frac{NM}{r^2}}+2{\displaystyle \frac{\nu _{,u}N\tau _{,u}T}{\nu }}\nu _{,u}{\displaystyle \frac{\nu _{,u}M}{r}}N\mu _{,u}`$ (45)
$`+8\pi \eta ^2\left[2e^{2(\gamma +\mu )}r(X^21)^2+{\displaystyle \frac{1}{\alpha }}e^{2\mu }\nu ^2(2P_{,u}QrQ^2)\right],`$
$`2T_{,u}`$ $`=`$ $`T_{,r}2{\displaystyle \frac{TN}{r\nu }}+2{\displaystyle \frac{\tau _{,u}N+\nu _{,u}T}{\nu }}+{\displaystyle \frac{TM}{r^2}}\tau _{,u}{\displaystyle \frac{\tau _{,u}M}{r}}T\mu _{,u},`$ (46)
$`2M_{,u}`$ $`=`$ $`M_{,r}+{\displaystyle \frac{M^2}{r^2}}2\mu _{,u}M2r\mu _{,u}+8\pi \eta ^2\left[e^{2\gamma }X^2P^2+2{\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}r^2(X^21)^2\right],`$ (47)
$`2(r+M)\gamma _{,r}`$ $`=`$ $`M_{,r}2{\displaystyle \frac{M}{r}}{\displaystyle \frac{M^2}{r^2}}+{\displaystyle \frac{T^2+N^2}{2\nu ^2}}+8\pi \eta ^2\left[R^2+{\displaystyle \frac{1}{\alpha }}e^{2\mu }\nu r^2Q^2\right],`$ (48)
$`2Q_{,u}`$ $`=`$ $`Q_{,r}{\displaystyle \frac{QM}{r^2}}+Q\mu _{,u}{\displaystyle \frac{Q\nu _{,u}}{\nu }}{\displaystyle \frac{P_{,u}N}{r^2\nu }}+{\displaystyle \frac{QN}{r\nu }}+{\displaystyle \frac{P_{,u}}{r^2}}+{\displaystyle \frac{P_{,u}M}{r^3}}\alpha {\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}{\displaystyle \frac{PX^2}{r}},`$ (49)
$`2R_{,u}`$ $`=`$ $`R_{,r}X_{,u}{\displaystyle \frac{X_{,u}M}{r}}+{\displaystyle \frac{RM}{r^2}}R\mu _{,u}4{\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}rX(X^21)e^{2\gamma }{\displaystyle \frac{XP^2}{r}}.`$ (50)
The corresponding first order system in the outer region is given by
$`\nu _{,y}`$ $`=`$ $`2{\displaystyle \frac{N}{y}},`$ (51)
$`\tau _{,y}`$ $`=`$ $`2{\displaystyle \frac{T}{y}},`$ (52)
$`\mu _{,y}`$ $`=`$ $`2yM,`$ (53)
$`P_{,y}`$ $`=`$ $`2{\displaystyle \frac{Q}{y^5}},`$ (54)
$`X_{,y}`$ $`=`$ $`2{\displaystyle \frac{R}{y}},`$ (55)
$`2N_{,u}`$ $`=`$ $`{\displaystyle \frac{1}{2}}y^3\left(N_{,y}2yNM2{\displaystyle \frac{T^2N^2}{y\nu }}\right)y^2\nu _{,u}MN\mu _{,u}+2{\displaystyle \frac{\nu _{,u}N\tau _{,u}T}{\nu }}\nu _{,u}`$ (57)
$`+8\pi \eta ^2\left[2e^{2(\gamma +\mu )}{\displaystyle \frac{(X^21)^2}{y^2}}+{\displaystyle \frac{1}{\alpha }}e^{2\mu }\nu ^2\left(2P_{,u}Q{\displaystyle \frac{Q^2}{y^2}}\right)\right],`$
$`2T_{,u}`$ $`=`$ $`{\displaystyle \frac{1}{2}}y^3\left(T_{,y}2yTM+4{\displaystyle \frac{TN}{y\nu }}\right)T\mu _{,u}y^2\tau _{,u}M+2{\displaystyle \frac{\tau _{,u}N+\nu _{,u}T}{\nu }}\tau _{,u},`$ (58)
$`2M_{,u}`$ $`=`$ $`{\displaystyle \frac{1}{2}}y^3(M_{,y}2yM^2)2{\displaystyle \frac{\mu _{,u}}{y^2}}2\mu _{,u}M+8\pi \eta ^2\left[e^{2\gamma }X^2P^2+2{\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}{\displaystyle \frac{(X^21)^2}{y^4}}\right],`$ (59)
$`2(y^2M+1)\gamma _{,y}`$ $`=`$ $`y^2M_{,y}+4yM+2y^3M^2{\displaystyle \frac{N^2+T^2}{y\nu ^2}}16\pi \eta ^2\left({\displaystyle \frac{R^2}{y}}+{\displaystyle \frac{1}{\alpha }}e^{2\mu }\nu {\displaystyle \frac{Q^2}{y^5}}\right),`$ (60)
$`2Q_{,u}`$ $`=`$ $`{\displaystyle \frac{1}{2}}y^3\left(Q_{,y}+2yQM2{\displaystyle \frac{QN}{y\nu }}\right)+y^4P_{,u}y^4{\displaystyle \frac{P_{,u}N}{\nu }}{\displaystyle \frac{\nu _{,u}Q}{\nu }}`$ (62)
$`+y^6P_{,u}M+Q\mu _{,u}\alpha {\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}y^2PX^2,`$
$`2R_{,u}`$ $`=`$ $`{\displaystyle \frac{1}{2}}y^3(R_{,y}2yRM)X_{,u}R\mu _{,u}y^2X_{,u}Me^{2\gamma }y^2XP^24{\displaystyle \frac{e^{2(\gamma +\mu )}}{\nu }}{\displaystyle \frac{X(X^21)}{y^2}}.`$ (63)
In order to solve these equations we must supplement them by appropriate initial and boundary conditions. We start by considering boundary conditions on the axis. In general we find the code is more stable if one imposes boundary conditions on the radial derivatives rather than the variables themselves. For the variables $`\nu `$, $`\tau `$ and $`X`$ we therefore impose the required boundary conditions on the initial data, but in the subsequent evolution we impose the weaker condition that their radial derivatives are finite on the axis. This ensures that the evolution equations propagate the axial conditions given on the initial data. For the variable $`\mu `$ we impose the condition that $`M`$ is zero on the axis which is equivalent to the rather weak condition that $`r^2\mu _r`$ vanishes there. The inverse power of $`r`$ in the definition of $`Q`$ makes it unsuitable to specify the value of this quantity at $`r=0`$ so in this case we work with the variable directly and require that $`P=1`$ on the axis. Finally the variable $`\gamma `$ is given by a purely radial equation, so in this case we must specify the value on the axis which is chosen to be zero to ensure elementary flatness. Therefore at $`r=0`$ we require
$`N`$ $`=`$ $`0,`$ (64)
$`T`$ $`=`$ $`0,`$ (65)
$`M`$ $`=`$ $`0,`$ (66)
$`\gamma `$ $`=`$ $`0,`$ (67)
$`P`$ $`=`$ $`1,`$ (68)
$`R`$ $`=`$ $`0.`$ (69)
For the boundary conditions at null infinity we know that regular solutions of the cylindrical wave equation have radial derivatives that decay faster than $`1/r`$ so that we may take the variables $`N`$, $`T`$ and $`R`$, which satisfy a wave type equation, to vanish at $`y=0`$. The asymptotics of $`\mu `$ are slightly different due to the additional power of $`r`$ in the radial derivative (similar to the spherically symmetric wave equation) but for a regular solution $`\mu _{,y}`$ vanishes at null infinity. The $`P`$ equation does not satisfy a wave type equation due to the inverse power of $`r`$ but has asymptotic behavior given by a modified Bessel function. The physically relevant finite solution has exponential decay so in this case one may impose the condition that $`Q=0`$ at $`y=0`$. Hence we require the solution to satisfy the following boundary conditions at $`y=0`$
$`N`$ $`=`$ $`0,`$ (70)
$`T`$ $`=`$ $`0,`$ (71)
$`\mu _{,y}`$ $`=`$ $`0,`$ (72)
$`Q`$ $`=`$ $`0,`$ (73)
$`R`$ $`=`$ $`0.`$ (74)
These boundary conditions are sufficient to determine the solution of the first order system (39)–(63) while suppressing the unphysical solutions which are singular on the axis or null infinity. Note that $`\gamma `$ is determined by the constraint equation (12), which is a first order ODE, and thus only needs one boundary condition.
## IV Testing the code
In paper I we used the CCM code to reproduce several exact vacuum solutions to an accuracy of about $`10^5`$ with second order convergence. We have also checked the properties of the codes for the static cosmic string in Minkowski spacetime and coupled to gravity. Both clearly showed second order convergence. Here we will focus on testing the implicit dynamic code. We have carried out four different independent tests, namely
* Reproducing the non-rotating vacuum solution of Weber and Wheeler ,
* Reproducing the rotating vacuum solution of Xanthopoulos ,
* Using the results for the static cosmic string (paper I) as initial data and checking that the system stays in its static configuration,
* Convergence analysis for the string hit by a Weber–Wheeler wave.
Two additional tests arise in a natural way from the field equations and the numerical scheme. As described above there are two additional field equations which are algebraic consequences of the other field equations. We have verified that these equations are satisfied to second order accuracy ($`\mathrm{\Delta }r^2`$). Furthermore the numerical scheme calculates the residuals of the algebraic equations to be solved, which have thus been monitored in test runs. They are satisfied to a much higher accuracy (double precision machine accuracy), so the total error is dominated by the truncation error of the second order differencing scheme. Another independent test is the comparison with the explicit CCM code which yields good agreement for as long as the latter remains stable. The four main tests are now described in more detail.
### A The Weber–Wheeler wave
In the first test we evolve the analytic solution given by Weber and Wheeler , which describes a gravitational pulse of the + polarisation mode. This solution has two free parameters, $`a`$ and $`b`$, which can be interpreted as the width and amplitude of the pulse. The equations together with a more detailed discussion have been given in paper I.
We prescribe $`\nu `$ as initial data according to the analytic expressions obtained for $`a=2`$ and $`b=0.5`$ and set the other free variables to zero, while $`\gamma `$ is calculated via quadrature from the constraint equation (12). In Figure 2 we show the deviation of the numerical results from the analytical one for $`N=1920`$ grid points (320 points in the inner, 1600 points in the outer region) and a Courant factor of 0.5 with respect to the inner region. The convergence analysis (see below) shows that this number of points provides sufficient resolution in the outer region while still keeping computation times at a tolerable level. All computations presented in this work have been obtained with these grid parameters, unless stated otherwise. The code stays stable for much longer time intervals than shown in the figure, but does not reveal any further interesting features as the analytic solution approaches its Minkowskian values and the error goes to zero.
### B The rotating solution of Xanthopoulos
Xanthopoulos derived an analytic vacuum solution for Einstein’s field equations in cylindrical symmetry containing both the + and $`\times `$ polarisation mode. Its analytic form in terms of our metric variables and a more detailed discussion has been given in paper I. The solution has one free parameter $`a`$ which is set to one in this calculation.
The error of our numerical results is displayed in Figure 3, where we have used the same grid parameters as in the Weber–Wheeler case. Again we have run the code for longer times and found that the error approaches zero. We conclude that the code reproduces both analytic vacuum solutions with excellent accuracy comparable to that of the CCM code and exhibits long term stability.
### C Evolution of the static cosmic string
The tests described above only involve vacuum solutions, so the matter part of the code and the interaction between matter and geometry has not been tested. An obvious test involving matter and geometry is to use the result for the static cosmic string in curved spacetime as initial data and evolve this scenario. All variables should, of course, remain at their initial values. We have evolved the static string data for our standard grid and the parameter set, $`\alpha =1`$ and $`\eta =0.2`$, which corresponds to a strong back-reaction of the string on the metric. The results are shown in Figure 4. The system stays in its static configuration with high accuracy over a long time interval.
### D Convergence analysis
Our investigation of the interaction between the cosmic string and gravitational waves will focus on the string being hit by a wave of the Weber–Wheeler type. In order to check this scenario for convergence we have run the code for the parameter set $`\eta =0.2`$, $`\alpha =1`$, $`a=2`$, $`b=0.5`$ for different grid resolutions, where $`a`$ and $`b`$ are again the width and amplitude of the Weber–Wheeler wave. In our case it is of particular interest to investigate the time dependence of the convergence to see what resolution we need in order to obtain reliable results for long runs. We calculate the convergence rate in the same way as in the static case (cf. paper I), but this time the $`\mathrm{}_2`$-norm is a function of time. So for each variable we have
$`\mathrm{}_2[\mathrm{\Delta }\mathrm{\Psi }^N](u)`$ $`=`$ $`\sqrt{{\displaystyle \frac{[\mathrm{\Psi }_k^N(u)\mathrm{\Psi }_k^{4320}(u)]^2}{N}}},`$ (75)
where the upper label “4320” indicates that the high resolution reference solution has been calculated for $`N=4320`$ grid points. In Figure 5 we show the convergence factor $`\mathrm{}_2[\mathrm{\Psi }^{1920}]/\mathrm{}_2[\mathrm{\Psi }^{2880}]`$ as a function of $`u`$ for $`\nu `$, $`\mu `$, $`\gamma `$, $`P`$ and $`X`$. The initial data for $`\tau `$ is identically zero for this scenario and stays zero during the evolution. The number of grid points is increased by a factor of 1.5 here (instead of the more commonly used 2) to reduce the computation time. Only points common to all grids have been used in the sum in equation (75). For second order convergence we would expect a convergence factor of $`1.5^2`$. Although the results in Figure 5 show weak variations with $`u`$, second order convergence is clearly maintained for long runs. In each case the outer region contains 5 times as many grid points as the inner region (e.g. $`n_1=320`$, $`n_2=1600`$ for the $`N=1920`$ case). The reason for this is that in the dynamical evolutions $`X`$ and especially $`P`$ exhibit significant spatial variations out to large radii. Due to the compactification, the spatial resolution of our grid decreases as we move towards null infinity and to resolve the spatial variations of the string variables out to sufficiently large radii we therefore have to introduce a large number of grid points in the outer region. No such problems occur in the inner region. If significantly fewer grid points are used in the outer region for this analysis, the convergence properties of the string variables can deteriorate to roughly first order level.
## V Time dependence of the string variables
### A Static cosmic strings excited by gravitational waves
The scenario we are going to investigate in this section is an initially static cosmic string hit by a gravitational wave of Weber–Wheeler type. For this purpose we use the static results with two modifications as initial data. Firstly the static metric function $`\nu _0`$ is multiplied by the exact Weber–Wheeler solution to simulate the gravitational wave pulse. Thus we guarantee that initially the cosmic string is indeed in an equilibrium configuration provided the wave pulse is located sufficiently far away from the origin and its interaction with the string is negligible. Ideally the numerical calculation would start with the incoming wave located at past null infinity. In order to approximate this scenario, we found it was sufficient to use the large negative initial time $`u_0=20`$. The second modification is to calculate $`\gamma `$ from the constraint equation (12) to preserve consistency with the Einstein field equations. In Figure 6 the corresponding initial data for $`\nu `$, $`P`$ and $`X`$ are shown for the parameter set $`\eta =10^3`$, $`\alpha =1`$, $`a=2`$ and $`b=0.5`$. From now on we will refer to these values as “standard parameters” and only specify parameters if they take on non-standard values. Note that $`\tau `$ vanishes on the initial slice in this case and stays identically zero throughout the evolution. The case of rotating gravitational waves hitting a cosmic string will be analysed in a future publication. The time evolution of the “standard configuration” is shown in Figure 7 where we plot $`\nu `$, $`\mu `$, $`\gamma `$, $`P`$ and $`X`$ as functions of $`w`$ at times $`20`$, $`10`$, $`0`$, $`2`$ and $`10`$. While the wave pulse is still far away from the origin, its interaction with the cosmic string is negligible (dotted lines). When it reaches the core region, however, it excites the cosmic string and the scalar and vector field start oscillating (dashed curves). After being reflected at the origin, the wave pulse travels along the outgoing characteristics and the metric variables $`\nu `$, $`\mu `$ and $`\gamma `$ quickly settle down into their static configuration which is close to Minkowskian values for $`\eta =10^3`$. The vector and scalar field of the cosmic string, on the other hand, continue ringing albeit with a different character. Whereas the oscillations of the scalar field $`X`$ are dominant in the range $`r2`$ and have significantly decayed at $`u=10`$ as shown in the figure, the vector field oscillations propagate to large radii and fall off very slowly (solid curves). This behavior is also illustrated in the right panel of Figure 8 which shows a contour plot of $`P`$ as a function of $`(u,r)`$ out to $`r=50`$. We shall see below, that the oscillations of $`P`$ will also decay eventually and the cosmic string will asymptotically settle back into its equilibrium configuration.
### B Frequency analysis
We will now quantitatively analyse the oscillations of the scalar and vector field of the cosmic string. Since we are working in rescaled coordinates, time and distance are measured in units of $`1/\sqrt{\lambda }\eta `$ and frequency in its inverse. To avoid complicated notation, however, we will omit the units from now on unless the meaning is unclear. In order to measure frequencies, we Fourier analyse the time evolution of the scalar and vector field for fixed radius $`r`$. Figure 9 shows $`P`$ and $`X`$ for standard parameters as functions of time at $`r=1`$ together with the corresponding power spectra. The plots for $`X`$ show a characteristic frequency $`f_X=0.43`$ whereas for $`P`$ we find a strong peak at $`f_P=0.16`$. The spectrum for $`P`$, however, also shows a strong mode at $`f=0.43`$. We have calculated similar power spectra for a large class of parameter sets and discovered this effect on numerous occasions — in addition to a strong peak at the characteristic frequency of $`P`$ or $`X`$ there is a second maximum at the characteristic frequency of the other field. In general the characteristic mode of $`X`$ resulted in stronger peaks at smaller radius, that of $`P`$ was stronger at larger radii. We attribute this feature to the interaction between the scalar and vector component of the cosmic string. The variation of the relative strength of the oscillations with radius confirms the corresponding observation in Figure 7. The accuracy of the measurements of the frequencies is limited by the resolution of the Fourier spectra which again is limited by the time interval covered in the evolution and, thus, by computation time. For tolerable computation times we get an accuracy $`\mathrm{\Delta }f0.01`$ which corresponds approximately to one bin in the frequency spectra.
In order to investigate the dependency of the oscillations on $`\alpha `$, $`\eta `$, $`a`$, $`b`$ and the radial position $`r`$, we have varied each parameter over at least two orders of magnitude while keeping the other parameters at standard values. We have found the following dependencies:
* The frequencies of both $`X`$ and $`P`$ did not show any variations with $`\eta `$ for $`\eta <0.1`$. (Note that $`\eta `$ does, however, appear in the units). For larger values of $`\eta `$, the non-linear interaction between string and geometry becomes dominant and we did not find a simple relation between frequency maxima and parameters.
* The variation of the parameters $`a`$ and $`b`$, the width and amplitude of the Weber–Wheeler pulse, has no measurable effect on the frequencies of $`X`$ and $`P`$, but only determined the amplitude of the oscillations. A narrow, strong pulse leads to larger amplitudes.
* For small $`r`$ the oscillations in $`X`$ are stronger, whereas those for $`P`$ dominate at large $`r`$. The frequency values, however, do not depend on the radius. For radii greater than 10 the oscillations in $`X`$ had decayed so strongly that we could no longer measure its frequency.
* To first order approximation the frequency of the scalar field, $`f_X`$, is independent of $`\alpha `$ over the observed range $`0.1\alpha 10`$. $`f_P`$, however, shows a strong dependence on $`\alpha `$ which is also illustrated in Figure 8, where we compare contour plots of $`P`$ obtained for $`\alpha =0.2`$ and $`1`$. The frequency is significantly larger for $`\alpha =1`$.
In Figure 10 we show $`f_P`$ and $`f_X`$ as functions of $`\alpha `$ together with the power laws obtained from a linear regression of the corresponding double logarithmic data. In the range $`5\alpha 8`$ (where $`\alpha =8`$ corresponds to the critical coupling), the expected values of $`f_P`$ and $`f_X`$ become similar and we observed only one maximum in the corresponding power spectra. Therefore a classification with respect to the vector or scalar origin of the frequencies is not obvious. These values (shown by filled lozenges in Figure 10) have not been used in the regression analysis. We obtain power law indices $`\sigma _X=0.00`$ and $`\sigma _P=0.50`$, so that
$`f_X`$ $``$ $`\mathrm{const}.`$ (76)
$`f_P`$ $``$ $`\sqrt{\alpha }.`$ (77)
Whereas the fitted curve for $`f_P`$ coincides with the observed values to high accuracy, Figure 10 shows a small but significant deviation of the measured $`f_X`$ from the fitted constant line. Indeed the sizable range over $`\alpha `$ for which we observe only one frequency indicates significant non-linear interaction similar to the effect of phase locking in non-linear systems of ODEs . If we consider $`f_X`$ and $`f_P`$ to be given by (76) and (77) in terms of the rescaled unphysical coordinates and we transform this back into physical units using $`\alpha =e^2/\lambda `$, we arrive at the following relations for the physical variables
$`f_X`$ $``$ $`\sqrt{\lambda }\eta ,`$ (78)
$`f_P`$ $``$ $`e\eta .`$ (79)
As shown in up to constant factors $`\sqrt{\lambda }\eta `$ and $`e\eta `$ are the masses of the scalar and the vector field, $`m_X`$ and $`m_P`$, and we conclude that $`X`$ and $`P`$ have characteristic frequencies
$`f_X`$ $``$ $`m_X,`$ (80)
$`f_P`$ $``$ $`m_P.`$ (81)
Since the frequencies for $`X`$ and $`P`$ seem only to depend upon the respective masses we have attempted to confirm these results by considering the oscillations of a cosmic string in two further scenarios. Firstly since the frequencies do not depend upon the Weber–Wheeler pulse we take as initial data the static values for the metric variables but excite the string by adding a Gaussian perturbation to either the $`X`$ or $`P`$ static initial values. The evolution is then computed using the fully coupled system. Secondly since the frequencies do not seem to depend upon the strength of the coupling to the gravitational field we have completely decoupled the gravitational field and considered the evolution of a cosmic string in Minkowski spacetime. The initial data is taken to be that for a static string in Minkowski spacetime with a Gaussian perturbation to either the $`X`$ or $`P`$ values. The evolution is then computed using the equations for a dynamical string in a Minkowskian background \[equations (85) and (86) of paper I\]. In both cases we find the same frequencies, to within an amount $`\mathrm{\Delta }f=0.01`$, that we found in the original case of the fully coupled system excited by a Weber–Wheeler pulse. Furthermore the frequencies did not depend on the location or shape of the field perturbation nor upon the choice of $`X`$ or $`P`$ as the perturbed field.
### C The long term behavior of the dynamic string
The time evolution shown in Figures 7 and 8 indicate that the oscillations of the cosmic string excited by gravitational waves gradually decay and metric and string settle down into an equilibrium state. We have calculated a very long run ($`20u410`$) to investigate the long term behavior in detail. The unphysically large value of $`\eta =0.1`$ is chosen for this calculation in order to guarantee a strong interaction between spacetime geometry and the cosmic string. In Figure 11 we show the difference $`\mathrm{\Delta }f:=f_{\mathrm{evol}}f_{\mathrm{stat}}`$ between the time-dependent $`\nu `$, $`\mu `$, $`\gamma `$ and $`X`$ and their corresponding static results obtained for the same parameters. For the vector field $`P`$ a similar 3-dimensional plot would require an extreme resolution to properly display the oscillations of the vector field (cf. Figure 7). For this reason we proceed differently and calculate the $`\mathrm{}_2`$-norm and the maximum of $`\mathrm{\Delta }P`$ for each slice $`u=\mathrm{const}`$. Both functions are plotted versus time in Figure 11. The incoming wave pulse can clearly be seen as a strong deviation of $`\nu `$ from the static function. The pulse excites the cosmic string and is reflected at the origin at $`u=0`$. The metric variables and the scalar field $`X`$ then quickly reach their equilibrium values. The oscillations in $`P`$ decay on a significantly longer time scale which is also evident in Figures 7 and 8 and the $`\mathrm{}_2`$-norm of $`\mathrm{\Delta }P`$ slowly approaches zero. Significantly longer runs than shown here are prohibited by the required computational time, but the results indicate that $`P`$ will also eventually reach its equilibrium configuration.
## VI Conclusion
In this paper we have described the details of the implicit, fully characteristic, numerical scheme which is used to solve the field equations for a cosmic string coupled to gravity. A feature of the cosmic string equations is that they admit exponentially diverging unphysical solutions. By using a Geroch decomposition it is possible to reformulate the problem in terms of fields which describe the string on an asymptotically flat $`2+1`$-dimensional metric as well as two auxiliary fields $`\nu `$ and $`\tau `$ which describe the gravitational degrees of freedom. We can then introduce a conformally compactified radial coordinate $`y`$ which allows us to include null infinity as part of the numerical grid. As well as avoiding the need to introduce artificial outgoing boundary conditions at the edge of the grid this approach has the advantage that we can also enforce boundary conditions for the string variables at null infinity which rule out the unphysical solutions. The use of the geometrically defined Geroch variables also improves the long term stability of the code compared to the use of metric variables.
The code has been shown to reproduce the results of two exact vacuum solutions, the Weber–Wheeler solution which describes a pulse of gravitational radiation with just the $`+`$ polarisation state, and a solution due to Xanthopoulos which describes a gravitational wave with both the $`+`$ and $`\times `$ polarisation states. The code has also been shown to reproduce the results of the static cosmic string code in that initial data corresponding to a static solution do not change when evolved in time using the dynamical code. For both the exact vacuum solutions and the static initial data the code shows excellent long term stability. Finally a time dependent convergence analysis demonstrates clear second order convergence of the code.
After demonstrating the reliability of the code we use it to analyse the interaction between an initially static cosmic string and a Weber–Wheeler type pulse of gravitational radiation. We find that the gravitational wave excites the string and causes the string variables $`X`$ and $`P`$ to oscillate before the configuration slowly settles back into its equilibrium state. In terms of the unphysical rescaled variables we find that the frequencies of the oscillations are essentially independent of the value of the coupling constant $`\eta `$ and of the width and amplitude of the Weber–Wheeler pulse. We also find that the frequency of $`X`$ is independent of $`\alpha `$ while that of $`P`$ is proportional to $`\sqrt{\alpha }`$. When this result is translated back into the physical units we find that the frequency of the scalar field is proportional to the mass of the scalar field and the frequency of the vector field is proportional to the mass of the vector field. This result is confirmed by investigating two further scenarios. Firstly we consider the evolution of static initial data for the string coupled to the gravitational field, but with a Gaussian perturbation to one of the string variables, and secondly we consider the same thing but in a Minkowskian background with the gravitational field decoupled. In both cases we obtain the same relationship between the frequencies and the mass.
Having investigated the interaction between a Weber–Wheeler type pulse of gravitational radiation and the cosmic string, the next obvious step is to consider the interaction between the string and a pulse of gravitational radiation with both polarisation states present. The code will also be a valuable tool in comparing the numerical results with those that one gets from using perturbation theory or the thin string limit to approximate the behavior of a cosmic string coupled to gravity.
###### Acknowledgements.
We would like to thank Ray d’Inverno for helpful discussions and Denis Pollney for help with GRTensor II. |
warning/0003/cond-mat0003418.html | ar5iv | text | # Current and power spectrum in a magnetic tunnel device with an atomic size spacer
## I Introduction
Recent interest in single-electron tunneling in ferromagnetic double tunnel junctions is stimulated by expected potential applications at microelectronics and by new phenomena observed in such systems. In order to have a device operating at room temperature the single electron charging energy $`E_c=e^2/2C`$ should be much larger than the thermal energy $`k_BT`$. It can be achieved decreasing the capacitance $`C`$ of the metallic spacer, which is proportional to its size. In a small metallic spacer a discreteness of the energy spectrum can be relevant and a separation of energy levels $`\mathrm{\Delta }Ek_BT`$. Such the situation was studied numerically just recently.
In the present paper we would like to investigate sequential tunneling in an extreme case, when the spacer particle has only a single electron level available for the tunneling process. This simplified model gives us possibility to gain a better insight into spin dependent tunneling processes and to solve the problem analytically. We will show that Coulomb interactions between electrons with different spins can lead to new effects. In some circumstances due to the Coulomb blockade effect the device can operate as a diode, in others it can show the negative differential resistance (NDR). The power spectrum analysis will be performed to understand correlations between currents for electrons of different spins and the transition from the sub-Poissonian to the super-Poissonian current noise in the ferromagnetic device.
## II Model and general derivations
Let us specify the system considered in detail. The separation between the ferromagnetic metallic electrodes is large and therefore, there is no direct electron tunneling between them. The electronic transport can be only via electronic states of the spacer particle placed between the electrodes. The particle can be a molecule (e.g. $`C_{60}`$), or a semiconductor quantum dot, in which the relevant energies are $`\mathrm{\Delta }E,E_ck_BT`$. For a small applied voltage $`V`$ ($`eV\mathrm{\Delta }E,E_c`$) electronic transport is only through a single electronic level $`E_0`$. Such the model was considered for a nonmagnetic device in Ref.\[\] and we generalize it for a ferromagnetic case including tunneling channels for electrons with opposite spin directions. The tunneling process for an electron with spin $`\sigma `$ through the left ($`j=1`$) and the right ($`j=2`$) junction is described by the net tunneling rates $`\gamma _{j\sigma }`$, which are assumed to be small $`\mathrm{}\gamma _{j\sigma }k_BT`$. This relation implies that the corresponding tunnel resistances $`R_{j\sigma }`$ are much larger than the quantum resistance $`R_Q=h/2e^2`$ and electronic transport can be described within the sequential tunneling approach. Since $`\mathrm{\Delta }E`$ is large, the tunneling process can be considered elastic (there is no thermalization of electrons on the spacer particle, which was usually assumed in the single electron transistor with a large metallic grain) . We also neglect fluctuations of the position of the electronic level $`E_0`$, which can be caused by thermal and electrostatic fluctuations of the environment.
Our model seems to be familiar to that considered recently for the Kondo effect in quantum dots. A condition for a development of the Kondo resonance is a buildup of many-body correlations between the dot and the electrodes, which can be achieved when electronic waves are coherently scattered on a magnetic impurity. It is in contrast to the present situation, where coupling between the particle and the electrodes is weak and electron tunneling events are uncorrelated and incoherent.
### A Stationary currents
Electronic transport is governed by the master equation
$`{\displaystyle \frac{d}{dt}}\left[\begin{array}{c}p_{}\\ p_{}\\ p_0\end{array}\right]=\widehat{M}\left[\begin{array}{c}p_{}\\ p_{}\\ p_0\end{array}\right],`$ (7)
where $`p_{}`$ and $`p_{}`$ denotes the probability to find an electron with the spin $`\sigma =`$ and $``$, $`p_0`$ \- the probability for an empty state $`E_0`$. Of course, the total probability $`p_{}+p_{}+p_0=1`$. The matrix $`\widehat{M}`$ is given by
$`\widehat{M}=\left[\begin{array}{ccc}\mathrm{\Gamma }_{}^{}& 0& \mathrm{\Gamma }_{}^+\\ 0& \mathrm{\Gamma }_{}^{}& \mathrm{\Gamma }_{}^+\\ \mathrm{\Gamma }_{}^{}& \mathrm{\Gamma }_{}^{}& \mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+\end{array}\right],`$ (11)
where $`\mathrm{\Gamma }_\sigma ^\pm =\mathrm{\Gamma }_{1\sigma }^\pm +\mathrm{\Gamma }_{2\sigma }^\pm `$, $`\mathrm{\Gamma }_{j\sigma }^\pm =f_j^\pm \gamma _{j\sigma }`$ are the total tunneling rates to ($`+`$) and off ($``$) the particle level $`E_0`$, $`f_j^\pm =\{1+\mathrm{exp}[\pm (E_0E_F(1)^jeV_j)/k_BT]\}^1`$. The voltage $`V`$ is applied to the left electrode and the voltage drop across the left and the right junction is $`V_1=C_2V/C`$ and $`V_2=C_1V/C`$, respectively. Here, $`C_j`$ denotes the capacitance of the $`j`$-th tunnel junction and $`C=C_1+C_2`$.
At the stationary state the probability $`p_\sigma `$ and $`p_0`$ are determined from the master equation (7) with the left hand side equal to zero, and the result is
$$p_\sigma =\frac{\mathrm{\Gamma }_\sigma ^+\mathrm{\Gamma }_\sigma ^{}}{\gamma _{}\gamma _{}\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+},p_0=\frac{\mathrm{\Gamma }_{}^{}\mathrm{\Gamma }_{}^{}}{\gamma _{}\gamma _{}\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+},$$
(12)
where $`\gamma _\sigma =\gamma _{1\sigma }+\gamma _{2\sigma }`$. The current through the left junction for electrons with the spin $`\sigma `$ is the difference of the tunneling current flowing to (+) and from ($``$) the particle
$`I_{1\sigma }I_{1\sigma }^+I_{1\sigma }^{}=e\left[\mathrm{\Gamma }_{1\sigma }^+p_0\mathrm{\Gamma }_{1\sigma }^{}p_\sigma \right]`$ (13)
$`=e(f_1^+f_2^+){\displaystyle \frac{\gamma _{1\sigma }^+\gamma _{2\sigma }^+\mathrm{\Gamma }_\sigma ^{}}{\gamma _{}\gamma _{}\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+}}.`$ (14)
Since there are no electronic relaxation processes on the particle, it results from the current conservation rule that $`I_{1\sigma }=I_{2\sigma }`$ for each electronic channel.
In magnetic tunnel junctions the resistance depends on the relative configuration of magnetic moments in the electrodes and this effect is known as the tunnel magnetoresistance (TMR). The value of TMR is given by the ratio $`TMR=(I_PI_{AP})/I_{AP}`$, where $`I_P`$ and $`I_{AP}`$ are the tunneling currents in the parallel (P) and the antiparallel (AP) configuration of the magnetic moments in the electrodes. It is convenient to express the tunneling rate coefficients in the form $`\gamma _{1\sigma }=\gamma _0(1\pm P_1)`$ and $`\gamma _{2\sigma }=\gamma _0\alpha (1\pm P_2)`$, where the sign $`+`$ ($``$) corresponds to the spin $`\sigma =`$ ($``$), $`P_1`$ and $`P_2`$ is the magnetic polarization of the left and the right electrode, $`\alpha `$ denotes the asymmetry between the potential barriers. Using Eq.(13) one gets
$$TMR=\frac{(1f_1^+f_2^+)4\alpha P_1P_2}{(1+\alpha )^2(P_1+\alpha P_2)^2(f_1^++\alpha f_2^+)^2+(P_1f_1^++\alpha f_2^+P_2)^2}.$$
(15)
For comparison we present the results for noninteracting electrons, i.e when the single electron charging energy $`E_c=0`$. In this limit the double occupancy of the level $`E_0`$ is allowed. The current through the left junctions for electrons with the spin $`\sigma `$ is then
$$I_{1\sigma }^0=e(f_1^+f_2^+)\frac{\gamma _{1\sigma }\gamma _{2\sigma }}{\gamma _\sigma }$$
(16)
and TMR
$$TMR^0=\frac{4\alpha P_1P_2}{(1+\alpha )^2(P_1+\alpha P_2)^2}.$$
(17)
Comparison of both the expressions for TMR \[Eq.(15) and (17)\] shows that Coulomb interactions can significantly increase the value of the magnetoresistance.
### B Fluctuations
Fluctuations in the system are studied within the generation-recombination approach for multi-electron channels. The Fourier transform of the correlation function of the quantity $`X`$ can be expressed as
$`S_{XX}(\omega )2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑te^{i\omega t}\left[X(t)X(0)X^2\right]`$ (18)
$`=4{\displaystyle \underset{n,m}{}}X_n\left[P(n,m;\omega ){\displaystyle \frac{p_n}{i\omega }}\right]X_mp_m,`$ (19)
where $`p_m`$ is the stationary value of the probability $`\widehat{p}`$ at the state $`m`$ \[given by Eq.(12)\], $`X_m`$ is the value of $`X`$ at this state. The conditional probability $`P(n,m;t)`$ to find the system in the state $`n`$ at time $`t`$, if it was in the initial state $`m`$ at $`t=0`$, satisfies the master equation (7), and its Fourier transform is given by $`P(n,m;\omega )=[i\omega \widehat{M}]_{nm}^1`$. The elements of the Green’s function $`G(n,m;\omega )[i\omega \widehat{M}]_{nm}^1p_n/i\omega `$ can be determined directly by matrix inversion and the result is
$$\widehat{G}(\omega )=\frac{\widehat{A}^+}{i\omega \lambda _+}\frac{\widehat{A}^{}}{i\omega \lambda _{}},$$
(20)
where $`\lambda _\pm =(\gamma _{}\gamma _{}\pm \mathrm{\Delta })/2`$ are the nonzero eigenvalues of the matrix $`\widehat{M}`$, $`\mathrm{\Delta }=\sqrt{(\gamma _{}\gamma _{})^2+4\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+}`$,
$`\widehat{A}^r={\displaystyle \frac{1}{D}}\left[\begin{array}{ccc}\mathrm{\Gamma }_{}^{}a_,^r& \mathrm{\Gamma }_{}^+a_,^r& \mathrm{\Gamma }_{}^+a_{,0}^r\\ \mathrm{\Gamma }_{}^{}a_,^r& \mathrm{\Gamma }_{}^+a_,^r& \mathrm{\Gamma }_{}^+a_{,0}^r\\ \mathrm{\Gamma }_{}^{}a_{,0}^r& \mathrm{\Gamma }_{}^{}a_{,0}^r& \mathrm{\Gamma }_{}^+a_{,0}^r\mathrm{\Gamma }_{}^+a_{,0}^r\end{array}\right]`$ (24)
corresponding to $`\lambda _r`$ ($`r=\pm `$), $`a_{\sigma ,\sigma }^r=\lambda _r\gamma _\sigma +\gamma _\sigma ^2+\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+`$, $`a_{\sigma ,\sigma }^r=\mathrm{\Gamma }_\sigma ^{}(\lambda _r+\gamma _{}+\gamma _{})`$, $`a_{\sigma ,0}^r=(\lambda _r+\gamma _\sigma )\mathrm{\Gamma }_\sigma ^{}+\mathrm{\Gamma }_\sigma ^+\mathrm{\Gamma }_\sigma ^{}`$, and $`D=\mathrm{\Delta }(\gamma _{}\gamma _{}\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+)`$. The Green’s function (20) is not singular for $`\omega 0`$ and therefore, one can easily separate the amplitudes of the noise resulting from fluctuation processes characterized by the relaxation time $`\tau _r=1/\lambda _r`$.
The fluctuations of the charge and the spin are expressed as
$`S_{NN}(\omega )=4e^2{\displaystyle \underset{\sigma ,\sigma ^{}}{}}G_{\sigma \sigma ^{}}(\omega )p_\sigma ^{}=`$ (25)
$`{\displaystyle \frac{4e^2\mathrm{\Gamma }_{}^{}\mathrm{\Gamma }_{}^{}}{\mathrm{\Delta }(\gamma _{}\gamma _{}\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+)^2}}{\displaystyle \underset{\sigma ,r}{}}r{\displaystyle \frac{\mathrm{\Gamma }_\sigma ^+\mathrm{\Gamma }_\sigma ^{}(\lambda _r+\mathrm{\Gamma }_\sigma ^{})}{i\omega \lambda _r}},`$ (26)
$`S_{MM}(\omega )=4{\displaystyle \frac{\mu _B^2}{4}}{\displaystyle \underset{\sigma ,\sigma ^{}}{}}\sigma \sigma ^{}G_{\sigma \sigma ^{}}(\omega )p_\sigma ^{}=`$ (27)
$`{\displaystyle \frac{\mu _B^2\mathrm{\Gamma }_{}^{}\mathrm{\Gamma }_{}^{}}{\mathrm{\Delta }(\gamma _{}\gamma _{}\mathrm{\Gamma }_{}^+\mathrm{\Gamma }_{}^+)^2}}{\displaystyle \underset{\sigma ,r}{}}r{\displaystyle \frac{\mathrm{\Gamma }_\sigma ^+(\gamma _\sigma +\mathrm{\Gamma }_\sigma ^+)(\lambda _r+\gamma _\sigma +\mathrm{\Gamma }_\sigma ^+)}{i\omega \lambda _r}},`$ (28)
where $`\mu _B`$ is the Bohr magneton.
The correlations between the currents $`I_{j\sigma }`$ and $`I_{j^{}\sigma ^{}}`$ in the tunnel junction $`j`$ and $`j^{}`$ for the electrons with the spin $`\sigma `$ and $`\sigma ^{}`$ are described by the power spectrum
$$S_{I_{j\sigma }I_{j^{}\sigma ^{}}}(\omega )=\delta _{jj^{}}\delta _{\sigma \sigma ^{}}S_{j\sigma }^{Sch}+S_{I_{j\sigma }I_{j^{}\sigma ^{}}}^c(\omega ),$$
(29)
where
$$S_{j\sigma }^{Sch}2e(I_{j\sigma }^++I_{j\sigma }^{})=2e^2\left[\mathrm{\Gamma }_{j\sigma }^+p_0+\mathrm{\Gamma }_{j\sigma }^{}p_\sigma \right]$$
(30)
is the high frequency ($`\omega \mathrm{}`$) limit of the shot-noise (the Schottky noise), which is the sum of the components corresponding to the tunneling current flowing to and from the particle. The frequency dependent part is expressed as
$`S_{I_{j\sigma }I_{j^{}\sigma ^{}}}^c(\omega )=2e^2(1)^{jj^{}}\{[\mathrm{\Gamma }_{j\sigma }^+G_{0\sigma ^{}}(\omega )\mathrm{\Gamma }_{j\sigma }^{}G_{\sigma \sigma ^{}}(\omega )]\mathrm{\Gamma }_{j^{}\sigma ^{}}^+p_0+[\mathrm{\Gamma }_{j^{}\sigma ^{}}^+G_{0\sigma }(\omega )\mathrm{\Gamma }_{j^{}\sigma ^{}}^{}G_{\sigma ^{}\sigma }(\omega )]\mathrm{\Gamma }_{j\sigma }^+p_0`$ (31)
$`+[\mathrm{\Gamma }_{j\sigma }^{}G_{\sigma 0}(\omega )\mathrm{\Gamma }_{j\sigma }^+G_{00}(\omega )]\mathrm{\Gamma }_{j^{}\sigma ^{}}^{}p_\sigma ^{}+[\mathrm{\Gamma }_{j^{}\sigma ^{}}^{}G_{\sigma ^{}0}(\omega )\mathrm{\Gamma }_{j^{}\sigma ^{}}^+G_{00}(\omega )]\mathrm{\Gamma }_{j\sigma }^{}p_\sigma \}.`$ (32)
The shot noise of the total current (including the displacement currents as well) is given by
$`S_{II}={\displaystyle \underset{j,j^{}}{}}{\displaystyle \frac{C_1^2C_2^2}{C^2C_jC_j^{}}}{\displaystyle \underset{\sigma ,\sigma ^{}}{}}[\delta _{jj^{}}\delta _{\sigma \sigma ^{}}S_{j\sigma }^{Sch}+S_{I_{j\sigma }I_{j^{}\sigma ^{}}}^c(\omega )].`$ (33)
## III Results
The analysis of the results we begin from a simplified situation, when the electrodes are made of paramagnetic metals. Next the device with ferromagnetic electrodes is considered. Since Coulomb interactions break the electron-hole symmetry, one can expect that the characteristics of the device for $`E_0<E_F`$ are different from those for $`E_0>E_F`$. Therefore, the both situations are considered separately.
### A Paramagnetic case
In the system with paramagnetic electrodes both the channels for electrons with the spin $``$ and $``$ are equivalent and the tunneling rates $`\gamma _j=\gamma _j=\gamma _j`$, $`\gamma _{}=\gamma _{}=\gamma `$, $`\mathrm{\Gamma }_{}^\pm =\mathrm{\Gamma }_{}^\pm =\mathrm{\Gamma }^\pm `$. The total current $`I_1=2e(f_1^+f_2^+)\gamma _1\gamma _2/[\gamma +\mathrm{\Gamma }^+]`$ differs form that for noninteracting electrons by the factor $`\mathrm{\Gamma }^+`$ in the denominator, which results form Coulomb interactions. In low temperatures there is a current blockade for the voltage within the range $`C/C_2<eV/(E_0E_F)<C/C_1`$. Dynamics of the fluctuations are characterized by the eigenvalues $`\lambda _+=\gamma +\mathrm{\Gamma }^+`$ and $`\lambda _{}=\gamma \mathrm{\Gamma }^+`$. Using Eqs.(29)-(31) and (20)-(24) one can derive the correlation function between the currents for electrons with the same spin as
$`S_{I_1I_1}(\omega )=2e^2{\displaystyle \frac{\gamma _1(f_1^+\mathrm{\Gamma }^{}+f_1^{}\mathrm{\Gamma }^+)}{\gamma +\mathrm{\Gamma }^+}}`$ (34)
$`2e^2{\displaystyle \frac{\gamma _1^2}{\gamma +\mathrm{\Gamma }^+}}\left[{\displaystyle \frac{f_1^+f_1^{}(\mathrm{\Gamma }^{})^2}{\omega ^2+\lambda _+^2}}{\displaystyle \frac{a_{}}{\omega ^2+\lambda _{}^2}}\right]`$ (35)
and between the different spins
$$S_{I_1I_1}(\omega )=2e^2\frac{\gamma _1^2}{\gamma +\mathrm{\Gamma }^+}\left[\frac{f_1^+f_1^{}(\mathrm{\Gamma }^{})^2}{\omega ^2+\lambda _+^2}+\frac{a_{}}{\omega ^2+\lambda _{}^2}\right],$$
(36)
where $`a_{}=(1+f_1^+)[f_1^+(\mathrm{\Gamma }^{+2}+2\gamma \mathrm{\Gamma }^+\gamma ^2)2(\mathrm{\Gamma }^+)^2]`$. Thus, the power spectrum of the total current through the left junction is expressed by
$`S_{I_1I_1}(\omega )=4e^2{\displaystyle \frac{\gamma _1(f_1^+\mathrm{\Gamma }^{}+f_1^{}\mathrm{\Gamma }^+)}{\gamma +\mathrm{\Gamma }^+}}`$ (37)
$`+8e^2{\displaystyle \frac{\gamma _1^2}{\gamma +\mathrm{\Gamma }^+}}{\displaystyle \frac{a_{}}{\omega ^2+\lambda _{}^2}}.`$ (38)
The noise corresponding to the eigenvalue $`\lambda _+`$ is completely cancelled. In the high-voltage regime the above formulae are much simpler, e.g. for $`V>0`$ the current $`I_1=2e\gamma _1\gamma _2/(\gamma _1+2\gamma _2)`$ and the Fano factor
$$_{11}\frac{S_{I_1I_1}(\omega =0)}{2eI_1}=1\frac{4\gamma _1\gamma _2}{(\gamma _1+2\gamma _2)^2}.$$
(39)
(see also \[\] and references therein).
For comparison in the case of noninteracting electrons ($`E_c=0`$) there are two independent channels and the power spectrum can be written as
$`S_{I_{1\sigma }^0I_{1\sigma }^0}^0(\omega )=2e^2{\displaystyle \frac{\gamma _{1\sigma }[f_1^+\mathrm{\Gamma }_\sigma ^{}+f_1^{}\mathrm{\Gamma }_\sigma ^+]}{\gamma _\sigma }}`$ (40)
$`4e^2{\displaystyle \frac{\gamma _{1\sigma }^2[f_1^+(\mathrm{\Gamma }_\sigma ^{})^2+f_1^{}(\mathrm{\Gamma }_\sigma ^+)^2]}{\gamma _\sigma (\omega ^2+\gamma _\sigma ^2)}}`$ (41)
for electrons with spin $`\sigma `$. For the paramagnetic electrodes and in the limit of a large positive $`V`$ one gets \[from Eqs.(16) and (40)\] the total current $`I_1^0=2e\gamma _1\gamma _2/(\gamma _1+\gamma _2)`$ and the Fano factor $`_{11}^0=12\gamma _1\gamma _2/(\gamma _1+\gamma _2)^2`$
Let us present also the correlation function between the currents through different tunnel junctions
$`\text{Re}[S_{I_1I_2}(\omega )]=4e^2{\displaystyle \frac{\gamma _1\gamma _2}{\gamma +\mathrm{\Gamma }^+}}{\displaystyle \frac{b_{12}}{\omega ^2+\lambda _{}^2}},`$ (42)
where $`b_{12}=\mathrm{\Gamma }^{+2}(4+f_1^++f_2^+2f_1^+f_2^+)+(f_1^++f_2^++2f_1^+f_2^+)(\mathrm{\Gamma }^{}\mathrm{\Gamma }^+)\gamma `$. Now, using Eq.(33) one gets the total power spectrum of the device
$`S_{II}(\omega )=4e^2{\displaystyle \frac{(C_2^2\gamma _1f_1^++C_1^2\gamma _2f_2^+)\mathrm{\Gamma }^{}+(C_2^2\gamma _1f_1^{}+C_1^2\gamma _2f_2^{})\mathrm{\Gamma }^+}{C^2(\gamma +\mathrm{\Gamma }^+)}}`$ (43)
$`+8e^2{\displaystyle \frac{\mathrm{\Gamma }_{12}^2(\mathrm{\Gamma }^{+2}+2\gamma \mathrm{\Gamma }^+\gamma ^2)\mathrm{\Gamma }_{12}\gamma _{12}\mathrm{\Gamma }^22\gamma _{12}^2\mathrm{\Gamma }^{+2}}{(\gamma +\mathrm{\Gamma }^+)(\omega ^2+\lambda _{}^2)}},`$ (44)
where $`\mathrm{\Gamma }_{12}=(C_2\gamma _1f_1^+C_1\gamma _2f_2^+)/C`$, $`\gamma _{12}=(C_2\gamma _1C_1\gamma _2)/C`$. One can check that in the zero-frequency limit $`S_{I_1I_1}(0)=`$Re$`[S_{I_1I_2}(0)]=S_{I_2I_2}(0)`$. Therefore, the Fano factors $`_{11}=_{12}=_{22}`$, which in the high-voltage range can be simply expressed as $`=14\gamma _1\gamma _2/(\gamma _1+2\gamma _2)^2`$.
We are also interested in charge and spin fluctuation induced by the flowing current. Using the formula (25) and (27) for the paramagnetic device one can write the charge-charge and the spin-spin correlation function as
$`S_{NN}(\omega )={\displaystyle \frac{8e^2\mathrm{\Gamma }^+\mathrm{\Gamma }^{}}{(\gamma +\mathrm{\Gamma }^+)(\omega ^2+\lambda _{}^2)}}`$ (45)
$`S_{MM}(\omega )={\displaystyle \frac{2\mu _B^2\mathrm{\Gamma }^+\mathrm{\Gamma }^{}}{(\gamma +\mathrm{\Gamma }^+)(\omega ^2+\lambda _+^2)}}.`$ (46)
From a frequency dependence of the correlation functions $`S_{NN}`$ and $`S_{MM}`$ one can assign the relaxation time corresponding to the charge and the spin fluctuations as $`\tau _{charge}=1/\lambda _{}`$ and $`\tau _{spin}=1/\lambda _+`$, respectively. One can check that the same result for the correlation functions can be derived from the two-level generation-recombination approach using $`S_{XX}(\omega )=4\text{var}(X)\tau /(\omega ^2\tau ^2+1)`$, where var$`(X)=X^2X^2`$ is the variance of the quantity $`X`$. Since $`\tau _{spin}>\tau _{charge}`$ then spin fluctuations occur in a low frequency regime, while the charge fluctuations in higher frequencies. The amplitude of the spin noise $`S_{MM}(\omega =0)`$ is larger than $`S_{NN}(\omega =0)`$ (in some cases the difference can be a few orders of magnitudes ). In the paramagnetic system the spin fluctuations, however, do not contribute to the current shot noise. The frequency dependence of the power spectrum (37) has then a Lorentzian form with the relaxation time $`\tau _{charge}`$.
### B Ferromagnetic electrodes and $`E_0<E_F`$
Let us first consider the ferromagnetic double tunnel barrier device, in which the particle level is below the Fermi level of the electrodes. A typical voltage dependence of the current is shown in Fig.1a. The $`I`$-$`V`$ function has a step like shape, with the current blockade for small voltages \[in the range $`C/C_1<eV/(E_FE_0)<C/C_2`$\] and the plateaux in the limit of large voltages, in which
$`I_1=\{\begin{array}{ccc}e\frac{\gamma _1\gamma _1(\gamma _2+\gamma _2)}{\gamma _{}\gamma _{}\gamma _2\gamma _2}\hfill & \text{for}& V(E_FE_0)/e,\\ e\frac{\gamma _2\gamma _2(\gamma _1+\gamma _1)}{\gamma _{}\gamma _{}\gamma _1\gamma _1}\hfill & \text{for}& V(E_FE_0)/e.\end{array}`$ (49)
We remind that according to our assumptions $`|V|\mathrm{\Delta }E,E_c`$ and the tunneling rates $`\gamma _{j\sigma }`$ are independent of $`V`$, even for the so called high voltages when the currents are given by Eq.(49). The current intensities (49) for large positive and negative voltages are different, in contrast to the case of noninteracting electrons, where both the $`I`$-$`V`$ steps are equal. Fig.1a shows that the height of the steps depends on the magnetic asymmetry of the electrodes, and an increase of the magnetic polarization $`P_1`$ in the left electrode reduces the current for $`V>0`$. If this electrode is made of a half-metallic ferromagnet (i.e. for $`P_1=1`$ and $`\gamma _1=0`$) the conducting channel corresponds only to electrons with the spin $``$, and there is the Coulomb blockade $`I_1=0`$ in low temperatures ($`k_BT(E_FE_0)`$) for any positive voltage. An electron with the spin $``$, which has tunneled form the right electrode into the particle, is captured there forever. The electron can neither tunnel to the left nor to the right electrode, and blocks the conducting channel for electrons with the spin $``$. Electronic transport can only occur for large negative voltages. Such the device works as a diode.
Using Eq.(49) one finds
$$TMR=\frac{4\alpha P_1P_2}{1P_1^2+2\alpha 2\alpha P_1P_2}$$
(50)
in the limit $`V(E_FE_0)/e`$. For comparison the value for noninteracting electrons
$$TMR^0=\frac{4\alpha P_1P_2}{1P_1^2+2\alpha 2\alpha P_1P_2+\alpha ^2(1P_2^2)},$$
(51)
is much smaller, especially in the system with asymmetric tunnel junctions ($`\alpha 1`$). One can say that Coulomb interactions enhance the value of TMR.
The power spectrum on the conducting step (for $`V>0`$) is given by
$`S_{I_1I_1}(\omega )=2e^2{\displaystyle \frac{\gamma _1\gamma _1(\gamma _2+\gamma _2)}{\gamma _{}\gamma _{}\gamma _2\gamma _2}}`$ (52)
$`{\displaystyle \frac{4e^2\gamma _1\gamma _1(\gamma _2+\gamma _2)}{\mathrm{\Delta }(\gamma _{}\gamma _{}\gamma _2\gamma _2)^2}}{\displaystyle \underset{r}{}}r\lambda _r{\displaystyle \frac{\lambda _ra+b}{\omega ^2+\lambda _r^2}}`$ (53)
where $`a=\gamma _1\gamma _1(\gamma _2+\gamma _2)`$ and $`b=\gamma _1^2\gamma _2(\gamma _2\gamma _1)+\gamma _1^2\gamma _2(\gamma _2\gamma _1)2\gamma _1\gamma _1\gamma _2\gamma _2`$. The eigenvalue in this case is $`\lambda _r=(\gamma _{}\gamma _{}+r\mathrm{\Delta })/2`$, and $`\mathrm{\Delta }=\sqrt{(\gamma _{}\gamma _{})^2+4\gamma _2\gamma _2}`$. The voltage dependence of the Fano factor is presented in Fig.1b. One can show that the zero-frequency power spectrum $`S_{I_jI_j^{}}(\omega =0)`$ corresponding to the currents through different tunnel junctions are equal, and thus, the Fano factors $`_{11}=_{12}=_{22}`$ for any model parameters (for any transition rates $`\gamma _{j\sigma }`$ at any voltage). In the regime of high-voltage its value is $`=1+2b/(\gamma _{}\gamma _{}\gamma _2\gamma _2)^2`$. If the coefficient $`b`$ is negative, then $`<1`$ and the noise is of the sub-Poissonian type. It occurs for $`2\alpha P_1^2(1P_2^2)<(1P_1P_2)(1P_1^2)`$. The transition from the sub-Poissonian to the super-Poissonian type of the current shot noise is a continuous process. In order to understand it we plotted in Fig.2 the frequency dependent part of the correlation functions $`S_{I_{1\sigma }I_{1\sigma ^{}}}^c(\omega =0)`$ \[given by Eq.(31)\] for the currents of electrons with different spins through the left junction in the high-voltage limit. One can expect competition between tunneling processes for electrons with the spin $``$ and $``$, which leads to an enhancement of the current noise. For simplicity the right electrode is taken paramagnetic, i.e. the source electrode can emit electrons with the same transition rate ($`\gamma _2=\gamma _2`$). The drain electrode is ferromagnetic and therefore, there is an asymmetry between the out-going channels for electrons with opposite spin directions, which is described by the magnetic polarization $`P_1`$. For $`P_1=0`$, the functions are equal $`S_{I_1I_1}^c(0)=S_{I_1I_1}^c(0)=S_{I_1I_1}^c(0)`$ and negative. It means that all tunneling events are anti-correlated, which leads to a reduction of the noise. An increase of the polarization $`P_1`$ increases the tunneling rate $`\gamma _1`$ for electrons with the spin $``$, they can faster leave the particle. Electrons with the opposite spin ($``$) spend a long time on the particle. It effects in the spin accumulation, which is responsible for an increase of $`S_{I_1I_1}^c(0)`$ and $`S_{I_1I_1}^c(0)`$. Their values can cross zero and achieve maxima for $`P_11`$. The function $`S_{I_1I_1}^c(0)`$ is always negative (for $`P_1>0`$). The process results an enhancement of the shot noise and the transition to the super-Poissonian range. The maximum value of the Fano factor $`=1+2\gamma _2/\gamma _1`$ occurs for the left electrode made of a half-metallic ferromagnet ($`P_1=1`$). Fig.2 shows also that a large asymmetry factor $`\alpha 1`$ between the left and the right tunnel barrier can prefer the transition to the super-Poissonian shot noise (see the dashed curves corresponding to $`\alpha =10`$).
### C Ferromagnetic electrodes and $`E_0>E_F`$
In the case of $`E_0>E_F`$ one can expect similar characteristics of our device to those presented above for $`E_0<E_F`$. It is really the case, but only for the high-voltage regime, where the $`I`$-$`V`$ curve has plateaux, whose level is given by Eq.(49). Fig.3 presents the voltage dependence of the current and the Fano factor. (Since the curves in the range of negative $`V`$ are very similar to those from Fig.1, we present the dependences for $`V>0`$ only). It is seen a resonant-like peak of the current in the range of moderate voltages, at $`E_0eV_2E_F`$ (i.e. for $`V/(|E_0E_F|/e)2`$ in Fig.3a). Its height can be much above the plateau level in the device with large asymmetry of the tunnel junctions. The most pronounced peak is for the device with the left electrode made of a half-metallic ferromagnet ($`P_1=1`$, $`\gamma _1=0`$). The total current, in this case, can be written as
$$I_1=e\frac{(f_1^+f_2^+)f_2^{}\gamma _1\gamma _2}{(1f_1^+f_2^+)\gamma _1+(1f_2^{+2})\gamma _2}.$$
(54)
It is worth noticing, that in this limit the current (54) and the occupation probability $`p_{}`$, $`p_{}`$, $`p_0`$ are independent of the transition rate $`\gamma _2`$. The current peak is the resonant-like transition of electrons through the particle level and the current blockade effect in the low and the high-voltage range. For a small $`V`$ the position of the particle level $`E_0eV_2`$ is above the Fermi level $`E_F`$ of the right electrode and electrons can not tunnel to the particle, whereas in a high-voltage range there is a Coulomb blockade of the conducting channel by an electron with spin $``$ captured on the particle. The width of the current peak depends on the smearing of the Fermi surface and decreases with a decreasing temperature.
The $`I`$-$`V`$ curve (54) resembles that obtained in the case of resonant tunneling through double barrier in semiconductors (the Esaki diode). The nature of the both tunneling effects is, however, different. In the present case the negative differential resistance (NDR) is caused by Coulomb interactions between electrons on the particle (by the Coulomb blockade effect). In the Esaki diode the charge accumulation in the well is irrelevant for electronic transport and the NDR results from a shift of the conduction band of the source electrode out of the resonant tunneling range (see \[\], which considered coulomb interactions in resonant tunneling as well). The width of the peak depends in the Esaki diode on the electronic structure of the device. It can be smeared due to fluctuations of the bottom of the potential well. In our model the position of $`E_0`$ is fixed and the broadening of the peak results only from the thermal distribution of electrons around the Fermi level.
Fig.3b shows the voltage dependence of the Fano factor. Its value is below unity in the low-voltage range and rapidly increases when $`E_0eV_2`$ crosses the Fermi level $`E_F`$ (i.e. for $`V/(|E_0E_F|/e)>2`$ in Fig.3b). The increase of $``$ is only in a narrow range of $`V`$, in the same in which the NDR effect is observed. In the high-voltage regime the noise is super-Poissonian for the most situations exhibited in Fig.3b. The voltage dependence of the Fano factor in the present case (see the curve for $`P_1=1`$ in Fig.3b) is qualitatively different from that in the resonant tunneling diode , where $``$ shows a large peak in the NDR region. The origin of the Fano peak is activation of interaction-induced fluctuations of the band bottom in the quantum well, when the system passes to the off-resonant electronic transport. As we have explained already in the previous section, the high value of $``$ in our system is related with the asymmetry of the conducting channels for electrons with the opposite spin directions.
Flowing electrons induce the charge and the spin fluctuations on the particle with the characteristic frequencies $`1/\tau _{charge}=\lambda _{}`$ and $`1/\tau _{spin}=\lambda _+`$, respectively. These fluctuations should be seen in the current noise. Therefore, we separate the Schottky term $`S_I^{Sch}`$ from the current noise and perform the spectral decomposition of the frequency dependent part $`S_{II}^c(\omega )`$. The total power spectrum can be expressed as
$$S_{II}(\omega )=S_I^{Sch}+S_{II}^{c+}(\omega )+S_{II}^c(\omega ),$$
(55)
where $`S_I^{Sch}=(C_2^2S_{I1}^{Sch}+C_1^2S_{I2}^{Sch})/C^2`$ and
$$S_{II}^{c\pm }(\omega )=\pm \left(\frac{C_1C_2}{C}\right)^2\underset{j,j^{}}{}\frac{\lambda _\pm }{C_jC_j^{}}\frac{\lambda _\pm a_{jj^{}}+b_{jj^{}}}{\omega ^2+\lambda _\pm ^2}.$$
(56)
The coefficient $`a_{jj^{}}`$ and $`b_{jj^{}}`$ are determined from Eq.(31), (20), (24) and (12). The voltage dependences of these terms for $`\omega =0`$ are presented in Fig.4. The system is the same as studied above (Fig.3), in which the asymmetry between the tunnel barriers is $`\alpha =10`$. The transition rates $`\gamma _{2\sigma }`$ are larger than $`\gamma _{1\sigma }`$, and therefore $`S_{I2}^{Sch}>S_{I1}^{Sch}`$. In the low-voltage range the term $`S_{I2}^{Sch}`$ is very large and dominates in $`S_I^{Sch}`$. The Fano factor $`=[S_I^{Sch}+S_{II}^{c+}(0)+S_{II}^c(0)]/(2eI)`$ is, however, below unity. In the considered system we change the magnetic polarization $`P_1`$, which influences of $`S_{I1}^{Sch}`$, but it is irrelevant for $`S_I^{Sch}`$. It explains, why all the curves in Fig.4a are so close to each other.
Fig.4b and 4c show the terms $`S_{II}^{c+}(0)`$ and $`S_{II}^c(0)`$ corresponding to the contribution of the spin and the charge fluctuations to the current noise. They are negative in the low-voltage range and positive for larger voltages. This indicates a change of current correlations when the particle level crosses the Fermi level ($`E_0eV_2E_F`$). The value $`S_{II}^{c+}(0)`$ strongly increases with an increase of the magnetic polarization $`P_1`$. Since $`S_{II}^c`$ and $`S_I^{Sch}`$ (see Fig.4c and 4a) are weakly dependent on $`P_1`$, it is evident that $`S_{II}^{c+}`$ is responsible for an enhancement of the Fano factor. Frequency dependent measurements of the current noise can confirm our prediction, that low frequency fluctuations dominate in the super-Poissonian noise in ferromagnetic tunnel junctions.
## IV SUMMARY
Summarizing, our sequential tunneling studies, performed in the ferromagnetic double barrier device with the atomic size particle, showed a few interesting effects. First, Coulomb interactions lead to an enhancement of the TMR effect. Second, an electron-hole symmetry is broken in the system, due to Coulomb interactions. The characteristics of the ferromagnetic device, with the electronic state $`E_0`$ of the spacer particle below the Fermi level $`E_F`$ of the electrodes, are qualitatively different from those for the case of $`E_0>E_F`$. We showed that the system, in which $`E_0<E_F`$ and one electrode is made of a half-metallic ferromagnet, can operate as a diode. When $`E_0>E_F`$ the device showed the NDR effect, which is better pronounced for ferromagnetic electrodes with different magnetic polarizations. Third, the transition from the sub-Poissonian to the super-Poissonian current noise is a continuous process, which depends on the magnetic asymmetry between the tunneling channels for electrons with the spin $``$ and $``$. The asymmetry between the left and the right tunnel barrier can facilitate the transition to the super-Poissonian range. Spin fluctuations are relevant for the super-Poissonian current noise and they are activated in the Coulomb blockade regime. The charge fluctuations are responsible for the sub-Poissonian current noise. The spin and the charge fluctuations have distinct relaxation times $`\tau _{spin}>\tau _{charge}`$, which can be observed in frequency dependent measurements of the power spectrum in a low and in a high-frequency range, respectively.
###### Acknowledgements.
The paper is supported from the State Committee for Scientific Research Republic of Poland within Grant No. 2 P03B 075 14. |
warning/0003/physics0003048.html | ar5iv | text | # I Introduction.
## I Introduction.
The uniqueness problem for the sources of the evoked potential in the brain is a relevant research question due to its role in the development of cerebral electric tomography, , , . Since long time ago, it is known that the general inverse problem of the determination of volumetric sources from the measurement of the potential at a surface is not solvable in general,. However, under additional assumptions about the nature of the sources, solutions can be obtained ,,. The supplementary assumptions can be classified in two groups: the physically grounded ones, which are fixed by the nature of the physical problem and the ones which are imposed by invoking their mathematical property of determining a solution, but having in another hand, a weak physical foundation. The resumed situation implies that the determination of physical conditions implying the uniqueness of the sources for the evoked potentials remains being an important subject of study. Results in this direction could avoid the imposition of artificial conditions altering the real information on the sources to be measured.
The question to be considered in this work is the uniqueness of the sources for evoked potentials under the assumption that these sources are localized over surfaces. This issue was also treated in Ref. by including also some specially defined volumetric sources. The concrete aim here is to present a derivation of the results enunciated in for the case of open surfaces and to generalize it for a wider set of surfaces including closed ones.
We consider that the results enunciated in Ref. are valid and useful ones. Even more, we think that a relevant merit of that paper is to call for the attention to the possibility for the uniqueness for classes of surface density of sources. Specifically, in our view, the conclusion stated there about the uniqueness of the sources of evoked potentials as restricted to sources distributed in open surfaces is effectively valid. In the present work, the central aim is to extend the result for a wider set of surfaces including closed ones by also furnishing an alternative way to derive the uniqueness result. The uniqueness problem for the special class of volumetric sources discussed in is not considered here in any way.
The physical system under consideration is conformed by various volumetric regions, each of them having a constant value of the conductivity, separated by surface boundaries at which the continuity equations for the electric current is obeyed. It should pointed out that the special volumetric sources examined in Ref. are not addressed here. The precise definition of the generators under examination is the following. The sources are assumed to be defined by continuous and smooth surface densities lying over a arbitrary but finite number of smooth open or closed surfaces. The unique constraint to be imposed on these surfaces is that there is no nesting among them. That is, there is no closed surface at which interior another open or closed of the surfaces resides. This class of supports expands the one considered in Ref. and in our view is sufficiently general to create the expectative for the practical applications of the results. It should be stressed that the boundaries between the interior metallic regions are not restricted by the ”non-nesting” condition. That is, the fact that the skull and the few boundaries between cerebral tissues can be visualized as nearly closed surface does not pose any limitation on the conclusion. The ”non-nesting” condition should be valid only for the surfaces in which the sources can be expected to reside. For example, if by any mean we are sure that the sources stay at the cortex surface, then the uniqueness result apply whenever the portion of the cortex implied does not contains any closed surface.
The paper is organized as follows. An auxiliary property is derived in the form of a theorem in the Section II. In Section III the proof of uniqueness for the kind of sources defined above is presented.
## II Green Theorem and Field Vanishing Conditions
Let us consider the potential $`\varphi `$ generated by a source distribution concentrated in the ”non-nested” set of open or closed surfaces defined in last Section, which at the same time are contained within a compact and simply connected spatial region $`R.`$The set $`R,`$ as explained before, is formed by various connected subregions $`R_i,i=0,1,\mathrm{}n`$ each of them filled with a metal having a constant conductivity $`\sigma _i`$. Also, let $`B_{ij}`$ the possibly but non necessarily existing, boundary between the subregions $`R_i`$ and $`R_j`$ and $`B_0`$the boundary of $`R.`$ For the sake of a physical picture, we can interpret $`B_0`$ as the surface of the skull, $`R`$ as the interior of the head and the subregions $`R_i`$ as the ones containing the various tissues within the brain. It is defined that the exterior space of the head corresponds to $`R_0`$. In addition, let $`S_i,i=1,\mathrm{}m`$ the surfaces pertaining to the arbitrary but finite set $`S`$ of non-nested open or closed surfaces in which the sources are assumed to be localized. The above mentioned definitions are illustrated in Fig.1.
Then, the Poisson equation satisfied by the potential $`\varphi `$ in the interior region of $`R`$ can be written as
$`^2\varphi \left(\stackrel{}{x}\right)`$ $`=`$ $`{\displaystyle \frac{g\left(\stackrel{}{x}\right)}{\sigma \left(\stackrel{}{x}\right)}},`$ (1)
$`g\left(\stackrel{}{x}\right)`$ $`=`$ $`\stackrel{}{}.\stackrel{}{J}\left(\stackrel{}{x}\right),`$ (2)
where$`\stackrel{}{J}`$ are the impressed currents (for example, generated by the neuron firings within the brain) and the space dependent conductivity is defined by
$$\sigma \left(\stackrel{}{x}\right)=\sigma _ifor\stackrel{}{x}R_i.$$
(3)
It should be noticed that the conductivities are different from zero only for the internal regions to $`R.`$ The vacuum outside is assumed to have zero conductivity and the field satisfying the Laplace equation. In addition outside the support of the sources where $`g=0`$ the Laplace equation is also satisfied.
The usual boundary conditions within the static approximation, associated to the continuity of the electric current at the boundaries, take the form
$$\sigma _i\frac{\varphi }{n_i}_{xB_{ij}}=\sigma _j\frac{\varphi }{n_j}_{xB_{ij}},$$
(4)
where $`n_i`$ symbolizes the directional derivative along a line normal to $`B_{ij}`$ but taken in the limit of $`x>B_{ij}`$from the side of the region $`R_i.`$
A main property is employed in this work in obtaining the claimed result. In the form of a theorem for a more precise statement it is expressed as
Theorem.
Let $`\varphi `$ is a solution of the Laplace equation within an open and connected spatial region $`R^{}`$. Assume that $`\phi `$have a vanishing electric field over an open section of certain smooth surface $`S^{}`$ which is contained in an open subset $`Q`$ of $`R^{}`$. Let the points of the boundaries between $`Q`$ and $`R^{}`$ have a minimal but finite distance among them. Then, the potential $`\varphi `$ is a constant over any open set contained in $`R^{}.`$
As a first stage in the derivation of this property, let us write the Green Theorem as applied to the interior of the open region $`Q`$defined in the Theorem 1 in which a field $`\phi `$ satisfies the Laplace equation. Then, the Green Theorem expresses $`\phi `$ evaluated at a particular interior point $`\stackrel{}{x}`$ in terms of itself and its derivatives at the boundary $`B_Q`$ as follows.
$$\phi \left(\stackrel{}{x}\right)=_{B_Q}𝑑\stackrel{}{s^{^{}}}.\left(\frac{1}{\left|\stackrel{}{x}\stackrel{}{x^{^{}}}\right|}\stackrel{}{_x^{^{}}}\phi \left(\stackrel{}{x^{^{}}}\right)\stackrel{}{_x^{^{}}}\left(\frac{1}{\left|\stackrel{}{x}\stackrel{}{x^{^{}}}\right|}\right)\phi \left(\stackrel{}{x^{^{}}}\right)\right)$$
(5)
where the integral is running over the boundary surface $`B_Q`$ which is described by the coordinates $`\stackrel{}{x^{^{}}}.`$ This relation expresses the potential as a sum of surface integrals of the continuous and bounded values of $`\phi `$ and its derivatives. Those quantities are in addition analytical in all the components of $`\stackrel{}{x},`$ if the point have a finite minimal distance to the points in $`B_Q.`$ These properties follow because $`Q`$ $`R^{}`$ and then, $`\phi `$ satisfies the Laplace equation in any open set in which $`Q`$ and its boundary is included. But, due to the finite distance condition among the point $`\stackrel{}{x}`$ and the points of $`B_Q`$, the expression (5) for $`\phi `$ should be an analytical function of all the coordinates of $`\stackrel{}{x}.`$ Figure 2 depicts the main elements in the formulation of the Green Theorem.
Further, let us consider that $`S^{}`$ is siting inside the region $`Q.`$ Then, as this surface is an equipotential and also the electric field over it vanishes, it follows that no line of force can have a common point with it. This is so because the divergence of the electric field vanishes, then it is clear that the existence of nonvanishing value of the electric field at another point of the line of force will then contradicts the assumed vanishing of the divergence. Therefore, the lines of forces in any sufficiently small open neighborhood containing a section of $`S^{}`$ should tend to be parallel to this surface on approaching it, or on another hand, the electric field should vanish. Next, it can be shown that in such neighborhoods the lines of forces can not tend to be parallel.
Let us suppose that lines of forces exist and tend to be tangent to the surface $`S^{}`$ and consider the integral form of the irrotational property of the electric field as
$$_C\stackrel{}{E}.d\stackrel{}{l}=_{C_1}\stackrel{}{E}.d\stackrel{}{l}+_{C_2}\stackrel{}{E}.d\stackrel{}{l}=_{C_1}\stackrel{}{E}.d\stackrel{}{l}=0$$
(6)
where the closed curve $`C`$ is constructed as follows: the piece $`C_1`$coincides with a line of force, the piece $`C_2`$ is fixed to rest within the surface $`S^{}`$ and the other two pieces necessary to close the curve are selected as being normal to the assumed existing family of lines of forces. The definitions are illustrated in Fig. 3. By construction, the electric field is colinear with the tangent vector to $`C_1`$ and let us assume that we select the segment of curve $`C_1`$for to have a sufficiently short but finite length in order that the cosine associated to the scalar product will have a definite sign in all $`C_1`$. This is always possible because the field determined by (5) should be continuous. Then Eq. (6) implies that the electric field vanish along all $`C_1`$ as a consequence of the integrand having a definite sign and then should vanish identically. Since this property is valid for any curve pertaining to a sufficiently small open interval containing any particular open section of $`S^{},`$ it follows that in certain open set containing $`S^{}`$there will be are no lines of forces, or what is the same, the electric field vanish.
To finish the proof of the theorem, it follows to show that if $`\phi `$ and the electric field vanish within a certain open neighborhood $`N,`$ included in an arbitrary open set $`O`$ pertaining to the region $`R^{}`$ in which the Laplace equation is obeyed, then $`\phi `$ and the electric field vanish in all $`O`$ . Consider first that $`Q`$ is an open set such that $`OQ`$ and also suppose that the smallest distance form a point in $`O`$ to the boundary $`B_Q`$ of $`Q`$ has the finite value $`\delta `$. Then, the Green Theorem (5) as applied to the region $`Q`$ expresses that the minimal radius of convergence of $`\phi `$considered as analytical function of any of the coordinates is equal or greater than $`\delta .`$
Imagine now a curve $`C`$ starting in an interior point $`P`$ of $`N`$ and ending at any point $`P_1`$ of $`O.`$ Assume that$`C`$ is formed by straight lines pieces (See Fig. 4). It is then possible to define $`\phi `$ as a function of the length of arc $`s`$ of $`C`$ as measured form the point $`P`$. It should be also valid that in any open segment of $`C,`$ not including the intersection point of the straight lines, the potential $`\phi `$is an analytical function of $`s.`$Furthermore, let consider$`C`$as partitioned in a finite number of segments of length $`\sigma <\delta .`$ Suppose also, that the intersection points of the straight lines are the borders of some of the segments. It can be noticed that $`\phi `$ vanishes in any segment of $`C`$ starting within $`N`$ because it vanishes in $`N`$exactly. Thus, if $`\phi `$and the electric field are not vanishing along all $`C,`$there should be a point over the curve in which the both quantities do not vanish for an open region satisfying $`s>s_o,`$ and vanish exactly for another open interval obeying $`s<s_o`$. However, in this case, all the derivatives of $`\phi `$ of the electric field over $`s`$ vanish at $`s_o.`$This property in addition with the fact that the Taylor series around $`s_o`$ should have a finite radius of convergence $`r>\delta ,`$ as it assumed in the Theorem 1, leads to the fact that $`\phi `$ and the electric field should vanish also for $`s>s_o.`$ Henceforth, the conclusion of the Theorem 1 follows: the potential $`\phi `$ and its corresponding electric field vanish at any interior point of $`R^{}.`$
## III Uniqueness of the Non-Nesting surface sources
Let us argue now the uniqueness of the sources which are defined over a set of non nested surfaces $`S`$ producing specific values of the evoked potential $`\varphi `$ at the boundary $`B_0`$of the region $`R.`$ For this purpose it will be assumed that two different source distributions produce the same evoked potential over $`B_0.`$ The electrostatic fields in all space associated to those sources should be different as functions defined in all space. They will be called $`\varphi _{1\text{ }}`$and $`\varphi _2.`$ As usual in the treatment of uniqueness problems in the linear Laplace equation, consider the new solution defined by the difference $`\varphi =`$ $`\varphi _1\varphi _2.`$ Clearly $`\phi `$ corresponds to sources given by the difference of the ones associated to$`\varphi _1`$ and $`\varphi _2.`$ It is also evident that $`\varphi `$ has vanishing values at $`B_0.`$ Then, since the sources are localized at the interior of $`R`$ and $`\varphi `$ satisfies the Laplace equation with zero boundary condition at $`B_0`$ and at the infinity, it follows that the field vanishes in all $`R_0,`$that is, in the free space outside the head. Therefore, it follows that the potential and the electric field vanish in all$`B_0`$when approaching this boundary from the free space ($`R_0).`$ The continuity of the potential, the boundary conditions (3) and the irrotational character of the electric field allows to conclude that $`\varphi `$ and the electric field also vanish at any point of $`B_0`$ but now when approaching it from any interior subregion $`R_i`$ having a boundary $`B_{i0}`$ with the free space. Moreover, if the boundary surface of any of these regions which are in contact with the boundary of $`R`$ is assumed to be smooth, then it follows from Theorem 1 that the potential $`\varphi `$ and its the electric field vanish in all the open subsets of $`R_i`$ which points are connected through its boundary $`B_{i0}`$with free space by curves non-touching the surfaces of $`S`$. It is clear that this result hold for all the open subsets of these $`R_i`$ in which Laplace equation is satisfied excluding those which are also residing inside one of the closed surfaces $`S_i`$ in the set $`S.`$
It is useful for the following reasoning to remark that if we have any boundary $`B_{ij}`$ between to regions $`R_{i\text{ }}`$ and $`R_j,`$ and the potential $`\varphi `$ and the electric field vanish in certain open (in the sense of the surface) and smooth regions of it, then Theorem 1 implies that the potential and the electric field also vanish in all the open subsets of $`R_i`$ and $`R_j`$ which are outside any of the closed surfaces in $`S.`$ Since the sources stay at the surfaces in $`S`$ the field $`\varphi `$ in some open region of $`R`$ included inside certain of the closed surfaces $`S_i`$ will not necessarily satisfy the Laplace equation in any interior point of $`R`$ and Theorem 1 is not applicable.
Let us consider in what follows a point $`P`$ included in a definite open vicinity of a subregion $`R_i.`$ Suppose also that $`P`$ is outside any of the closed surfaces in $`S`$ . Imagine a curve $`C`$which join $`P`$ with the free space and does not touch any of the surfaces in $`S`$. It is clear that, if appropriately defined, $`C`$ should intersect a finite number of boundaries $`B_{ij}`$ including always a certain one $`B_{j0}`$with free space. Let us also assume that $`C`$is adjusted in a way that in each boundary it crosses, the intersection point is contained in a smooth and open vicinity (in the sense of the surface) of the boundary (See Fig. 1 and 5). Then, it also follows that the curve $`C`$can be included in open set $`O_C`$ having no intersection with the non-nested surfaces in $`S.`$ This is so because the region excluding the interior of the closed surfaces in $`S`$ is also connected if the $`S_i`$are disjoint . But, from Theorem 1 it follows that $`\varphi `$and the electric field must vanish in all $`O_C.`$ This should be the outcome because the successive application of the Theorem 1 to the boundaries intersected by the curve $`C`$ permits to recursively imply the vanishing of $`\varphi `$and the electric field in each of the intersections of $`O_C`$ with the subregions $`R_i`$ through which $`C`$ passes. The first step in the recursion can be selected as the intersection of $`C`$ with $`B_{j0}`$ at a point which by assumption is contained in an open neighborhood of the boundary $`B_{j0}`$. As the electric field and $`\varphi `$ vanish at free space, the fields in the first of the considered intersection of $`Oc`$should vanish. This fact permits to define another open and smooth neighborhood of the next boundary intersected by $`C`$in which the field vanish and so on up to the arrival to the intersection with the boundary of the region $`R_i`$containing the ending of $`C`$at the original point P. Therefore, the electric field and the potential should vanish at an arbitrary point $`P`$ of $`R`$ with only two restrictions: 1) $`P`$ to be contained in an open neighborhood of some $`R_i`$ and 2) $`P`$ to reside outside any of the surfaces in $`S.`$ Thus, it is concluded that the difference solution $`\varphi `$ and its corresponding electric field, in all the space outside the region containing the sources vanish. Henceforth, it implies that the difference between the two source distributions also should be zero over any of the open surface in the set $`S.`$ This is necessary because the flux going out from any small piece of the considered surface is zero, which means that the assumed continuous density of surface sources exactly vanish. This completes the proof of the conclusion of Ref. in connection with sources supported by open surfaces. It only rests to show that the sources are also null over the closed $`S_i.`$
Before continuing with the proof, it is illustrative to exemplify from a physical point of view how the presence of nested surfaces among the $`S_i`$ destroys the uniqueness. For this aim let us let us consider that a closed surface $`S_i`$ has another open or closed of the surface $`S_j`$ properly contained inside it. That means that an open set containing $`S_j`$ is contained inside $`S_i.`$Imagine also that $`S_i`$ is interpreted as the surface of a metal shell connected to the ground; that is, to a zero potential and that the surface $`S_j`$ is the support of an arbitrary density of sources. As it is known from electrostatics theory, the charge density of a metal connected to the ground is always capable to create a surface density of charge at $`S_i`$ such that it exactly cancels the electric field and the potential at the outside of $`S_i,`$ in spite of the high degree of arbitrariness of the charge densities at the interior. That is, for nested surfaces in $`S`$, it is not possible to conclude the uniqueness, because at the interior of a nesting surface, and distributed over the nested ones, arbitrary source distributions can exist which determine exactly the same evoked potential at the outside boundary $`B_0`$.
Let us finally show that if no nesting exists the uniqueness also follows. Consider any of the closed surfaces, let say $`S_i.`$ As argued before $`\varphi `$ and the electric field vanish at any exterior point of $`S_i`$ pertaining to certain open set containing $`S_i.`$ Then, the field created by the difference between the sources associated to the two different solutions assumed to exist should be different from zero only at the interior region. That zone, in the most general situation can be filled by a finite number of metallic bodies with different but constant conductivities. The necessary vanishing of the interior field follows from the exact conservation of the lines of forces for the ohmic electric current as expressed in integral form by
$$𝑑\stackrel{}{s}.\sigma \left(\stackrel{}{x}\right)\stackrel{}{E}\left(\stackrel{}{x}\right)=0.$$
(7)
Let us consider a surface $`T`$defined by the all the lines of forces of the current vector passing through an arbitrarily small circumference $`c`$ which sits on a plane being orthogonal to a particular line of force passing through its center. Let the center be a point at the surface $`S_i`$ . Because, the above defined construction, all the flux of the current passing trough the piece of surface of $`S_i(`$which we will refer as $`p)`$ intersected by $`T`$ is exactly equal to the flux through any intersection of $`T`$with another surface determining in conjunction with $`p`$ a closed region. By selecting a sufficiently small radius for the circumference $`c`$ it can be noticed that the sign of the electric field component along the unit tangent vector to the central line of forces should be fixed. This is so because on the other hand there will be an accumulation of charge in some closed surface. Now, let us consider the fact that the electric field is irrotational and examine a line of force of the current density which must start at the surface $`S_i.`$ It should end also at $`S_i`$, because in another hand the current density will not be divergence less. After using the irrotational condition for the electric field in the form
$$_C\stackrel{}{E}.d\stackrel{}{l}=_{C_1}\stackrel{}{E}.d\stackrel{}{l}+_{C_2}\stackrel{}{E}.d\stackrel{}{l}=_{C_1}\stackrel{}{E}.d\stackrel{}{l}=0$$
(8)
in which $`C_1`$ is the line of force starting and ending at $`S_i`$ and $`C_2`$ is a curve joining the mentioned points at $`S_i`$ but with all its points lying outside $`S_i`$ where $`\varphi =\varphi _1`$-$`\varphi _2`$and the electric field vanish. Let us notice that the electric field and the current have always the same direction and sense as vectors, because the electric conductivity is a positive scalar. In addition, as it is argued above, the current can not reverse the sign of its component along the tangent vector of line of forces. Therefore, it follows that also the electric field can‘t revert the sign of its component along a line of force. Thus, the integrand of the line integral over the $`C_1`$ curve should have a definite sign at all the points, hence implying that $`\varphi `$and the electric field should vanish exactly in all $`C_1.`$ Resuming, it follows that the electric field vanish also at the interior of any of the closed surfaces $`S_i.`$ Therefore, the conclusion arises that the difference solution $`\varphi =\varphi _1`$-$`\varphi _2=0`$ in all the space, thus showing that the evoked potential at $`B_0`$ uniquely fixes the sources when they have their support in a set of non nesting surfaces $`S.`$
Acknowledgments
We would like to thank the helpful discussions with Drs. Augusto González , Jorge Riera and Pedro Valdés. One of the authors ( A.C.) also would like acknowledge the support for the development of this work given by the Christopher Reynolds Foundation (New York,U.S.A.) and the Center of Theoretical Studies of Physical Systems of the Clark Atlanta University (Atlanta, U.S.A). The support of the Associateship Programme of the Abdus Salam International Centre for Theoretical Physics (Trieste Italy) is also greatly acknowledged.
Figure Captions
Fig.1. An illustration of a simply connected region $`R`$ constituted in this case by only two simply connected subregions $`R_1`$ and $`R_2`$ having a boundary $`B_{12}.`$The boundary with free space is denoted by $`B_0.`$ The set of non-nesting surfaces $`S`$have four elements $`S_i`$ , $`i=1,\mathrm{..4}.`$ two of them open and other two closed ones. A piece wise straight curve $`C`$ joining any interior point $`P`$ of $`R`$ and a point $`O`$ in the free space is also shown.
Fig.2. Picture representing the region $`Q`$ in which a field $`\phi `$ satisfies the Laplace equation and its value at the point $`\stackrel{}{x}`$ is given by the Green integral (5).
Fig.3. The contour employed in the line integral in Eq. (6).
Fig.4. Picture of the region $`R_i`$and the open neighborhood $`N`$ in which the field $`\phi `$ vanish exactly . A piece wise straight line curve $`C`$joining a point $`PN`$and certain point $`P_1`$ in $`R_i`$ is also shown.
Fig.5. Scheme of the curve $`C`$ and the open region$`O_C`$ containing it. |
warning/0003/math0003142.html | ar5iv | text | # References
Deformation Quantization of Coadjoint Orbits
M. A. Lledó
Dipartimento di Fisica, Politecnico di Torino,
Corso Duca degli Abruzzi 24, I-10129 Torino, Italy, and
INFN, Sezione di Torino, Italy.
e-mail: lledo@athena.polito.it
## Abstract
A method for the deformation quantization of coadjoint orbits of semisimple Lie groups is proposed. It is based on the algebraic structure of the orbit. Its relation to geometric quantization and differentiable deformations is explored.
Let $`G`$ be a complex Lie group of dimension $`n`$ and $`G_R`$ a real form of $`G`$. Let $`𝒢`$ and $`𝒢_R`$ be their respective Lie algebras with Lie bracket $`[,]`$. As it is well known, $`𝒢_R^{}`$ has a Poisson structure,
$$\{f_1,f_2\}(\lambda )=<[(df_1)_\lambda ,(df_2)_\lambda ],\lambda >,f_1,f_2C^{\mathrm{}}(𝒢_R^{}),\lambda 𝒢_R^{}.$$
(1)
Choosing a basis $`\{X_1,\mathrm{}X_n\}`$ of $`𝒢_R`$ and its dual, $`\{\xi ^1,\mathrm{}\xi ^n\}`$, the Poisson bracket can be written as
$$\{f_1,f_2\}(x)=c_{ij}^k\lambda _k\frac{f_1}{x_i}\frac{f_2}{x_j},x=\underset{i=1}{\overset{n}{}}x_i\xi ^i𝒢^{}.$$
Notice that this Poisson bracket is never symplectic, in particular it is 0 at the origin. Under the action of $`gG_R`$, it satisfies $`g^{}\{f_1,f_2\}=\{g^{}f_1,g^{}f_2\}`$, so $`G_R`$ is a group of automorphisms of the Poisson algebra $`C^{\mathrm{}}(𝒢_R^{})`$. The action of $`G_R`$ on $`𝒢_R^{}`$ is not transitive, so $`𝒢_R^{}`$ is foliated in orbits. This foliation coincides with the foliation given by the Hamiltonian vector fields of (1). So the orbits of the coadjoint action of a Lie group are symplectic manifolds.
We want to describe formal deformations of the Poisson algebra $`C^{\mathrm{}}(\mathrm{\Theta })`$ ($`\mathrm{\Theta }`$ is a coadjoint orbit) or of some subalgebra of it. It is convenient to work with the complexification of the Poisson algebra.
An associative algebra $`𝒜_h`$ over $`[[h]]`$ is a formal deformation of a Poisson algebra $`(𝒜,\{,\})`$ over $``$ if there exists an isomorphism of $`[[h]]`$-modules $`\psi :𝒜[[h]]𝒜_h`$ satisfying the following properties:
a. $`\psi ^1(F_1F_2)=f_1f_2`$ mod($`h`$) where $`F_i𝒜_h`$ are such that $`\psi ^1(F_i)=f_i`$ mod$`(h)`$, $`f_i𝒜`$. (By mod$`(h)`$ we mean that the projections $`p:𝒜[[h]]𝒜[[h]]/h𝒜[[h]]`$ of both quantities coincide).
b. $`\psi ^1(F_1F_2F_2F_1)=h\{f_1,f_2\}`$ mod($`h^2`$).
An example of interest for our purposes is the polynomial algebra on $`𝒢^{}`$. A formal deformation of $`\text{Pol}(𝒢^{})`$ is given by the algebra $`U_h=T_{[[h]]}(𝒢)/_h`$, where $`T_{[[h]]}(𝒢)`$ is the tensor algebra over $`[[h]]`$ and $`_h`$ is the proper two sided ideal
$$_h=\underset{X,Y𝒢}{}T_{[[h]]}(𝒢)(XYYXh[X,Y])T_{[[h]]}(𝒢)T_{[[h]]}(𝒢).$$
The isomorphism $`\psi :\text{Pol}(𝒢^{})U_h`$ is not canonical. A possible choice is in terms of a Poincaré-Birkhoff-Witt basis,
$$\psi (x_{i_1}x_{i_2}\mathrm{}x_{i_k})=X_{i_1}X_{i_2}\mathrm{}X_{i_k},1i_1\mathrm{}i_kn.$$
(2)
Another choice is the symmetrizer map,
$$\text{Sym}(x_{i_1}x_{i_2}\mathrm{}x_{i_k})=\frac{1}{k!}\underset{\sigma S_k}{}X_{\sigma (i_1)}X_{\sigma (i_2)}\mathrm{}X_{\sigma (i_k)},$$
(3)
where $`S_k`$ is the group of permutations of order $`k`$.
Given a choice for $`\psi `$ one can define an associative product (star product) on $`𝒜[[h]]`$ by
$$a_\psi b=\psi ^1(\psi (a)\psi (b)).$$
Then, for any choice of $`\psi `$, $`(𝒜_h,_\psi )`$ is an algebra isomorphic to $`𝒜_h`$. With the star product we recover the semiclassical interpretation of the elements of the algebra as functions on the phase space. The star product can always be written as a formal series
$$a_\psi b=ab+\underset{n>0}{}h^nC_\psi ^n(a,b),$$
where $`C_\psi ^n`$ are some bilinear operators. Let $``$ and $`^{}`$ be two isomorphic star products,
$$ab=T^1(T(a)T(b)),T:𝒜[[h]]A[[h]].$$
It is clear that $`T`$ can be written as
$$T(a)=\underset{n0}{}h^nT^n(a),$$
and because of property a, $`T^0`$ must be an automorphism of the commutative algebra $`𝒜[[h]]`$. If $`T_0`$ is the identity we say that $``$ and $`^{}`$ are equivalent (or gauge equivalent) star products. (2) and (3) are two equivalent star products.
For $`𝒜`$ being the full algebra of $`C^{\mathrm{}}`$ functions on the Poisson manifold, if the operators $`C_\psi ^n`$ are bidifferential operators we say that the star product is differentiable. Gauge equivalence can be restricted to the class of differentiable star products by considering only differentiable $`T^n`$. Notice that the differentiability is a property of the particular star product and not of the formal deformation. The star products (2) and (3) can be extended to $`C^{\mathrm{}}(\mathrm{\Theta })`$ as differentiable star products, but we will see later an example of a star product corresponding to the same formal deformation which is not differentiable .
We will consider only semisimple Lie groups. The semisimple coadjoint orbits of $`G`$ on $`𝒢`$ are complex algebraic varieties defined over $``$. They are given by the invariant polynomials. If $`l`$ is the rank of $`G`$, we can choose $`l`$ homogeneous polynomials $`p_i(x)`$, $`i=1,\mathrm{},l`$ generating the subalgebra of invariant polynomials, $`[p_1,\mathrm{}p_n]`$. Then the semisimple coadjoint orbits are given by the algebraic equations
$$p_i(x)=c_i.$$
(4)
(see for example ). The intersection of the complex orbit with $`𝒢_R`$ is a real algebraic variety consisting on a finite number of connected components, which are orbits of the real form of the group. For the compact real form there is only one connected component.
It is easy to check that the star products 2 and 3 do not restrict well to the orbit, that is, in general
$$ap_i|_\mathrm{\Theta }0.$$
We want to know if there is some choice of $`\psi `$ that gives a star product which restricts to $`\mathrm{\Theta }`$.
In the approach of geometric quantization, the algebra of quantum observables is given by the quotient of $`U_h`$ by a certain ideal. This ideal, $`I_h`$, is prime and $`\text{Ad}_G`$-invariant, so there is a well defined action of $`G`$ on $`U_h/I_h`$. We summarize here the results of , where the quantization of the coadjoint orbits are obtained in terms of the quotient of the enveloping algebra by a prime, Ad<sub>G</sub>-invariant ideal. We consider the polynomial algebra over the real algebraic manifold (union of orbits) defined by (4),
$$\text{Pol}(\mathrm{\Theta })=\text{Pol}(𝒢^{})/_0,_0=\{p\text{Pol}(𝒢^{})/p|_\mathrm{\Theta }=0\}.$$
This is the Poisson algebra that we want to deform. We quote first a result from Varadarajan that we need.
Lemma (1). Let $`x𝒢^{}`$ be a regular element of $`𝒢^{}`$ (or equivalently, a point in which the centralizer has dimension equal to the rank of $`𝒢^{}`$). Then $`(dp_1)_x`$, …, $`(dp_l)_x`$ are linearly independent.
From now on we will restrict to regular orbits only. It is clear that in this case $`_0`$ is generated by $`p_1c_1,\mathrm{}p_lc_l`$. We consider now the elements in $`U_h`$ that are the image of $`p_i`$ by the symmetrizer, $`P_i=\text{Sym}(p_i)`$, called Casimir operators. $`P_i`$ are central elements in $`U_h`$ and they are also Ad<sub>G</sub>-invariant (the symmetrizer commutes with the action of $`G`$).
Let $`_h`$ be the ideal generated by $`P_iC_i(h)`$, where $`C_i(0)=c_i`$. Then $`U_h/_h`$ is a formal deformation of $`\text{Pol}(\mathrm{\Theta })`$. The technical assumption of regularity is needed to prove the existence of a $`[[h]]`$-module isomorphism $`\psi :\text{Pol}(\mathrm{\Theta })[[h]]U_h/_h`$, which is not obvious. The ideal $`_h`$ itself is Ad<sub>G</sub>-invariant, so $`G`$ has a natural action by automorphisms on the algebra $`U_h/_h`$. For special values of $`C_i(h)`$, $`_h`$ is in the kernel of an irreducible unitary representation of $`G_R`$.
This deformation of polynomials can be specialized for any value of $`h`$. For SU(2),
$$[H,X]=\mathrm{}2X,[H,Y]=\mathrm{}2Y,[X,Y]=\mathrm{}H,$$
the Casimir operator is
$$P=\frac{1}{2}(XY+YX+\frac{1}{2}H^2).$$
It was shown in that with the choice
$$C=l(l+\mathrm{})),l=\mathrm{}m/2,$$
the algebra obtained is the same than the one obtained in geometric quantization.
According to our definition, a star product in Pol$`(\mathrm{\Theta })`$ is given by a $`[[h]]`$-module isomorphism $`\stackrel{~}{\psi }:\text{Pol}(\mathrm{\Theta })[[h]]U_h/_h`$. In particular, to obtain a star product in Pol$`(𝒢^{})`$ which restricts to the orbit one should look for an isomorphism $`\psi :\text{Pol}(𝒢^{})[[h]]U_h`$ such that the following diagram commutes
$$\begin{array}{ccc}\text{Pol}(𝒢^{})[[h]]& \stackrel{\psi }{}& U_h\\ \pi & & \widehat{\pi }& & \\ \text{Pol}(\mathrm{\Theta })[[h]]& \stackrel{\stackrel{~}{\psi }}{}& U_h/I_h.\end{array}$$
(5)
$`\pi `$ and $`\widehat{\pi }`$ are the natural projections. If $`\psi (_0)=_h`$, then $`\stackrel{~}{\psi }`$ is defined uniquely by the diagram (5). An example of such star product is given in , where it is also shown that it is not differentiable. Since $`\widehat{\psi }`$ is not unique one could ask if there is some choice that renders it differentiable. This issue will be addressed in . |
warning/0003/astro-ph0003204.html | ar5iv | text | # Structure of the Large Magellanic Cloud from 2MASS
## 1 Introduction
Morphologically, the LMC is an irregular barred spiral galaxy with three spiral arms and an extended outer loop of stellar material de Vaucouleurs & Freeman (1973). Based on deprojection and photometric distances (e.g. de Vaucouleurs 1957, 1980), its disk is inclined at an angle of $`27^{}`$ to the plane of the sky. The disk exhibits solid body rotation out to $`2.5^{}`$ with a rotation center at $`5^h21^m,69^{}17^{}(1950)`$, about $`0.6^{}`$ north of the optical center of the bar. This kinematic signature is present in a variety of tracers: HI gas, planetary nebulae, HII regions, supergiants, CH stars, etc. Freeman et al. (1983) have examined kinematics of rich star clusters with ages between $`100`$ Myr and $`10`$ Gyr. They found that young clusters rotated with HI gas, while the older ones (SWB VII; Searle et al. 1980) formed a flattened rotating system with dispersion along the line of sight $`\sigma 18`$ km/s. A later study of more extended sample of outer LMC clusters Schommer et al. (1992) confirmed the absence of isothermal pressure-supported spheroid. All this has led to the standard view that the LMC is a geometrically thin object.
However, recent studies have suggested that the LMC may have an extended component. First, the evidence for a flattened spheroid population was found in the kinematics of old long-period variables Hughes et al. (1991). Kunkel et al. (1997) describe a population of carbon stars out to 12 kpc from the LMC center. These authors interpret these in the context of a thin disk model and derive a rotation curve and mass estimate. However, Weinberg (2000) argues that the LMC should be evolving rapidly in the Milky Way tidal field, based on both analytic calculations and n-body simulations. The tidal field causes the LMC disk axis to precess and torques disk orbits out of the disk plane, causing a strongly flared, spheroidal-like distribution in the outer Cloud and loss of stars and gas. This interaction leads to an a spatially extended population while roughly preserving the disk-like kinematic signature (i.e., small $`\sigma `$).
The detection of, or strict limits on, the predicted extended distribution would resolve these views and is one of the goals for the present study. Our star count analysis is based on fitting the projected spatial density of several LMC populations among those identified in Nikolaev & Weinberg (2000; hereafter NW00) based on their location in the color-magnitude diagram (CMD) of the field. The late-type giant populations are dominated by the LMC and may be used as tracers of the spatial structure of the Cloud. Each population is fitted by two models: 1) thin exponential disk and 2) spherical power law model. Our best-fit models based on circular disks give the inclination of the LMC to the line of sight $`i24^{}28^{}`$ and the position angle of the disk $`\theta 165^{}174^{}`$, in good agreement with previous estimates. The direction of the LMC disk inclination is also determined. In short, projected 2MASS star counts reproduce the standard LMC inferences.
The near-infrared 2MASS photometry easily discriminates carbon stars in the color-magnitude diagram. While the 2MASS single-epoch survey of the LMC does not provide variability information, we identified a region of the CMD populated nearly exclusively by carbon-rich AGB stars (Region J in NW00), which are also long-period variables (LPVs). This identification is reinforced by recent analyses of MACHO data Alves et al. (1998); Wood (1999), although roughly $`25\%`$ could be binaries Wood (1999). LPVs obey period-luminosity-color (PLC) relations (e.g., Feast et al. 1989), and based on the PLC relations, the LPVs in a narrow color range ($`1.6<JK_S<1.7`$) are standard candles with $`\sigma _K0.3^m`$. Their photometric distribution in selected LMC fields has a least three distinct components: a well-defined narrow distribution due to LMC disk and two secondary peaks at fainter and brighter magnitudes. The differential photometric distance to the central disk peak provides a direct determination of the inclination: $`42.3^{}\pm 7.2^{}`$. This value is consistent with but larger than those inferred by deprojecting isopleths. The secondary peak could in principle be due to stellar blends, geometric structure, interstellar reddening, distribution of periods, gradient in age and metallicity, or contaminating population of overtone pulsators in the sample. We carefully examine the possible origins of this secondary component, and conclude that most likely explanation is spatial structure. This interpretation is bolstered by the good match of the central peak, which includes known fundamental mode pulsators, with the established LMC disk inclination. This implies the presence of an extended component of the LMC, which may be as thick as 14 kpc along the line of sight. This component is thicker than the 2.8 kpc flattened spheroid suggested by Hughes et al. (1991) from kinematic data and may be streams of material rather than be smooth and well-mixed. The bright peak is then an intervening population at a distance of roughly 35 kpc and the existence of relatively young carbon-rich AGB stars suggests tidal debris.
The outline for the paper follows. In the main section, we briefly describe the 2MASS sample (§2), and present the star count (§3) and the standard candle (§4) analyses of the LMC. In particular, §4.1 describes the selection of standard candles, §4.2 gives the details of the standard candle analysis and §4.3 presents the main results. The major conclusions are summarized in §5.
## 2 Observations
The LMC field, (4<sup>h</sup>00<sup>m</sup> to 6<sup>h</sup> 56<sup>m</sup> in right ascension, $`78^{}`$ to $`60^{}`$ in declination, J2000.0) has been observed by 2MASS and is included in the most recent data release. Details of the data reduction and sample selection are described in NW00. Figure 1 shows both the projected spatial distribution of sources and the color-magnitude diagram of the field.
The color-magnitude diagram also shows the location of 12 regions analyzed in NW00. Total number of sources in our sample is 823,037.
NW00 examined the populations in selected CMD regions and associated the features with known populations of stars (cf. Table 2 of NW00). Here, we take an in-depth look at the spatial distribution of sources in seven of those regions (Table 1) which contain mostly LMC stars. These seven areas of the CMD account for a quarter of all sources in the field.
The projected spatial density distribution for each of the seven regions is shown in Figure 2. As expected, younger populations (Regions A and H) have relatively clumpy distributions which trace the spiral pattern of the LMC (cf. Figure 7 of Schmidt-Kaler 1977). Older stars, on the other hand, have smoother and more extended distributions with significant overdensity in the bar of the Cloud. Several well-known morphological features of the LMC are easily recognizable in Figure 2, e.g. 30 Doradus complex (an HII region near $`\alpha =5^h36^m`$, $`\delta =69^{}`$) and asymmetric outer loop in the south-eastern part of the LMC, traced by AGB stars.
## 3 Spatial Structure Using Parametric Maximum Likelihood
To quantify the spatial distribution of sources in six selected regions, we perform maximum likelihood (ML) analyses for thin exponential disk and spherical power-law models. The observed source counts for each population are binned in equatorial coordinates, $`n_{ij}^o,i=1,N;j=1,M`$. The scale was chosen by successively reducing the bin size until the inferred parameters converged. The final coordinate grid is spaced uniformly, with the bin size $`\mathrm{\Delta }\alpha =\mathrm{\Delta }\delta =1^{}`$. The ML scheme selects a parametric model for which the expected source counts $`n_{ij}^e`$ most closely match the observed source counts $`n_{ij}^o`$. The goodness-of-fit measure,
$$u^2=\underset{i,j}{}\frac{(n_{ij}^on_{ij}^e)^2}{n_{ij}^e},$$
is distributed as $`\chi ^2`$ with $`(N\times M1n_p)`$ degrees of freedom, where $`n_p`$ is the number of free parameters of the model (see below). The expected source counts $`n_{ij}^e`$ are obtained from the corresponding source density in each bin, $`\rho _{ij}^e`$, predicted by the model:
$$n_{ij}^e=\frac{N_{src}\rho _{ij}^e}{_{i,j}\rho _{ij}^e}$$
(see Appendix A for details).
The density of the exponential disk is described by
$$\rho =C\mathrm{exp}(r/R),$$
(1)
where $`R`$ is the disk scale length. The model has five free parameters: coordinates of the disk center $`\alpha _0`$, $`\delta _0`$, the scale length $`R`$, and two orientation angles: inclination $`i`$ ($`0^{}i<180^{}`$) and position angle (line of nodes) $`\theta `$ ($`0^{}\theta <360^{}`$). The position angle $`\theta `$ increases counterclockwise from the north (i.e. NESW). The density of the spherical power law model is described by
$$\rho =Sr^\nu .$$
(2)
The free parameters are positions of the LMC center, $`\alpha _0`$ and $`\delta _0`$, and power law index $`\nu `$.
The results of parametric fits are presented in Table 2 (for thin exponential disk model) and Table 3 (for spherical power law model). The best fit parameters in both tables can be grouped by the population age: results for young stars (A, H) and older stars (E, F, J, K) show good agreement among themselves. The error in centroid of each population is $`\text{ }<\text{ }0.1^{}`$. The near side position angles marked with ‘?’ are inaccurate due to clumpiness and compactness of the corresponding distributions.
The distributions of young OB stars and supergiants are clumpy and therefore poorly fitted by a smooth model. The centroid of these populations is $`1^{}`$ to the north of the optical center of the bar (defined by the center of symmetry of the bar, at $`\alpha _{1950}=81.0^{}`$, $`\delta _{1950}=69.8^{}`$), similar to the displacement found by de Vaucouleurs & Freeman (1973). The scale lengths $`R`$ derived for these populations are noticeably greater than the scale lengths for older populations and reflect the location of the distinct star forming activity. The position and inclination angles for these populations are mutually consistent. Our derived inclinations, $`i36^{}38^{}`$, are consistent with $`i=38.2^{}`$ found from the distribution of HI regions (Feitzinger et al. 1977), and also with $`i=36_5^{+2}`$ degrees found from Monte Carlo simulations of ultraviolet photopolarimetric maps of the western LMC Cole et al. (1999).
The older populations (M giants, AGB stars and LPVs) are well-represented by a smooth density law. The centroids for these populations are within $`0.4^{}`$ of each other on the sky and are close to the optical center of the bar. The scale lengths are $`R1.31.4`$ kpc (sources in Region E have the largest scale length among older stars, $`R1.6`$, which is probably due to Galactic M dwarfs contamination). The inferred inclinations are in the range $`i24^{}28^{}`$, in good agreement with previous determinations from star counts de Vaucouleurs (1955), $`i=(25\pm 5)^{}`$; distribution of star clusters, $`i=(25\pm 9)^{}`$ or HI isophotes, $`i=(27\pm 5)^{}`$ McGee & Milton (1966); or photographic R isophotes de Vaucouleurs (1957), $`i=(27\pm 2)^{}`$. These deprojection-based position angles for the group are mutually consistent, $`\theta 169^{}173^{}`$, and fall in the range $`\theta =160^{}180^{}`$ derived from surface photometry of the LMC by others (see, e.g. Table 2 of Schmidt-Kaler & Gochermann 1992).
The extended spatial coverage of the LMC field by 2MASS allows one to determine the absolute direction of the inclination, i.e. to determine the closest side of the LMC. To illustrate this point, we make two different test models of an exponential disk with $`\{\alpha _0,\delta _0,R,i,\theta \}`$ $`=`$ $`\{80^{},70^{},1.5,\pm 30^{},135^{}\}`$. Both disks are modeled with 3,000 point sources. The restored density contours are presented in Figure 3. The difference in the expected source counts is clearly seen in the outer regions. This suggests that for relatively spatially extended populations both the absolute value and the direction of the inclination can be reliably determined. On the other hand, if a population is relatively compact in the sky, the inferred direction of inclination may differ from the actual value (cf. Table 2). Based on the results for older populations, we see that the nearest side of the LMC is its eastern side, in agreement with previous results based on photometry of Cepheids in the Cloud de Vaucouleurs (1955); Gascoigne & Shobbrook (1978); Laney & Stobie (1986). Restricting our attention to Region J, which has little if any Galactic contamination, we can be sure that we are not affected by a gradient caused by the Galactic foreground.
The results of fits to a spherical power-law profile are described in Table 3. There are no significant difference between power-law exponents for various populations: all values are $`2.52.6`$. The centroid shift for younger populations, present in Table 2, is seen here as well. The underlying distribution, a disk and spheroid together, is described here as a single profile. A power-law disk profile, of course, will have an index $`1+\nu `$, where $`\nu `$ is the index for the spherical profile, and therefore some unknown combination of a disk and spheroid makes interpretation difficult. We note that these fits have poorer quality than the exponential fits and probably do not bear on the reality of a spheroidal population. A more independent test for spatially extended population is described below (§4.2).
We compared the results of our fits to similar fits by Hughes et al. (1991), who modeled the distribution of intermediate and old long-period variables (ILPV and OLPV, respectively) with exponential disk and power law models. They found $`\nu =1.8\pm 0.1`$ and $`R=1.6\pm 0.2`$ (OLPV), $`\nu =1.7\pm 0.1`$ and $`R=1.7\pm 0.2`$ (ILPV). Our results for similar populations (Regions E, F, G and J) imply $`\nu 2.5`$ and $`R1.4`$ kpc. While the scale lengths $`R`$ are mutually consistent, the best-fit power law exponents differ. Our best fit parameters are larger than the power-law model exponents of Hughes et al. which has $`\nu 2.0`$.
To summarize, the projected distribution of LMC populations observed by 2MASS is consistent with previous studies. We found the scale length of the LMC disk, $`R1.41.5`$ kpc, the inclination angle, $`i24^{}28^{}`$, and the direction of the LMC tilt in good agreement with existing estimates.
## 4 Standard Candle Analysis
We complement our previous analysis (§3) by incorporating photometric distances in addition to our CMD selection. Below, we describe the selection of standard candles, the details of the analysis, and the implications for the structure of the LMC derived from 2MASS photometry.
### 4.1 Selecting Standard Candles from 2MASS Data
Good standard candles satisfy three conditions: 1) they must be luminous and easily identified; 2) they must be sufficiently numerous and representative of the underlying structure, and 3) they must have small photometric dispersion as a class (that is, small $`\sigma _M`$). In NW00, we argued that stars in Region J of the 2MASS CMD are potentially good standard candles. Being brighter and redder than the RGB tip, most of these stars are carbon-rich thermally-pulsating AGB stars (TP-AGB). Recent data Alves et al. (1998); Wood (1999) suggests that most of these stars exhibit Mira-like variability. The fraction of variables in this region is close to $`100\%`$, although roughly $`25\%`$ could be binaries Wood (1999). Here, we will assume that all sources in Region J are carbon-rich long-period variables. Their red colors effectively discriminate against the population of oxygen-rich LPVs, since the latter rarely have $`JK_s>1.5`$ Hughes & Wood (1990). As long-period variables, these stars follow a linear period-luminosity-color relation (e.g. Feast et al. 1989). The luminosity of these stars can be characterized by their periods or near-infrared colors, and therefore, these stars can be used to probe the structure of the LMC along the line of sight.
The 2MASS sample of C-rich LPVs in the LMC contains 8229 stars. The surface map of Region J is shown in Figure 4.
The luminosity-color relation for these stars results in the well-defined ridge in the figure. In the absence of the period data, we cannot use standard PL or PC relations to calibrate the intrinsic brightness of these variables. Rather, we have to rely on luminosity-color relation. Figure 5 shows the sample of 79 oxygen- and carbon-rich Miras in the LMC Glass et al. (1990). Magnitude and color of each Mira are averaged over period and plotted on top of 2MASS color-magnitude digram. The luminosity-color relation (shown with the solid line) for 14 carbon Miras in the color interval bounded by vertical dashed lines is
$$<K_s>=(0.99\pm 0.80)<JK_s>+(12.36\pm 1.33),\sigma =0.38.$$
The average r.m.s is $`\sigma \text{ }<\text{ }0.3^m`$ for $`1.4<JK_s<1.7`$. Given this LC relation, we may argue that selecting LPVs from a reasonably narrow color range will result in sources with similar luminosities, i.e., standard candles. For our analysis, we choose color interval $`1.6<JK_s<1.7`$, sufficiently narrow to ensure similar luminosities of our standard candles and sufficiently broad to host enough sources (1385) for statistically meaningful inference.
The ridge-line fit above was based on the average LC relation. From analysis of light curves in Glass et al. (1990), the amplitudes of carbon-rich LPVs are $`\mathrm{\Delta }K\text{ }<\text{ }0.5`$ mag, so one may expect significant broadening of the LC relation due to random phase observations. However, as we will demonstrate below (§4.3), the width of main peak in the apparent luminosity function is only $`\sigma 0.2^m`$, even with random phase observations. This suggests that the phase-average LC relation is much tighter.
### 4.2 Method
Without prior knowledge of the true source distribution, which is likely to be irregular due to tidal interaction (e.g. Weinberg 2000) or sufficient characterization of the stellar populations to allow a non-parametric density estimation with all of the data, we study the photometric distribution in several fields in the LMC. The fields are located along two great circle arcs passing through the central region of the Cloud: Arc 1, parallel to the line of nodes, and Arc 2, perpendicular to the line of nodes (see Fig. 6).
The apparent luminosity function can be analyzed in terms of the centroid $`\overline{m}`$ and the width $`2\sigma _m`$ of the distribution. For a homogeneous standard candle population, the centroid of the distribution measures LMC distance and the width of the distribution gives an estimate of its line-of-sight depth. The error analysis of a standard apparent magnitude-absolute magnitude-distance relation gives
$$\sigma _m^2\sigma _M^2+\frac{4.72}{R_{LMC}^2}\sigma _r^2+\sigma _A^2+\sigma _{ph}^2,$$
(3)
where $`R_{LMC}=50`$ kpc is the average distance to the LMC, $`\sigma _M`$ is the intrinsic precision of our standard candles, $`\sigma _A`$ is the variance due to extinction, $`\sigma _r`$ is the geometric depth, and $`\sigma _{ph}`$ is the photometric error. Equation (3) states that the apparent brightness distribution is the convolution of the spatial density and the intrinsic luminosity function and therefore provides an upper limit to the geometrical depth.
### 4.3 Results and Interpretation
The brightness distributions of standard candles in selected fields are shown in Figures 7 and 8.
All fields in Figures 7 and 8 exhibit well-defined central peaks corresponding to the midplane of the LMC disk. We immediately notice that the narrowest features in the brightness distributions have widths $`\sigma _m0.2^m`$ (cf. Field 3). This is a direct evidence that our color-selected sources are standard candles at least as good as $`\sigma _M0.2^m`$. In fact, they are even better, because of the additional terms on the right-hand-side in equation (3). This suggests that carbon LPVs in the narrow ($`0.1`$ mag) color range are excellent standard candles, even observed at random phases.
#### 4.3.1 Analysis of the Distribution Centroids
The centroids of the distributions are consistent with the inclination of the LMC derived in §3. Since Arc 1 is parallel to the line of nodes, the stars in the fields along this arc should be roughly at the same distance. Hence, we do not expect any drift in the means $`\overline{m}`$ for stars in these fields. On the other hand, we expect a shift in the mean magnitude for fields along Arc 2, which is perpendicular to the line of nodes. Assuming the Eastern part of the Cloud is closer to us, sources in Eastern fields should be, on average, brighter than their counterparts in Western fields. Figures 7 and 8 confirm the expectations. The magnitude of the effect is shown in Figure 9 for both arcs as the function of the angular distance from the optical center of the bar.
To improve the signal-to-noise ratio, we take the average $`\overline{m}`$ of the mean magnitudes in three bands and plot the resulting averages normalized to the central field of the respective arc. We find that fields along the North-South line have similar mean magnitude, whereas Eastern fields are on average $`0.1^m`$ brighter and Western fields are $`0.2^m`$ fainter than central bar fields. The magnitude difference is similar to what Caldwell & Coulson (1986) found in their analysis of Cepheids in the Cloud.
Data shown in Figure 9 allow independent determination of the LMC inclination angle. We provide results of two methods: (1) the inclination of the best fit plane to the photometric distance estimate for each field; and (2) the mean inclination of each field paired with the central field. In Method 1, we compute the best fit line through spatial positions of the field centers. The estimates for each arc are $`0.9^{}\pm 0.3^{}`$ and $`42.3^{}\pm 7.2^{}`$, respectively. We observe that the inclinations along our Arc 1, as expected, are consistent with zero, while the inclination angles along Arc 2 are consistent with values in Table 2 at $`2\sigma `$ level. Alternatively, in Method 2 does not presuppose a tilted plane but rather computes the inclination of each field relative to the line of sight through the central field. We then compute the weighted average over all pairs. The resulting inclinations are $`1.0^{}\pm 4.1^{}`$ and $`34.4^{}\pm 2.7^{}`$ for the first and second arcs, respectively.
In principle, this magnitude drift could have been produced by a reddening gradient across the Cloud. However, the extinction maps by Oestreicher & Schmidt-Kaler (1996) do not show any systematic change in reddening in the West-East direction across the LMC. In addition, an analysis of the extinction across the entire cloud based on the location of the giant branch reveals little evidence for significant extinction on average outside of the inner square degree Nikolaev & Weinberg (2000). Moreover, if the magnitude drift were indeed due to reddening gradient, it would have had band-dependent signature. The ratios of magnitude offset in $`J`$, $`H`$, $`K_s`$ bands would be in direct proportion to the extinction coefficients $`A_J`$, $`A_H`$, $`A_K`$. Since this is not the case, we have to reject this possibility and interpret the drift as the true distance effect.
#### 4.3.2 Shape of the Distributions
Careful study of the density distributions in Figures 7 and 8 shows that some distributions have extended tails. For example, distributions in Fields 3 and 4 clearly show positive skewness. Figure 4 also shows the extension of the central ridge toward fainter magnitudes. To quantify the extended component, we model the distribution of LMC disk sources with a properly centered Gaussian. After subtracting the disk component in the central field (Field 3), we find a secondary, broader distribution which contains approximately 90 sources or about a quarter of the total number in the field. This distribution itself appears to have two distinct components. The strongest peak of this extended component is $`0.5`$ mag fainter than that of the disk distribution and a weaker peak is roughly $`0.8`$ mag brighter. In general, the tail toward fainter magnitudes is more pronounced, although the distribution extends to both sides.
Among the factors which could produce the extended component visible in Figures 7 and 8 are: 1) foreground population, 2) source blending, 3) interstellar reddening, 4) population of overtone pulsators, 5) distribution of periods, 6) age/metallicity variations, and 7) spatial density distribution<sup>1</sup><sup>1</sup>1We do not list spurious detection as the possibility, because the distribution s are the same in three bands, suggesting real feature. Foreground population can be rejected because the fraction of Galactic sources in Region J is nearly zero (see Table 1). Interstellar reddening would produce a band-dependent effect (i.e. longer tails in $`J`$ band than in $`K_s`$ band). We consider each of the remaining five possibilities in detail:
1. Source blending. Considering that the secondary peak runs parallel to the main LC relation (cf. Figure 4), at least one member in the blend must be a carbon star. However, blending a carbon star with an unresolved source would produce a source brighter than the carbon star. This means the secondary peak would be brighter than the main feature, which contradicts to data;
2. Population of overtone pulsators. The question about the pulsation mode of Mira variables seem to have been resolved recently with Miras unambiguously identified as fundamental mode pulsators Wood & Sebo (1996); Wood et al. (1998). It is conceivable, then, that our color-selected sample of standard candles includes a population of first- and higher-overtone LPVs. Based on Figures 7 and 8, this ‘secondary’ population should constitute about $`25\%`$ of the entire sample, and be fainter by $`0.40.5^m`$. We test this possibility by cross-correlating 2MASS data with the sample of Hughes & Wood (1990), which contains fundamental-mode LPVs only. Picking out sources in the color range $`1.6<JK_s<1.7`$ and plotting the histogram of their $`K_s`$ magnitudes reveals the same extended component, approximately $`0.5`$ mag fainter than the main peak. Since the feature is present even in the data without overtone pulsators, this explanation has to be rejected;
3. Distribution of periods. Even in case of zero-dispersion PLC relation for C-rich LPVs, the observed magnitude distribution could result from an intrinsic period distribution (in addition to overtone pulsations discussed above). Although this explanation remains a possibility because the physics of these stars remains uncertain in detail, such an effect is not apparent in current theoretical models, which confirms an intrinsic width of the instability strip but not multiple peaks. Conversely, as previously mentioned, the good fit of the primary peak to the disk inclination gives us confidence in a well-defined instability strip and width of this peak is consistent with other pulsators. Alternatively, the pulsation periods depend on mass and metallicity and distribution of periods implies the range of mass, ages and metallicities, and vice versa. Assuming a range of initial AGB masses leads to the range of intrinsic bolometric magnitudes for a single period Marigo et al. (1996). Marigo’s et al. theoretical evolutionary tracks show the difference in a few tenths of a magnitude for carbon Miras at a constant period, depending on the mass of the star. Fitting the apparent luminosity function to these tracks Cole (2000) suggests that the stronger peak may be due to younger ($`0.6`$ Gyr) and more massive stars ($`2M_{}`$), while the broader component is formed by older ($`2.8`$ Gyr) and less massive ($`1.2M_{}`$) stars. This appears consistent with observations: Frogel et al. (1990) found that carbon stars are present in globular clusters aged 100 Myr to about 3 Gyr and that carbon stars in intermediate-age LMC clusters are a few tenths of a magnitude fainter than those in young LMC clusters. Alves et al. (1998) reported that AGB variables in the intermediate age cluster NGC 1783 are about 0.5 mag brighter at given period than variables in ancient LMC globular cluster NGC 1898. We find two aspects of this possibility to be unsatisfying. First, it does not naturally explain the bright secondary peak. Second, the star formation rate subsequent to the 2.8 Gyr burst would have to be small in order to result in a well-defined feature;
4. Extended density distribution in the LMC. This appears to be the most natural explanation to the observed brightness distributions: it is band-independent, it does not rely on special star formation history, and straightforwardly explains both the brighter and fainter secondary distribution. The distinct feature at fainter magnitudes may be associated with the population kinematically distinct from the disk sources, recently found by Graff et al. (1999). This motivates a detailed kinematic follow-up of the small bright “clump” in the photometric distribution. Based on the offset, this population is 15 kpc from the LMC center and coincident with the intervening population detected by Zaritsky & Lin (1997). The existence of LPVs implies a relatively young population, and the broad area is consistent with tidally stripped material. A future detailed study will attempt to self-consistently model the disk, an extended spheroid component as suggested by Hughes et al. (1991), and test for other distinct features. Regardless of details, the broad photometric distribution suggests that LMC is geometrically thick along the line of sight. From Figures 7 and 8, the average thickness of the LMC is $`2\sigma _m0.7^m`$, which implies a thickness (after deconvolution) of $`14`$ kpc (for the LMC distance $`R_{LMC}=50`$ kpc). This makes the LMC as thick along the line of sight as it is wide across the sky, something which one would naively expect for a dwarf companion in the process of being tidally shredded. In summary, the spatial interpretation of these photometric distributions suggests the LMC consists of a centrally concentrated barred disk and some number of extended distributions, including a spheroid and tidally distorted and possibly stripped populations.
## 5 Summary
We have analyzed the spatial distributions of several LMC populations, identified in NW00 based on their location in the color-magnitude diagram. Quantitative analysis of the observed distributions includes parametric source-count and a preliminary standard candle analyses. Our major conclusions are:
* Projected star count analyses based 2MASS data yield scale lengths, deprojection-based inclinations, and position angles consistent with previous studies. The near-far degeneracy of the disk orientation is broken by perspective difference. The near side of the Cloud subtends a larger angle in the sky-and this allows us to determine that Eastern side of the disk is closer, in agreement with Cepheid-based results.
* We propose using carbon-rich LPVs in a narrow color range, $`1.6<JK_s<1.7`$ as standard candles to probe the structure of the LMC along the line of sight. Based on published light curves, their intrinsic magnitudes have a dispersion of $`\text{ }<\text{ }0.3^m`$ including the random phase of the observations. The width of the LMC disk in the observed photometric distribution ($`\sigma _m0.2`$) suggests that these are excellent standard candles.
* The photometric distribution of our standard candle sample reveals strong central peak with extended tails in both directions, with tail toward fainter magnitudes more pronounced. In the densest field, the distributions in all bands show two secondary peaks, $`0.5`$ mag fainter and $`0.8`$ mag brighter than the primary. We interpret the primary peak as due to the midplane of the LMC disk and examine various possibilities which could produce the tails and secondary peaks of observed distribution. We conclude that the photometric distribution of standard candles is most likely caused by a spatially extended stellar component and tidally stripped debris, consistent with a tidally disturbed dwarf companion. It is possible that the distribution reflects the intrinsic period distribution of LPVs but or distinct populations of masses (ages) and metallicities for these stars, but such explanations special populations or evolutionary histories.
* The distribution of apparent magnitudes of standard candles is consistent with tilted geometry derived from star counts. We derive a direct determination of the LMC disk inclination of $`42.3^{}\pm 7.2^{}`$, consistent with Laney & Stobie (1986) estimate of $`45^{}\pm 7^{}`$ and with results of Welch et al. (1987), $`37^{}\pm 16^{}`$. Interpreting the apparent luminosity function as due to real source density, we find evidence of the extended component of the LMC, with a width of approximately 14 kpc (for the LMC distance $`R_{LMC}=50`$ kpc).
This detection of this thick component implies that LMC may contain a kinematically distinct population as suggested by Graff et al. (1999) and/or an the extended component found by Hughes et al. (1991) and predicted for the tidal interaction with the Milky Way (Weinberg 2000). The existence of such a population will affect the self-lensing models of the LMC in the microlensing studies. These preliminary results suggest a number of projects for short- and long-term follow-up. An improved analysis of the 2MASS sample may be obtained with larger sample of standard candles, or with full three dimensional analysis based on several distinct standard candles using data from the entire survey rather than selected fields (work in progress). Assuming that the extended stars originated in the disk, both populations will appear rotationally supported. However, using the 2MASS photometric distribution as a guide, a combined disk and extended spheroid population should be kinematically separable and these stars are good candidates for future spectroscopic work.
## Acknowledgements
The authors would like to thank Peter Wood for MACHO data on variable stars and for help in interpretation of these data. We are grateful to Michael Skrutskie and Andrew Cole for careful reading of the manuscript and suggesting ways of improving it. We thank Shashi Kanbur for helpful advice regarding LPVs. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center, funded by the National Aeronautics and Space Administration and the National Science Foundation.
## Appendix A Parametric Maximum Likelihood
The expected source density for a bin in the direction ($`\alpha _i`$, $`\delta _j`$) is determined by integrating the LMC density model along the line of sight across the bin. This is given by the following integral:
$`\rho _{ij}^{exp}`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑tt^2{\displaystyle _{\delta _j\mathrm{\Delta }\delta /2}^{\delta _j+\mathrm{\Delta }\delta /2}}𝑑\delta \mathrm{cos}\delta {\displaystyle _{\alpha _i\mathrm{\Delta }\alpha /2}^{\alpha _i+\mathrm{\Delta }\alpha /2}}𝑑\alpha \rho (t,\alpha ,\delta )`$ (A1)
$``$ $`{\displaystyle _0^{\mathrm{}}}𝑑tt^2\rho (t,\alpha _i,\delta _j)\mathrm{cos}\delta _j\mathrm{\Delta }\alpha \mathrm{\Delta }\delta ,`$
where the last equality assumes that the bin size $`\mathrm{\Delta }\alpha `$, $`\mathrm{\Delta }\delta `$ is sufficiently small. We perform the integral in (A1) using 256-point Gaussian quadrature formula, with $`20`$ kpc and $`80`$ kpc as the integration limits. The underlying source density $`\rho (t,\alpha ,\delta )`$ is given either by equation (1) or by equation (2). The coordinate transformations for both models follow.
### A.1 Exponential Disk
To quantify $`\rho ()`$, we introduce the coordinate system $`\{x_0,y_0,z_0\}`$ which has the origin at the center of the LMC at $`\{t,\alpha ,\delta \}=\{R_{LMC},\alpha _0,\delta _0\}`$ and has $`z_0`$-axis toward the observer, $`x_0`$-axis antiparallel to the right ascension axis, and $`y_0`$-axis parallel to the declination axis. The coordinate transformations are given by
$`x_0`$ $`=`$ $`t\mathrm{cos}\delta \mathrm{sin}(\alpha \alpha _0)`$
$`y_0`$ $`=`$ $`t\mathrm{sin}\delta \mathrm{cos}\delta _0t\mathrm{cos}\delta \mathrm{sin}\delta _0\mathrm{cos}(\alpha \alpha _0)`$ (A2)
$`z_0`$ $`=`$ $`R_{LMC}t\mathrm{cos}\delta \mathrm{cos}\delta _0\mathrm{cos}(\alpha \alpha _0)t\mathrm{sin}\delta \mathrm{sin}\delta _0.`$
The coordinate system of the exponential disk, $`\{x^{},y^{},z^{}\}`$, is the same rectangular system as $`\{x_0,y_0,z_0\}`$, except rotated about $`z_0`$-axis by the position angle $`\theta `$ counterclockwise and about the new $`x^{}`$-axis by inclination angle $`i`$ clockwise. The coordinate transformations are given by
$`x^{}`$ $`=`$ $`x_0\mathrm{cos}\theta +y_0\mathrm{sin}\theta `$
$`y^{}`$ $`=`$ $`x_0\mathrm{sin}\theta \mathrm{cos}i+y_0\mathrm{cos}\theta \mathrm{cos}iz_0\mathrm{sin}i`$ (A3)
$`z^{}`$ $`=`$ $`x_0\mathrm{sin}\theta \mathrm{sin}i+y_0\mathrm{cos}\theta \mathrm{sin}i+z_0\mathrm{cos}i.`$
Because our resolution in photometric distance is larger than the disk thickness, the integral in equation (A1) may be simplified by assuming that the exponential disk is infinitely thin, i.e. $`z^{}=0`$ for all points of the disk. The contribution to the integral is zero everywhere, then, except the point where line of sight intercepts the plane of the disk. The value of $`t`$ at the intercept, $`\overline{t}`$, is
$`\overline{t}`$ $`=`$ $`R_{LMC}\mathrm{cos}i\times [\mathrm{cos}\delta \mathrm{sin}(\alpha \alpha _0)\mathrm{sin}\theta \mathrm{sin}i`$ (A5)
$`+\left(\mathrm{sin}\delta \mathrm{cos}\delta _0\mathrm{cos}\delta \mathrm{sin}\delta _0\mathrm{cos}(\alpha \alpha _0)\right)\mathrm{cos}\theta \mathrm{sin}i`$
$`(\mathrm{cos}\delta \mathrm{cos}\delta _0\mathrm{cos}(\alpha \alpha _0)+\mathrm{sin}\delta \mathrm{sin}\delta _0)\mathrm{cos}i]^1.`$
The values of $`x_0`$, $`y_0`$ and $`z_0`$ coordinates at the intercept point follow from equation (A.1). The radius of the intercept point in the plane of the exponential disk is given by
$$r=\sqrt{x^2+y^2}.$$
(A6)
Finally, the expected source density (A1) may be written as
$$\rho _{ij}^{exp}\overline{t}^2\mathrm{exp}(r/R)\mathrm{cos}\delta _j,$$
(A7)
where $`\overline{t}`$ and $`r`$ are given by equations (A5) and (A6), respectively.
### A.2 Power Law Model
Treatment of the spherical power law model is much simpler. Unlike the disk, there is no unique axis of symmetry and one may write the density (2) in the coordinates $`t`$, $`\alpha `$ and $`\delta `$ directly. It is straightforward to verify that
$$r=\sqrt{t^2+R_{LMC}^22R_{LMC}t\left[\mathrm{cos}\delta \mathrm{cos}\delta _0\mathrm{cos}(\alpha \alpha _0)+\mathrm{sin}\delta \mathrm{sin}\delta _0\right]}.$$
(A8) |
warning/0003/hep-th0003127.html | ar5iv | text | # Algebraic characterization of constraints and generation of mass in gauge theories
## 1 INTRODUCTION
In this paper we discuss a new approach to quantum gauge theories, from a group-theoretic perspective, in which mass enters the theory in a natural way. More precisely, the presence of mass will manifest through non-trivial responses
$$U\mathrm{\Psi }=D_{\stackrel{~}{T}}^{(m)}(U)\mathrm{\Psi },$$
(1)
of the wavefunctional $`\mathrm{\Psi }`$ under the action of gauge transformations $`U\stackrel{~}{T}`$, where we denote by $`D_{\stackrel{~}{T}}^{(m)}`$ a specific representation of the gauge group $`\stackrel{~}{T}`$ with index $`m`$. The standard case $`D_{\stackrel{~}{T}}^{(m)}(U)=1,U\stackrel{~}{T}`$ corresponds to the well-known ‘Gauss law’ condition, which also reads $`\mathrm{\Phi }_a\mathrm{\Psi }=0`$ for infinitesimal gauge transformations $`U1+\varphi ^a\mathrm{\Phi }_a`$. The case of Abelian representations $`D_{\stackrel{~}{T}}^{(\vartheta )}(U_n)=e^{in\vartheta }`$ of $`\stackrel{~}{T}`$, where $`n`$ denotes the winding number of $`U_n`$, leads to the well-known $`\vartheta `$-vacuum phenomena. We shall see that more general (non-Abelian) representations $`D_{\stackrel{~}{T}}^{(m)}`$ of the gauge group $`\stackrel{~}{T}`$ entail non-equivalent quantizations (in the sense of, e.g. ) and a generation of mass.
This non-trivial response of $`\mathrm{\Psi }`$ under gauge transformations $`U`$ causes a deformation of the corresponding Lie-algebra commutators and leads to the appearance of central terms proportional to mass parameters (eventually parametrizing the non-equivalent quantizations) in the algebra of constraints, which then become a mixture of first- and second-class constraints. As a result, extra (internal) field degrees of freedom emerge out of second-class constraints and are transferred to the gauge potentials to conform massive bosons (without Higgs fields!).
Thus, the ‘classical’ case $`D_{\stackrel{~}{T}}^{(m)}=1`$ is not in general preserved in passing to the quantum theory. Upon quantization, first-class constraints (connected with a gauge invariance of the classical system) might become second-class, a metamorphosis which is familiar when quantizing anomalous gauge theories. Quantum “anomalies” change the picture of physical states being singlets under the constraint algebra. Anomalous (unexpected) situations generally go with the standard viewpoint of quantizing classical systems and the avoidance of them is evident when quantizing, for example, Yang-Mills theory with chiral fermions, where a cancellation of gauge anomalies is apparently needed; however, these breakdowns, which sometimes are inescapable obstacles for canonical quantization, could be reinterpreted as normal (even essential) situations in a wider setting. Dealing with constraints directly in the quantum arena, this transmutation in the nature of constraints should be naturally allowed, as it provides new richness to the quantum theory.
This cohomological mechanism of mass generation makes perfect sense from a Group Approach to Quantization (GAQ ) framework, which is, at heart, an operator description of a quantum system. Thus, our essential ingredient to define a quantum system will be a given underlying symmetry algebra $`\stackrel{~}{𝒢}`$ rather than an action functional $`S`$, which is the standard starting point in the usual (“classically-oriented”) formulation of QFT.
In order to set the context, let us describe a simple, but illustrative, example of an abstract quantizing algebra $`\stackrel{~}{𝒢}`$ which eventually applies to a diversity of physical systems.
## 2 A SIMPLE ABSTRACT QUANTIZING ALGEBRA
Our particular algebra under study will be the following:
$`[X_j,P_k]`$ $`=`$ $`i\delta _{jk}I,`$
$`[\mathrm{\Phi }_a,\mathrm{\Phi }_b]`$ $`=`$ $`if_{ab}^c\mathrm{\Phi }_c+im_{ab}I,`$ (2)
$`[X_j,\mathrm{\Phi }_a]`$ $`=`$ $`i\stackrel{ˇ}{f}_{ja}^kX_k,[P_j,\mathrm{\Phi }_a]=i\stackrel{ˇ}{f}_{ja}^kP_k,`$
where $`X_j`$ and $`P_k`$ represent standard “position” and “momentum” operators, respectively, corresponding to the extended phase space $``$ of the preconstrained (free-like) theory; The operators $`\mathrm{\Phi }_a`$ represent the constraints which, for the moment, are supposed to close a Lie subalgebra $`\stackrel{~}{𝒯}`$ with structure constants $`f_{ab}^c`$ and central charges $`m_{ab}`$. We also consider a diagonal action of constraints $`\mathrm{\Phi }`$ on $`X`$ and $`P`$ with structure constants $`\stackrel{ˇ}{f}_{ja}^k`$ (non-diagonal actions mixing $`X`$ and $`P`$ lead to interesting “anomalous” situations which we shall not discuss here ). By $`I`$ we simply denote the identity operator, that is, the generator of the typical phase invariance $`\mathrm{\Psi }e^{i\beta }\mathrm{\Psi }`$ of Quantum Mechanics. At this stage, it is worth mentioning that we could have introduced dynamics in our model by adding a Hamiltonian operator $`H`$ to $`\stackrel{~}{𝒢}`$. However, we have preferred not to include it because, although we could make compatible the dynamics $`H`$ and the constraints $`\mathrm{\Phi }`$, the price could result in an unpleasant enlarging of $`\stackrel{~}{𝒢}`$, which would make the quantization procedure much more involved. Anyway, for us, the true dynamics (that which preserves the constraints) will eventually arise as part of the set of good operators (observables) of the theory (see below).
Note that a flexibility in the class of the constraints has being allowed by introducing arbitrary central charges $`m_{ab}`$ in (2). Thus, the operators $`\mathrm{\Phi }_a`$ represent a mixed set of first- and second-class constraints. Let us denote by $`𝒯^{(1)}=\{\mathrm{\Phi }_n^{(1)}\}`$ the subalgebra of first-class constraints, that is, the ones which do not give rise to central terms proportional to $`m_{ab}`$ at the right hand side of the commutators (2). The rest of constraints (second-class) will be arranged by conjugated pairs $`(\mathrm{\Phi }_{+\alpha }^{(2)},\mathrm{\Phi }_\alpha ^{(2)})`$, so that $`m_{+\alpha ,\alpha }0`$.
The simplest (‘classical’) case is when $`m_{ab}=0,a,b`$, that is, when all constraints are first class $`𝒯^{(1)}=𝒯=\stackrel{~}{𝒯}/u(1)`$ and wave functions are singlets under $`𝒯`$. However, the ‘quantum’ case $`m_{ab}0`$ entails non-equivalent quantizations with important physical consequences. This possibility indicates a non-trivial response (1) of the wave function $`\mathrm{\Psi }`$ under $`\stackrel{~}{𝒯}`$. That is, $`\mathrm{\Psi }`$ acquires a non-trivial dependence on extra degrees of freedom $`\varphi _\alpha ^{(2)}`$ (‘negative modes’ attached to pairs of second-class constraints), in addition to the usual configuration space variables $`x_j`$ (attached to $`X_j`$).
Let us formally outline the actual construction of the unitary irreducible representations of the group $`\stackrel{~}{G}`$ with Lie-algebra (2). Wave functions $`\mathrm{\Psi }`$ are defined as complex functions on $`\stackrel{~}{G}`$, $`\mathrm{\Psi }:\stackrel{~}{G}C`$, so that the (let us say) left-action
$$L_{\stackrel{~}{g}^{}}\mathrm{\Psi }(\stackrel{~}{g})\mathrm{\Psi }(\stackrel{~}{g}^{}\stackrel{~}{g}),\stackrel{~}{g}^{},\stackrel{~}{g}\stackrel{~}{G}$$
(3)
defines a reducible (in general) representation of $`\stackrel{~}{G}`$. The reduction is achieved by means of that maximal set of right restrictions on wave functions
$$R_{\stackrel{~}{g}_p}\mathrm{\Psi }=\mathrm{\Psi },\stackrel{~}{g}_pG_p,$$
(4)
(which commute with the left action) compatible with the natural condition $`I\mathrm{\Psi }=\mathrm{\Psi }`$. The right restrictions (4) generalize the notion of polarization conditions of Geometric Quantization and give rise to a certain representation space depending on the choice of the subgroup $`G_p\stackrel{~}{G}`$. For the algebra (2), a polarization subgroup can be $`G_p^{(P)}=F_P\times _sT_p`$, that is, the semi-direct product of the Abelian group of translations $`F_P`$ generated by $`_P\{P_k\}`$ (half of the symplectic generators in $``$) by a polarization subalgebra $`𝒯_p=\{\mathrm{\Phi }_n^{(1)},\mathrm{\Phi }_{+\alpha }^{(2)}\}`$ of $`\stackrel{~}{T}`$ consisting of first-class constraints and half of second-class constraints (namely, the ‘positive modes’). The polarization conditions (4) lead to the configuration-space representation made of wave functions $`\mathrm{\Psi }(x_j,\varphi _\alpha ^{(2)})`$ depending arbitrarily on the group coordinates on $`\stackrel{~}{G}/G_p`$ only. Thus, as mentioned above, wave functions transform non-trivially under the left-action $`L_\varphi \mathrm{\Psi }(\stackrel{~}{g})=D_{\stackrel{~}{T}}^{(m)}(\varphi )\mathrm{\Psi }(\stackrel{~}{g})`$ of $`\stackrel{~}{T}`$ according to a given representation $`D_{\stackrel{~}{T}}^{(m)}`$ like in (1). The physical Hilbert space is made of those wave functions $`\mathrm{\Psi }_{\mathrm{ph}.}`$ that transform as ‘highest-weight vectors’ under $`\stackrel{~}{T}`$, that is, they stay invariant under the left-action of first-class constraints and (let us say) negative second-class modes:
$`L_{\varphi _n^{(1)}}\mathrm{\Psi }_{\mathrm{ph}.}`$ $`=`$ $`\mathrm{\Psi }_{\mathrm{ph}.},n=1,\mathrm{},\mathrm{dim}(T^{(1)}),`$
$`L_{\varphi _\alpha ^{(2)}}\mathrm{\Psi }_{\mathrm{ph}.}`$ $`=`$ $`\mathrm{\Psi }_{\mathrm{ph}.},\alpha =1,\mathrm{},\mathrm{dim}(T/T^{(1)})/2,`$ (5)
which close the subgroup $`T_p\stackrel{~}{T}`$.
The counting of true degrees of freedom is as follows: polarized-constrained wave functions (5) depend arbitrarily on $`d=\mathrm{dim}(\stackrel{~}{G})\mathrm{dim}(G_p)\mathrm{dim}(T_p)1`$ reduced-space coordinates (we are subtracting the phase coordinate $`e^{i\beta }`$). The algebra of observables of the theory, $`\stackrel{~}{𝒢}_{\mathrm{good}}𝒰(\stackrel{~}{𝒢})`$ (the enveloping algebra), has to be found inside the normalizer of constraints, that is:
$$[\stackrel{~}{𝒢}_{\mathrm{good}},𝒯_p]𝒯_p.$$
(6)
From this characterization, the subalgebra of first-class constraints $`𝒯^{(1)}`$ become a horizontal ideal (a gauge subalgebra ) of $`\stackrel{~}{𝒢}_{\mathrm{good}}`$. The Hamiltonian operator has to be found inside $`\stackrel{~}{𝒢}_{\mathrm{good}}`$ by using extra physical arguments.
In what follows, the quantization of massless and massive non-Abelian Yang-Mills, linear Gravity and Abelian two-form gauge field theories are developed from this new approach, where a cohomological origin of mass is pointed out.
## 3 UNIFIED QUANTIZATION OF MASSLESS AND MASSIVE VECTOR AND TENSOR BOSONS
Let us start with the simplest case of the electromagnetic field. Let us use a Fourier parametrization
$`A_\mu (x)`$ $``$ $`{\displaystyle \frac{d^3k}{2k^0}[a_\mu (k)e^{ikx}+a_\mu ^{}(k)e^{ikx}]},`$
$`\mathrm{\Phi }(x)`$ $``$ $`{\displaystyle \frac{d^3k}{2k^0}[\phi (k)e^{ikx}+\phi ^{}(k)e^{ikx}]},`$ (7)
for the vector potential $`A_\mu (x)`$ and the constraints $`\mathrm{\Phi }(x)`$ (the generators of local $`U(1)(x)`$ gauge transformations). The Lie algebra $`\stackrel{~}{𝒢}`$ of the quantizing electromagnetic group $`\stackrel{~}{G}`$ has the following form
$`[a_\mu (k),a_\nu ^{}(k^{})]`$ $`=`$ $`\eta _{\mu \nu }\mathrm{\Delta }_{kk^{}}I,`$
$`[\phi (k),\phi ^{}(k^{})]`$ $`=`$ $`k^2\mathrm{\Delta }_{kk^{}}I,`$
$`[a_\mu ^{}(k),\phi (k^{})]`$ $`=`$ $`ik_\mu \mathrm{\Delta }_{kk^{}}I,`$
$`[a_\mu (k),\phi ^{}(k^{})]`$ $`=`$ $`ik_\mu \mathrm{\Delta }_{kk^{}}I,`$ (8)
where $`\mathrm{\Delta }_{kk^{}}=2k^0\delta ^3(kk^{})`$ is the generalized delta function on the positive sheet of the mass hyperboloid and $`k^2=m^2`$ is the squared mass. Constraints are first-class for $`k^2=0`$ and constraint equations $`\phi \mathrm{\Psi }=0=\phi ^{}\mathrm{\Psi }`$ keep 2 field degrees of freedom out of the original 4, as corresponds to a photon. For $`k^20`$, constraints are second-class and the restrictions $`\phi \mathrm{\Psi }=0`$ keep 3 field degrees of freedom out of the original 4, as corresponds to a Proca field.
For symmetric and anti-symmetric tensor potentials $`A_{\mu \nu }^{(\pm )}`$, the algebra is the following :
$`[a_{\lambda \nu }^{(\pm )}(k),a_{\rho \sigma }^{(\pm )}(k^{})]`$ $`=`$ $`N_{\lambda \nu \rho \sigma }^{(\pm )}\mathrm{\Delta }_{kk^{}}I,`$
$`[\phi _\rho ^{(\pm )}(k),\phi _\sigma ^{(\pm )}(k^{})]`$ $`=`$ $`k^2M_{\rho \sigma }^{(\pm )}(k)\mathrm{\Delta }_{kk^{}}I,`$
$`[a_{\lambda \nu }^{(\pm )}(k),\phi _\sigma ^{(\pm )}(k^{})]`$ $`=`$ $`ik^\rho N_{\lambda \nu \rho \sigma }^{(\pm )}\mathrm{\Delta }_{kk^{}}I,`$
$`[a_{\lambda \nu }^{(\pm )}(k),\phi _\sigma ^{(\pm )}(k^{})]`$ $`=`$ $`ik^\rho N_{\lambda \nu \rho \sigma }^{(\pm )}\mathrm{\Delta }_{kk^{}}I,`$ (9)
where $`M_{\rho \sigma }^{(\pm )}(k)\eta _{\rho \sigma }\kappa _{()}\frac{k_\rho k_\sigma }{k^2}`$ and $`N_{\lambda \nu \rho \sigma }^{(\pm )}\eta _{\lambda \rho }\eta _{\nu \sigma }\pm \eta _{\lambda \sigma }\eta _{\nu \rho }\kappa _{(\pm )}\eta _{\lambda \nu }\eta _{\rho \sigma }`$, with $`\kappa _{(+)}=1`$ and $`\kappa _{()}=0`$. For the massless $`k^2=0`$ case, all constraints are first-class for the symmetric case, whereas the massless, anti-symmetric case possesses a couple of second-class constraints:
$$[\stackrel{ˇ}{k}^\rho \phi _\rho ^{()}(k),\stackrel{ˇ}{k}^{}{}_{}{}^{\sigma }\phi _{\sigma }^{()}(k^{})]=4(k^0)^4\mathrm{\Delta }_{kk^{}}I,$$
(10)
where $`\stackrel{ˇ}{k}^\rho k_\rho `$. Thus, first-class constraints for the massless anti-symmetric case are $`𝒯_{()}^{(1)}=\{ϵ_\mu ^\rho \phi _\rho ^{()},ϵ_\mu ^\rho \phi _\rho ^{()}\},\mu =0,1,2,`$ where $`ϵ_\mu ^\rho `$ is a tetrad which diagonalizes the matrix $`P_{\rho \sigma }=k_\rho k_\sigma `$; in particular, we choose $`ϵ_3^\rho \stackrel{ˇ}{k}^\rho `$ and $`ϵ_0^\rho k^\rho `$. There are $`2=108`$ true degrees of freedom for the symmetric case (a massless graviton) and $`1=65`$ for the anti-symmetric case (a pseudo-scalar particle).
For $`k^20`$, all constraints are second-class for the symmetric case, whereas, for the anti-symmetric case, constraints close a Proca-like subalgebra which leads to three pairs of second-class constraints, and a pair of gauge vector fields $`(k^\lambda \phi _\lambda ^{()},k^\lambda \phi _\lambda ^{()})`$. The constraint equations keep $`6=104`$ field degrees of freedom for the symmetric case (massive spin 2 particle \+ massive scalar field —the trace of the symmetric tensor), and $`3=63`$ field degrees of freedom for the anti-symmetric case (massive pseudo-vector particle).
For non-Abelian $`SU(N)`$ Yang-Mills theories in the Weyl gauge $`A^0=0`$ there is still a residual gauge invariance $`T=\mathrm{Map}(\mathrm{}^3,SU(N))`$. The basic commutators between the non-Abelian vector potentials $`A_a^j(x),j=1,2,3;a=1,\mathrm{},N^21`$, the electric field $`E_a^j(x)`$ and the (Gauss law) constraints $`\mathrm{\Phi }_a(x)`$ are
$`[A_a^j(x),E_b^k(y)]`$ $`=`$ $`i\delta _{ab}\delta ^{jk}\delta (xy)I,`$
$`[\mathrm{\Phi }_a(x),\mathrm{\Phi }_b(y)]`$ $`=`$ $`if_{ab}^c\delta (xy)\mathrm{\Phi }_c(x)`$
$`if_{ab}^c{\displaystyle \frac{\lambda _c}{r^2}}\delta (xy)I,`$
$`[A_a^j(x),\mathrm{\Phi }_b(y)]`$ $`=`$ $`if_{ab}^c\delta (xy)A_c^j(x)`$
$`{\displaystyle \frac{i}{r}}\delta _{ab}_x^j\delta (xy)I,`$
$`[E_a^j(x),\mathrm{\Phi }_b(y)]`$ $`=`$ $`if_{ab}^c\delta (xy)E_c^j(x),`$ (11)
where $`r`$ is the coupling constant and $`\lambda _{ab}=f_{ab}^c\lambda _c`$ is a mass matrix $`(\lambda m^3`$). Let us denote by $`c\mathrm{dim}(T^{(1)})`$ and $`\tau N^21`$ the dimensions of the rigid subgroup of first-class constraints and $`SU(N)`$ respectively. Unpolarized wave functions $`\mathrm{\Psi }(A_a^j,E_a^j,\varphi _a)`$ depend on $`n=2\times 3\tau +\tau `$ field coordinates in $`d=3`$ dimensions; polarization equations introduce $`p=c+\frac{nc}{2}`$ independent restrictions on wave functions, corresponding to $`c`$ non-dynamical coordinates in $`T^{(1)}`$ and half of the dynamical ones; finally, constraints (5) impose $`q=c+\frac{\tau c}{2}`$ additional restrictions which leave $`f=npq=2c+3(\tau c)`$ field degrees of freedom (in $`d=3`$). These fields correspond to $`c`$ massless vector bosons (2 polarizations) attached to $`T^{(1)}`$ and $`\tau c`$ massive vector bosons. In particular, for the massless case, we have $`c=\tau `$, since constraints are first-class (that is, we can impose $`q=\tau `$ restrictions) and constrained wave functions have support on $`f_{m=0}=3\tau \tau =2\tau f_{m0}`$ arbitrary fields corresponding to $`\tau `$ massless vector bosons. The subalgebra $`𝒯^{(1)}`$ corresponds to the unbroken gauge symmetry of the constrained theory. There are distinct symmetry-breaking patterns $`\stackrel{~}{𝒯}𝒯^{(1)}`$ according to the different choices of mass-matrices $`\lambda _{ab}=f_{ab}^c\lambda _c`$ in (11). |
warning/0003/astro-ph0003338.html | ar5iv | text | # Large-Scale Cosmic Shear Measurements
## 1. Introduction
Weak lensing provides a potentially powerful probe of mass fluctuations in the Universe (Gunn (1967); Mellier (1999) and references therein). Three independent groups have recently presented estimates of the shear variance from deep ‘blank-field’ CCD imaging surveys. Van Waerbeke et al. (2000) (hereafter vWME+) measured the shear variance in circular cells of radii ranging from $`0^{}.7`$ to $`3^{}.5`$; Bacon et al. (2000) (hereafter BRE) measured the shear variance in square cells of side $`8^{}.0`$ and Wittman et al. (2000) (hereafter WTK+) have provided estimates of the shear-shear correlation function at separations $`3^{}.25`$, $`8^{}.5`$ and $`22^{}.0`$. Here we present shear variance measurements from $`1.5`$ square degrees of deep photometry obtained as part of our ongoing weak lensing survey. We find results which are broadly in good agreement with the recently published estimates.
## 2. The Data
The data were taken at the 3.6m CFHT telescope using the $`8192\times 8192`$ pixel UH8K camera at prime focus. The camera delivers a field size of $`0^{}.5`$ with $`0^{\prime \prime }.207`$ pixels. Our survey strategy has been to target blank fields in six widely separated areas for ease of scheduling, and in each area we plan to make $`6`$ or so pointings scattered over a region of extent $`3^{}`$. In January 1999 the UH8K was replaced by the CFH12K camera. By that time we had completed 1 hour integrations in both V and I for two pointings in each of three areas. The field names, centers and also the estimated seeing are given in table 1. One of the 8 devices in the mosaic (lying in the NW corner of our images) has very poor charge transfer efficiency and the data from this device were discarded. After further masking of regions around bright stars the total useful solid angle per field is $`0.16`$ square degrees.
## 3. Data Reduction
The data were reduced much as for our MS0302 supercluster observations Kaiser et al. (1999). After flat fielding, an object finder was applied to each image, and a set of bright but non-saturated stars extracted for registration purposes. The positions of these stars, along with celestial coordinates from the USNOA Monet (1998) catalog, were used to find a mapping from image pixel coordinates to orthographic sky coordinates.
Images of the stars were analyzed to generate a model for the point spread function (PSF) $`g(𝐱;𝐫)`$, this being the 2-dimensional profile of a star with centroid at $`𝐫`$ measured in coordinates $`𝐱`$ being measured relative to the centroid. The model is a sum of 2-D image valued modes $`g_i(𝐱)`$: $`g(𝐱;𝐫)=_ic_i(𝐫)g_i(𝐱)`$ with coefficients $`c_i(𝐫)`$ which are low order polynomials in star position. We found a 1st order model to be adequate to describe the variation of the PSF with position on the chip (though see below).
The astrometric solutions, the PSF models, and also standard star observations were first used to create a set of photometrically calibrated ‘homogenized’ images which were degraded to have a common identical PSF. These images were compared in order to identify cosmic rays and other transient events. Next, for each raw image we generated a re-circularized image by convolving with a 90-degree rotated version of the PSF model. These, as well as the raw images, were then warped to sky coordinates, with previously identified cosmic rays being removed, and the stacks of images combined to provide a quilt of overlapping images. This procedure results in three images: a median of the raw images (in which the PSF is generally non-circular), a median of the re-circularized images, and also a sky-noise image. The final summed images were sampled with $`0^{\prime \prime }.15`$ pixel size.
The final object catalogs were obtained by applying hfindpeaks to the median averaged raw images, this program having been modified in order to allow properly for the non-trivial noise correlations in the final images. After applying aperture photometry analysis, shapes of the objects were measured as described in ?). The essential result of this is a polarization vector $`q_\alpha =M_{\alpha lm}d^2xw(x)x_lx_mf(𝐱)`$ formed as a combination of weighted second moments of the re-circularized image, and a polarizability tensor which describes the response of the polarization to gravitational shear. The weight function $`w(x)`$ was taken to be a Gaussian ball of 2 pixels in scale length. For convenience, the quantities actually generated were the normalized polarization $`\widehat{q}_\alpha =q_\alpha /\sqrt{q_\beta q_\beta }`$ and a polarizability $`Q`$ defined such that for galaxies with this shape and size, the expectation value of $`\widehat{q}_\alpha `$ is $`\widehat{q}_\alpha =Q\gamma _\alpha `$. The $`q_\alpha `$ values were corrected for artificial shear introduced in the image warping by slight errors in our astrometric solution.
Finally a selection on significance of $`4\nu 100`$ was applied to select a ‘faint galaxy’ catalog. The corresponding magnitude limits are somewhat fuzzy since significance depends on both size and magnitude. The counts (number per magnitude interval) turn over at about $`m_\mathrm{I}24`$ and $`m_\mathrm{V}25`$. Numbers of objects in the final catalogs are shown in table 3. The density of objects on the sky is very similar to that in the images obtained by vWME+, WTK+, and slightly higher than the density of objects in the BRE sample, so our shear variance estimates should be more or less directly comparable.
Given a single galaxy, a fair (but very noisy) estimate of the shear is $`\gamma _\alpha =\widehat{q}_\alpha /Q`$. However, since the normalized polarization response $`Q`$ varies with shape and size of the object (small objects having very little response for example), to measure the mean shear in a region containing $`N`$ galaxies one should weight the individual estimates by $`Q^2`$ so the optimal mean shear estimate is $`\overline{\gamma }_\alpha =Q\widehat{q}_\alpha /Q^2`$. This assumes that one has little prior knowledge of the of galaxy redshifts. If the averaging region contains a large number of galaxies, as is the case for the cells considered here, one can replace $`Q^2`$ by $`NQ^2`$ where the $`Q^2`$ is an average over all of the galaxies in the catalog. The optimal mean shear is then the average of weighted shear values for the individual galaxies: $`\gamma _G=\omega _G\widehat{q}_G/Q_G`$, with normalized weight $`\omega _GQ_G^2/Q^2`$. The quantity $`Q^2`$ is a useful measure of the image quality. It is equal to the inverse shear variance per galaxy, so, for instance, the statistical uncertainty in the mean shear $`\overline{\gamma }`$ measured from a sample of $`N`$ galaxies is $`\sigma _{\overline{\gamma }}^2=1/(NQ^2)`$. The values of $`Q^2`$ are also given in table 3.
Preliminary results of this analysis Wilson et al. (1999) in the form of estimates of the net shear for each of the six fields and for catalogs generated from the I and V images separately gave shear values typically of about 1%, but with a few larger values. These, however, showed little correlation between the two passbands, suggesting that the results were contaminated by some systematic error. Examination of the re-circularized images of the stars in the fields with seemingly spurious shear values revealed at least a major part of the problem. The mean stellar polarization was found to vary systematically with magnitude. This is to be expected for very bright stars where the pixels saturate and charge begins to bleed along the slow direction of the CCD. The effect found here had the same signature (a trend for $`q_1`$ to become negative for bright objects) but appeared at a low level and, unexpectedly, for stars much fainter than the saturation limit. The result was that in some cases our PSF model fitting procedure, which weighted stars according to brightness, did not correctly recircularize the faint stars as it should (since we require the PSF appropriate for faint galaxies). The effect seemed to be variable, and also tended to be associated with particular chips. As a simple fix, we fit the residual polarizations for the faintest stars ($`I>18`$, $`V>21`$) to a 4th order spatial polynomial and then used a smear polarizability analogous to that defined by ?) to correct the galaxy $`q_\alpha `$ values. This reduced the spurious shear values considerably.
## 4. Cosmic Shear Variance
We have chosen to focus here on a single simple statistic: the variance of the shear averaged in cells of various sizes. This is the statistic used by vWME+ and BRE and is simply related to the shear covariance function presented by WTK+. The cell averaged shear variance is also simply computable from the spectrum of mass fluctuations (e.g. Kaiser (1992)), so this provides a useful link between observation and theory. Now each weighted shear estimate $`\gamma _G`$ consists of a random intrinsic component $`\gamma _{G\mathrm{int}}`$ and a ‘cosmic’ component proportional to the integral of the tidal field along the line of sight. Modeling the cosmic shear as the sum over a set of statistically independent screens, we have
$$\gamma _G=\gamma _{G\mathrm{int}}+\omega _G\underset{S}{}\gamma _S(𝜽_G)\beta (z_G,z_S)$$
(1)
where $`\gamma _S(𝜽)`$ is the shear field for the $`S`$th screen and for fictitious sources at infinite distance, and $`\beta \mathrm{max}(0,1D_{SG}/D_{OG})`$ is the usual ratio of angular diameter distances.
This model allows one to compute the variance of the shear averaged over galaxies falling in a cell on the sky. We will also be interested in the co-variance of shear measured in different passbands. Consider the mean shear $`\overline{\gamma }_P=(1/N_P)\gamma _P`$ for a specific cell and for galaxies found in two passbands $`P=A,B`$. Averaging over realizations of random intrinsic shear values, and also averaging over an ensemble of realizations of cosmic shear screens, yields
$$\overline{𝜸}_A\overline{𝜸}_B=\frac{N_{AB}}{N_AN_B}𝜸_A𝜸_B+\underset{S}{}\omega _A\beta _{AS}\omega _B\beta _{BS}\overline{\gamma }_S^2$$
(2)
where $`N_{AB}`$ is the number of objects in the cell which were detected in both passbands. This formula is also valid when $`A`$ and $`B`$ are the same. The expectation value of the dot product of the cell averaged shear is therefore equal to a noise term plus a cosmic term which is a sum of the cell-averaged shear variances for the screens. Interestingly, the noise term involves the total shear variance $`𝜸_A𝜸_B`$, containing both intrinsic and cosmic contributions, which is convenient since this is the quantity that one can actually measure. In obtaining (2) we assumed that the faint galaxies are randomly distributed on the sky, this being motivated by the fact that the angular correlation function is very small, with $`w(\theta )<10^2`$ on all relevant scales.
To implement this, for each field we averaged the shear in a grid of contiguous square cells of side $`L`$, and those cells in the lower quartile of occupation number were discarded. To obtain an estimate of the cosmic shear variance we then computed for each field
$$\overline{𝜸}^2_{AB}=\frac{1}{n_{\mathrm{cells}}}\underset{\mathrm{cells}}{}\left[\overline{𝜸}_A\overline{𝜸}_BN_{AB}𝜸_A𝜸_B/(N_AN_B)\right].$$
(3)
The shear covariance functions are $`𝜸_\mathrm{I}𝜸_\mathrm{I}=1/Q^2_\mathrm{I}`$ and $`𝜸_\mathrm{V}𝜸_\mathrm{V}=1/Q^2_\mathrm{V}`$ whereas $`𝜸_\mathrm{I}𝜸_\mathrm{V}`$ was estimated by correlating the shears for objects which were detected in both the I and V catalogs. The shear variances estimated from the separate fields were then averaged together to obtain a final cosmic shear variance with uncertainty estimated from the scatter of the field estimates about the mean.
The diagonal components of $`\overline{𝜸}^2_{AB}`$ provide estimates for the shear variance for the respective passbands, with strength roughly proportional to the square of the mean distance to the galaxies (assuming a spectrum of mass fluctuations with index $`n1`$), and the off-diagonal components should lie somewhere in between. The results are shown in figure 4 and in table 4 and deviate somewhat from this expectation: the I-V cross correlation lies systematically below both the I- and V-band shear variance estimates. These results are robust to changes in the order of the polynomial in the stellar polarization model, and they are not caused by a few discrepant cells.
The difference between the I and V band shear variance may be due to differences in the redshift distributions, some evidence for which was found by ?) in their study of MS1054. However, we typically find about 60% of the galaxies are detected in both passbands, so this requires fairly high redshifts for the blue galaxies. For example, assume that the I-band sample has a redshift distribution like that measured by Cowie (personal communication) in the range $`23<\mathrm{I}_{\mathrm{AB}}<24.0`$, but that the V band sample contains an additional 40% population of higher redshift galaxies. If we place these at redshift 3, we find that the V-band shear variance is about a factor 2 higher than the I-band, much as seen in figure 4. However the I-V cross-correlation is then predicted to be about 30% higher than the I-band variance, which is not seen.
The simplest interpretation of these results is that the shear inferred from the I- and V-band data separately has been inflated by residual systematic errors of some kind. The level of these errors is on the order of 2 percent rms shear on scales of a few arc-minutes, falling to somewhat below the 1 percent level on $`30^{}`$ scales. If so, the most reliable estimate of the cosmic shear variance is provided by the I-V cross-correlation since systematic errors which are uncorrelated between the passbands will cancel out. Of course there is no guarantee that the cross-correlation is not affected by some source of error which is common to both passbands, a specific example of which is artificial shear arising from intrinsic correlation of galaxy shapes in clusters etc. due to tidal effects.
The I-V shear variance estimator is shown with an expanded vertical scale in figure 4. Also shown are the recently announced results. The BRE result is shown as presented in their paper and with total error estimate including cosmic variance. The vWME+ circular cell average shear are plotted against $`L=\sqrt{\pi }\theta `$. The vWME+ error bars are statistical only. WTK+ presented estimates of the ellipticity correlation function $`C_1(\theta )=ϵ_1(0)ϵ_1(\theta )`$. We have converted their $`C_1(\theta )`$ to an equivalent shear variance using formulae from ?) with $`ϵ=2\gamma `$ and assuming a spectral index $`n=1`$. The lower panel shows the variance multiplied by averaging box size $`L`$. For a $`n=1`$ spectrum, corresponding to a mass auto-correlation function $`\xi (r)1/r^2`$, this quantity should independent of scale.
At small scales $`<10`$ arcmin there seems to be remarkably good agreement between the independent estimates. Note that the measurements were made using three separate observing facilities. At $`L=3^{}.75`$ we find $`\overline{𝜸}^22.510^4`$. This about a factor 4-5 lower than the prediction for a light-traces mass $`\mathrm{\Omega }_m=1`$ cosmology, and an effective redshift for the background galaxies $`z_{\mathrm{eff}}=1`$ Kaiser (1992); Jain & Seljak (1997).
At larger scales the shear variance we find falls below that of WTK+. Their largest scale estimates appear to conflict with our null result at about the 2-sigma level. Our large-angle results are also smaller than the $`\mathrm{\Omega }_\mathrm{m}=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ theoretical model predictions.
## 5. Discussion
For an effective background galaxy redshift of $`z_{\mathrm{eff}}1.0`$ these measurements probe mass fluctuations in a shell peaked at $`z0.4`$. At this redshift the $`30^{}`$ field size corresponds to a comoving distance of about $`6h^1`$Mpc, so the cell variances presented here probe scales in the range $`0.46h^1`$Mpc. On the smaller end of this scale we find very good agreement with recently announced estimates from other groups, and also with canonical cosmological theory predictions. It is hard to definitively rule out the possibility that the small angle measurements are inflated by systematic errors, but one can safely rule out theories such as light-traces mass high density models which predict shear variance a factor $`5`$ higher than our results.
On larger scales our measurements are extremely precise, yet we find only a null detection for our largest cells. These results show that on large scales the rms shear is at most a fraction of a percent. The apparent discrepancy between these results and the theoretical predictions is quite interesting, and suggests a steepening of the mass correlation function at scales $`12h^1`$Mpc. More data are needed however to definitively confirm this.
## 6. Acknowledgements
The results here were extracted from data taken at the Canada France Hawaii Telescope. The analysis was supported by NSF grants AST95-00515, AST99-70805. GW gratefully acknowledges financial support from the estate of Beatrice Watson Parrent and from Mr. & Mrs. Frank W. Hustace, Jr. We thank Peter Schneider and Gary Bernstein for helpful suggestions. |
warning/0003/cond-mat0003329.html | ar5iv | text | # Anomalous temperature dependence of magnetic quantum oscillations in CeBiPt
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## Abstract
Shubnikov–de Haas (SdH) and Hall-effect measurements of CeBiPt and LaBiPt reveal the presence of simple and very small Fermi surfaces with hole-like charge carriers for both semimetals. In the magnetic material, CeBiPt, we observe a strong temperature dependence of the SdH frequency. This highly unusual effect might be connected with an internal exchange field within the material and a strongly spin-dependent scattering of the charge carriers.
\]
The Shubnikov–de Haas (SdH) effect, i.e., quantum oscillations in the magnetoresistivity $`\rho (B)`$, was discovered about 70 years ago in the semimetal Bi . Since then, innumerable studies on the SdH effect and the related de Haas–van Alphen (dHvA) effect in the magnetization were reported for many metals, metallic compounds, and semimetals. Both effects arise from the oscillations of the free energy of the electrons as with increasing field $`B`$ a sudden change of population of Landau cylinders occurs whenever a Landau cylinder is pushed through an extremal cross-section $`A`$ of the Fermi surface perpendicular to the magnetic field. The oscillations are periodic in 1/$`B`$ and the frequency $`F`$ is directly proportional to $`A`$, i.e., $`F=(\mathrm{}/2\pi e)A`$. The theory of the dHvA effect is well established , and also the SdH effect, reflecting the more involved transport properties, is principally understood . In recent years, SdH and dHvA effects have been exploited to establish the Fermi surface of complex materials with strong electron correlations such as heavy-fermion systems . In many cases, the effective mass of the carriers in a given band of the Fermi surface can be determined from the temperature dependence of the amplitude of the oscillations. Here we report on a new observation, namely a strongly temperature-dependent frequency of the SdH oscillations. This anomalous $`T`$ dependence of $`F`$ was found for a certain field orientation in semimetallic CeBiPt, while it is absent in the homologous - except for the Ce 4$`f`$ electron \- semimetal LaBiPt.
The single crystals were grown at Hiroshima University with the Bridgman technique in a hermetically sealed Mo crucible from the starting materials Ce (m5N, Ames Laboratory), Bi (m5N), Pt (m3N), and La (m4N). To avoid oxidation of elemental Ce or La during handling, CePt or LaPt were first prepared by argon-arc melting and the appropriate amount of Bi then was added for single-crystal growth. The crucible, sealed under Ar atmosphere, was heated to 1350C in a furnace (with an intermediate halt at 500C for 2 h) and after 12 hours was slowly cooled by moving it out of the central zone of the furnace with 1 mm/h. This cooling process from 1350C to 20C took 6 days, then the furnace was switched off. Whether the 1-1-1 compounds or 3-4-3 compounds like Ce<sub>3</sub>Bi<sub>4</sub>Pt<sub>3</sub> are formed, depends very sensitively on the excess amount of Bi. The 1-1-1 compounds show the F$`\overline{4}`$3m cubic structure previously determined for polycrystalline samples with no indications of second phases. The lattice constant 6.8338 (25)$`\mathrm{\AA }`$ found for CeBiPt is in good agreement with the previously reported value . For the LaBiPt crystal, Laue-diffraction pictures showed some mosaicity of the sample.
The longitudinal and transverse magnetoresistance (Hall effect) were measured with a standard <sup>3</sup>He cryostat setup up to 15 T in Karlsruhe, and up to 28 T at the High Magnetic Field Laboratory in Grenoble. Both sets of data agree with each other in region of overlap. Six gold wires were glued with graphite paste to the samples, thereby enabling to measure simultaneously the longitudinal ($`\rho _{xx}`$) and transverse ($`\rho _{xy}`$) magnetoresistance. These resistances were measured by use of a low-frequency ($`16`$ Hz) lock-in technique for normal and reversed field orientations which allowed a well-defined separation of $`\rho _{xx}`$ and $`\rho _{xy}`$. The dHvA signal of the LaBiPt sample was measured by means of a capacitance cantilever torque magnetometer. All sample holders could be rotated in situ around one axis.
Fig. 1(a) and (b) show the magnetoresistance $`\rho (B)`$ for CeBiPt and for LaBiPt with field along \[1 0 0\] for several temperatures. The current direction was perpendicular to $`B`$. In both cases, clear oscillations are seen. For LaBiPt at $`T`$ = 0.4 K, the strong decrease of $`\rho `$ below about 1 T (with a 50% value at 0.45 T) indicates the onset of superconductivity with $`T_c`$ = 0.88 K \[see inset of Fig. 1(b) for $`\rho (T)`$\]. The rather large critical-field slope ($`\mathrm{d}B_{c2}/\mathrm{d}T)_{T_c}`$ -1 T/K hints at a small coherence length. For some field directions we observed a beat pattern for LaBiPt which might be attributed to the mosaicity mentioned above. We can, however, not exclude the existence of a second extremal Fermi surface of small area in LaBiPt.
For CeBiPt, a first unusual observation is the sharp drop of the magnetoresistance at low fields and low temperatures \[Fig. 1(a)\]. Below about $`T_N=1`$ K, CeBiPt is in an antiferromagnetically ordered state as was evidenced by specific-heat and magnetization measurements for our samples . Below $`T_N`$, the field-dependent magnetization, $`M`$, has a maximum in $`\mathrm{d}M/\mathrm{d}B`$ at about 0.3 T . Therefore, the negative magnetoresistance at low fields presumably is caused by antiferromagnetic fluctuations which become reduced in an applied magnetic field.
The most striking observation for CeBiPt is the shift of the oscillating SdH signal with temperature. This can be made more apparent when the resistivity is converted to conductivity via $`\sigma _{xx}=\rho _{xx}/(\rho _{xx}^2+\rho _{xy}^2)`$ making use of the simultaneously measured transverse magnetoresistance $`\rho _{xy}`$ . A smooth background longitudinal conductivity $`\sigma _0`$ can be fitted to the $`\sigma _{xx}`$ data to obtain the SdH signal $`\mathrm{\Delta }\sigma /\sigma _0=(\sigma _{xx}\sigma _0)/\sigma _0`$ as shown in Fig. 2. The inset of Fig. 2 displays the sequence of the position of 1/$`B`$ of the oscillation maxima and minima.
The oscillations are indeed periodic in $`1/B`$ within our resolution and for the limited field range. The oscillation frequency $`F`$ is directly obtained from the slope of $`1/B`$ vs. oscillation index $`n`$. $`F`$ decreases from 58.2 T at 0.4 K to about 35 T at 10.3 K. Additional data taken in Karlsruhe in fields up to 15 T are fully in line with this decrease (see Fig. 3). The total change of $`F`$ between 0.4 and 10.3 K corresponds to an apparent reduction of the Fermi-surface cross-section by more than 50%! The inset of Fig. 3 shows the decrease of the amplitude $`A`$ of the SdH oscillations with increasing $`T`$ as measured at 10.5 T. $`A`$ can be described quite well with the standard expression $`AT/\mathrm{sinh}[(14.69(T/B)(m^{}/m_0))]`$, where $`T`$ and $`B`$ are expressed in K and T, respectively . We obtain the effective mass $`m^{}=0.24m_0`$, where $`m_0`$ is the free-electron mass. A simple consistency check serves to test whether we are indeed dealing with a quantum-oscillation phenomenon. From the field dependence of $`\mathrm{\Delta }\sigma /\sigma _0`$ at fixed $`T`$ one can estimate the charge-carrier scattering rate $`\tau ^1`$. The Dingle temperature $`T_D4`$ K obtained from the field dependence of the SdH amplitude (at lower fields) corresponds to $`\tau ^1=3.3\times 10^{12}`$ s<sup>-1</sup>. Together with the carrier density $`n_h=7.7\times 10^{17}`$ cm<sup>-3</sup> obtained from the Hall effect at liquid-He temperature we obtain in the simple Drude model $`\rho 3.7`$ m$`\mathrm{\Omega }`$cm which is perfectly in line with the measured resistivity.
The angular dependence of $`F`$ for CeBiPt is depicted in Fig. 4(a) for $`T`$ = 0.43 K and 4.21 K. The anomalous $`T`$ dependence is found only for field along the \[1 0 0\] and \[0 1 0\] directions where $`F`$ is found to be somewhat larger, i.e., about 58 T. Between approximately $`\theta =20^{}`$ and $`70^{}`$, where $`\theta `$ is the angle from the \[1 0 0\] direction, $`F`$ is constant at 30 T and furthermore independent of $`T`$. For $`B`$ along the \[1 1 1\] direction a very low SdH frequency of about 20 T - but again independent of $`T`$ \- was found. LaBiPt shows a temperature-independent angular dependence of $`F`$ \[Fig. 4(b)\]. The SdH and the dHvA frequency are within error bars identical and somewhat larger ($`F`$ increases from 65 T for $`B`$ along \[1 0 0\] to 95 T for $`B`$ along \[1 1 0\]) than $`F`$ for CeBiPt. The large magnetization of CeBiPt due to the Ce $`4f`$ electrons prevented us from taking dHvA data of the torque magnetization for this compound. Nevertheless, the very good agreement between SdH and dHvA data for LaBiPt lends support to the SdH data for CeBiPt. For both materials we were able to observe a SdH signal over the whole angular range. This suggests that the Fermi surfaces are simple single-connected hole pockets. The volume enclosed by the Fermi surface is estimated to comprise only about $`1.5\times 10^4`$ of the volume of the first Brillouin zone for CeBiPt. This is consistent with the low (hole-like) charge-carrier concentration of $`n_h=7.7\times 10^{17}`$ cm<sup>-3</sup>. Within a free-electron picture, $`n_h`$ corresponds to a Fermi energy of $`E_F=\mathrm{}^2(3\pi ^2n_h)^{2/3}/2m^{}=12.8`$ meV. Assuming a circular Fermi-surface cross section, i.e., $`A=\pi k_F^2`$ with the Fermi wavevector $`k_F`$, this results in a SdH frequency of $`F=m^{}F/\mathrm{}e27`$ T in nice agreement with the experimental $`F`$ between $`20^{}`$ and $`70^{}`$.
The main point of the present investigation is the observation of a temperature dependence of the quantum-oscillation frequency which is found only for the Ce-based metal. To our knowledge such an effect has never been observed before. As a possible origin for this behavior one could assume a temperature-dependent charge-carrier density which would lead to a change in the Fermi-surface topology. However, the simultaneously measured Hall constant $`R_H=1/n_he`$, is independent of temperature.
Another possibility would be a reconstruction of the Brillouin zone. This might be caused, e.g., by an antiferromagnetic ordering, which introduces an extra periodicity into the lattice. The effect of magnetic ordering has previously been observed in a field-dependent change of the dHvA frequency in NdIn<sub>3</sub> or in an unusual temperature dependence of the dHvA amplitude in YbAs , SmSb , and SmAgSb<sub>2</sub> . As mentioned, CeBiPt undergoes an antiferromagnetic transition at about 1 K. However, no abrupt change of $`F`$ with $`T`$ around $`T_N`$ or $`B`$ but a rather smooth variation over a large $`T`$ range is observed. This indicates that a Brillouin-zone reconstruction must be ruled out as a cause for the observed frequency change.
One important fact which is evident from our investigation is the absence of any unusual effect for the non-magnetic sister compound LaBiPt. Therefore, it is clear that the magnetism of the Ce atoms affects the magnetic quantum oscillations. For a sample possessing an internal magnetization, the magnetic induction $`B_i=\mu _0(H+M)`$ is different from the externally applied field $`B=\mu _0H`$. The SdH signal is proportional to $`\mathrm{sin}(2\pi F/B_i)`$. Therefore, a temperature-dependent magnetization would cause a change in the SdH frequency. However, in order to explain the experimentally observed increase of the SdH frequency with decreasing temperature, $`M`$ would have to become smaller at lower temperatures. This is contrary to the usual behavior of $`M`$ and not in line with the measured magnetization which is about $`\mu _0M=0.2`$ T at $`\mu _0H=12`$ T and $`T=1.7`$ K and decreases with increasing temperature. Moreover, the magnitude of $`\mu _0M`$ compared to $`B`$ is much too small to account for the observed change of $`F`$.
An anomalous situation may occur when we suppose that a temperature-dependent (and possibly also a field-dependent) phase difference $`\mathrm{\Phi }`$ exists between the spin-up and spin-down oscillations and that, in addition, their amplitudes may be different . The latter can easily occur because the scattering by magnetic impurities (or magnetic sublattices) is in general appreciably spin dependent. The usual splitting of the spin-up and spin-down Landau levels gives rise to a phase shift of the oscillations by $`\mathrm{\Phi }=\pm \pi gm^{}/m_0`$, where $`g`$ is the $`g`$ factor of the conduction electrons. If an antiferromagnetic exchange interaction $`B_{ex}`$ is present, the phase difference can be expressed as $`\mathrm{\Phi }=\pm \pi (gB_{ex}/B)m^{}/m_0`$ . The superposition of the spin-up and spin-down oscillations gives a signal which is $`M=a_{}\mathrm{sin}(\psi +\mathrm{\Phi }/2)+a_{}\mathrm{sin}(\psi \mathrm{\Phi }/2)`$, with $`\psi =(2\pi F/B)\pm \pi /4`$ and the spin-dependent amplitudes $`a_{}`$ and $`a_{}`$. In case of largely different amplitudes, the frequency of the oscillating signal becomes $`F^{}=F\pm B_{ex}/4`$. Adopting this scenario would, however, imply a change of $`B_{ex}`$ with $`T`$ by several 10 T which appear unlikely. On the other hand, the apparent decrease of the oscillation frequency in $`1/B`$ above about 0.12 T<sup>-1</sup> at low $`T`$ does suggest an influence of internal and/or exchange fields. Calculations of spin splitting in low-carrier-density systems are necessary to check the viability of this scenario.
In summary, we have presented an unusual $`T`$ dependence of SdH oscillations in CeBiPt which are due to the magnetic Ce $`4f`$ electrons. Simple estimates of the scattering rate derived from the Dingle temperature and of the volume enclosed by the Fermi surface yield good agreement with the corresponding quantities obtained from other measurements. This lends confidence to the assignment of the observed apparent $`T`$ dependence being a feature associated with the magnetic quantum oscillations in CeBiPt. However, theoretical work is needed to unravel the origin of this new phenomenon.
We would like to thank A. Schröder and J. Sereni for many helpful discussions. This work was supported by the Deutsche Forschungsgemeinschaft, SFB 195, and the TMR Programme of the European Community under contract No. ERBFMGECT950077. |
warning/0003/hep-ph0003037.html | ar5iv | text | # Ultra High Energy Cosmic Rays from Axion Stars
## Abstract
We propose a model in which ultra high energy cosmic rays are produced by collisions between neutron stars and axion stars. The acceleration of such a cosmic ray is made by the electric field, $`10^{15}(B/10^{12}\text{G})\text{eV}\text{cm}^1`$, which is induced in an axion star by the strong magnetic field $`B>10^{12}`$ G of a neutron star. As we have shown previously, similar collisions generate gamma ray bursts when the magnetic field is much smaller, e.g. $`10^{10}`$ G. If we assume that the axion mass is $`10^9`$ eV, we can explain huge energies of the gamma ray bursts $`10^{54}`$ erg as well as the ultra high energies of the cosmic rays $`10^{20}`$ eV. In addition, it turns out that these axion stars are plausible candidates for MACHOs. We point out a possibility of observing monochromatic radiations emitted from the axion stars.
preprint: Nisho-00/1
The origin of ultra high energy cosmic rays ( UHECRs ) is one of most mysterious puzzles in astrophysics. UHECRs with energies $`10^{20}`$ eV can not travel in distance more than $`50`$ Mpc owing to the interactions between UHECRs and the cosmic backgrand radiations. Observations, however, show that there are no visible candidates for the generators of such UHECRs in the arrival directions of UHECRs. On the other hand, dark matter in the Universe is one of most mysterious puzzles in cosmology. Axion is a plausible candidate for the dark matter. Probably, some of axions may form boson stars ( axion stars ) in the present Universe by gravitational cooling or gravitational collapse of axion clumps formed at the period of QCD phase transition. Here we explain a generation mechanism of UHECRs and discuss a production rate of UHECRs assuming collisions between axions stars and neutron stars. We also show that the collisions cause emission of observable monochromatic radiations from the axion stars.
We have previously proposed a possible generation mechanism of gamma ray burst ( GRB ). According to the mechanism the collision between axion star and neutron star generates a gamma ray burst; the axion star dissipates its mass energy very rapidly under the strong magnetic field of the neutron star. Thus, the energy released in the collision is given by the mass $`M_a`$ of the axion star. Typically it is given by $`M_a10^5M_{}(10^5\text{eV}/m_a)`$ where $`m_a`$ denotes the axion mass. In the analysis we have taken the mass such as $`m_a10^5`$ eV as suggested observationally in standard invisible axion models. The mass, $`M_a10^{49}`$ erg, corresponding to this choice, however, is not enough to explain a huge energy observed in GRB980123 even if jet is assumed in the GRB.
In this paper, we explain an acceleration mechanism of UHECRs assuming the axion mass, $`m_a10^9`$ eV, although the choice is not conventional. The essence of the acceleration is that a strong electric field $`10^{15}(B/10^{12}\text{G})\text{eV}\text{cm}^1`$ is induced in the axion star when it is exposed to the magnetic field $`B>10^{12}`$ G of a neutron star. This electric field can accelerate charged particles and makes them obtain the huge energies $`10^{20}`$ eV. Furthermore, it turns out that the energy $`10^{54}`$ erg observed in the GRB can be also explained assuming a moderate jet of the GRB. This is because the mass of the axion star is given by $`10^1M_{}10^{53}`$ erg with the assumption of $`m_a`$. Thus the collisions between the axion stars and the neutron stars are possible sources of both UHECRs and GRBs. Additionally, it turns out that the axion star is a plausible candidate for MACHO because the value of the mass is suitable for explaining the observations of MACHOs. Since all of baryonic candidates for MACHOs have been argued to have serious difficulties, nonbaryonic ones like the axion stars are favored. We discuss that the collisions generate monochromatic radiations with a frequency $`m_a/2\pi 2.4\times 10^5`$ Hz. With the detection of such radiations we can test our model and determine the axion mass.
As we shall show below, we need to separate two cases of the collisions, the collision with a neutron star possessing relatively weak magnetic field $`10^{10}`$ G and the one with a neutron star possessing relatively strong magnetic field $`>10^{12}`$ G. In the first case the axion star collides directly with the neutron star and dissipates its whole energy inside of the outercrust of the neutron star. The ultra high energy cosmic rays are not produced in this case but gamma ray bursts are produced. On the other hand, in the second case the axion star never collide directly with the neutron star because it evapolates before the collision. This is because the axion star induces a stronger electric field when it is exposed to the stronger magnetic field of the neutron star. But such a strong electric field decays very rapidly owing to electron-positron pair creations when the strength of the electric field goes beyond a critical value. Namely, as the axion star approaches the neutron star, the strength of the magnetic field around it increases gradually and reaches a critical value beyond which the electric field induced in the axion star decays very rapidly. It means that the axion star itself decays very rapidly. In this case the ultra high energy cosmic rays are produced. Gamma ray bursts might be also produced, but they are emitted in a cone with much small solid angle and their duration is very short ( less than millisecond ).
Let us first explain briefly axion stars ( ASs ) and how they induce strong electric fields under a magnetic field of a neutron star. The AS is a coherent object of the real scalar field $`a(x)`$ describing the axion. It is represented by a solution of the equation of the axion field coupled with gravity. An approximate form of the solution is given by
$$a(x)=f_{PQ}a_0\mathrm{sin}(m_at)\mathrm{exp}(r/R_a),$$
(1)
where $`t`$ ( $`r`$ ) is time ( radial ) coordinate and $`f_{PQ}`$ is the decay constant of the axion. The value of $`f_{PQ}`$ is constrained conventinally from cosmological and astrophysical considerations such that $`10^{10}`$ GeV $`<f_{PQ}<`$ $`10^{12}`$ GeV ( the axion mass $`m_a`$ is given in terms of $`f_{PQ}`$ such that $`m_a10^7\text{GeV}^2/f_{PQ}`$ ). However, when we assume unconventionally entropy productions below the temperature $`1`$ GeV in the early Universe, we may be released from the constraints. Hereafter we assume that $`f_{PQ}10^{16}`$ GeV or $`m_a10^9`$ eV.
In the formula, $`R_a`$ represents the radius of an AS which has been obtained numerically in terms of mass $`M_a`$ of the AS; $`R_a=6.4m_{pl}^2/m_a^2M_a1.6\times 10^5\text{cm}m_9^2(10^1M_{}/M_a)`$, with $`m_9=m_a/10^9`$ eV and $`m_{pl}`$ is Planck mass. Similarly, the amplitude $`a_0`$ in eq(1) is represented such as $`a_0=1.73\times 10^2(10^5\text{cm}/R_a)^2m_9^1`$. Therefore, we find that the solution is parameterized by one free parameter, either one of the mass $`M_a`$ or the radius $`R_a`$. It is also important to note that the solution is not static but oscillating with the frequency of $`m_a/2\pi `$. It has been demonstrated that there is no static regular solution of the real scalar massless field coupled with gravity. This may be general property of the real scalar massive field.
The AS mass is determined by physical conditions under which the AS has been formed; how large amount of cloud of axions are cooled gravitationally to form the AS, etc.. The situation is similar to other stars such as neutron stars or white dwarfs. A typical mass scale in these cases is the critical mass; stars with masses larger than the critical mass collapse gravitationally into more compact ones or black holes. In the case of the AS, there also exists a critical mass $`M_c`$ which is given by $`M_c10^1M_{}m_9^1`$ where $`M_{}`$ represents the solar mass. Therefore, we adopt this critical mass $`M_c`$ as a typical mass scale of the ASs. ( The critical mass is the maximal mass which ASs can take. Thus, their masses, in general, are smaller than this one. Since energies released in GRBs by the ASs are given by masses themselves, the maximal energy in GRBs is given by the critical mass ).
Although the gravitational cooling for star formation is in general ineffective because it is too slow process, it has been shown that the cooling is very effective for the real scalar axion field. Thus, the axion stars can be easily formed gravitationally in a gas of the axions. It is reasonable to assume that the most of the axions in the Universe forms the axion stars.
Let us explain how an AS induces an electric field under a magnetic field $`\stackrel{}{B}`$ of a neutron star. Owing to the interaction between the axion and the electromagnetic field described by $`L_{a\gamma \gamma }=c\alpha a(x)\stackrel{}{E}\stackrel{}{B}/f_{PQ}\pi `$, where the value of $`c`$ ( of the order of unity ) depends on axion models, the Gauss law is modified such that $`\stackrel{}{}\stackrel{}{E}=c\alpha \stackrel{}{}(a(x)\stackrel{}{B})/f_{PQ}\pi +\text{“matter”}`$. The last term “matter” denotes electric charges of ordinary matters. The first term in the right hand side represents an electric charge made of the axion. ( This charge is oscillating and so there exist a corresponding oscillating current, $`J_a=c\alpha _ta(x)\stackrel{}{B}/f_{PQ}\pi `$. Thus, radiations are emitted by the AS, which might be observable. We will discuss it later. ) Thus, we find that the electric field $`\stackrel{}{E_a}`$ is induced; $`\stackrel{}{E_a}=c\alpha a(x)\stackrel{}{B}/f_{PQ}\pi `$ with $`\alpha 1/137`$, Numerically, its strength is given by
$$E_a10^{15}\text{eV}\text{cm}^1B_{12}m_9,$$
(2)
with $`B_{12}=B/10^{12}`$ G, where we have used the solution in eq(1) for the critical mass. Obviously, the spatial extention of the field is given by the radius $`R_a1.6\times 10^5m_9^1\text{cm}`$ of the AS.
This electric field is oscillating with the frequency, $`m_a/2\pi 2.4\times 10^5m_9\text{Hz}`$. Thus a particle with electric charge Ze can be accelerated in a direction within the half of the period, $`\pi /m_a`$ or in a distance $`R_a`$ ( $`\pi /m_a\times \text{light velocity}`$ ) by this field, unless it collides with other particles within the period. Thus, the energy $`\mathrm{\Delta }E`$ obtained by the particle is given by
$$\mathrm{\Delta }E=\text{Ze}E_a\times \pi /m_a\times \text{light velocity}10^{20}\text{Z}\text{eV}B_{12}.$$
(3)
Therefore, the electric field can accelerate the charged particle so that its energy reaches $`10^{20}\text{Z}`$ eV. These charged particles may be baryons or electron-positron pairs produced by the decay of the electric field itself as discussed soon below.
Comment is in order. It seems apparently from eq(3) that stronger magnetic fields yield cosmic rays with higher energies. Stronger magnetic fields, however, induce stronger electric fields, which are unstable against electron-positron pair creations. Therefore, strong magnetic fields do not necessarily yield cosmic rays with higher energies.
If the particles collide with other particles on the way of acceleration, in other words, their mean free paths are shorter than $`R_a`$, they can not obtain such high energies. It is easy to see, however, that the mean free paths of quarks or leptons with much higher energies than their masses are longer than $`R_a`$ in magnetosphere of neutron star. This is because since cross sections, $`\sigma `$, of quarks or leptons with such energies $`E`$ behaves such as $`\sigma 1/E^2`$, mean free paths $`1/n\sigma `$ is longer than $`R_a10^5`$ cm unless number density $`n`$ of particles around the AB is extremely large ( i.e. $`n>10^{40}/\text{cm}^3`$ ).
As is well known, the strong electric field is unstable against electron-positron pair creations. This implies that AS decays into the pairs when the electric field induced in AS is sufficiently strong.
Let us estimate the decay rate of the field and show that the AS decays before colliding with a neutron star whose magnetic field at surface is stronger than $`10^{12}`$ G. We also show that the AS can collide directly with a neutron star whose magnetic field is relatively weak $`10^{10}`$ G.
The decay rate $`R_d`$ of the field per unit volume and per unit time is given by
$$R_d=\frac{\alpha E_a^2}{\pi ^2}\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{exp}(m_e^2\pi n/eE_a)}{n^2}$$
(4)
where $`m_e`$ denotes electron mass. The rate is very small for an electric field much weaker than $`m_e^2\pi 4\times 10^{16}`$ eV/cm. The electric field of AS, however, can be very strong and it can be comparable to $`m_e^2\pi `$. Therefore, the rate is much large. Numerically, it is given by $`R_d10^{47}B_{12}^2m_9^2\text{cm}^3\text{s}^1_{n=1}^{\mathrm{}}\mathrm{exp}(0.7\times 10^2n/B_{12}m_9)/n^2`$. Since the spatial extention of the field is approximately given by $`10^5m_9^1`$ cm, the total decay rate $`W`$ of the field in AS is $`10^{62}B_{12}^2m_9^1\text{s}^1_{n=1}^{\mathrm{}}\mathrm{exp}(0.7\times 10^2n/B_{12}m_9)/n^2`$. Numerically, it reads
$$W10/\text{s}\text{for}B_{12}=0.5,W10^6/\text{s}\text{for}B_{12}=0.55,W10^{32}/\text{s}\text{for}B_{12}=1$$
(5)
with $`m_9=1`$.
Therefore, we find that the AS decays very rapidly ( or almost suddenly ) when it approaches a region where the strength of the magnetic field reaches a critical value of about $`10^{12}`$ G. Hence, the AS evapolates before colliding with the neutron star whose magnetic field at the surface is stronger than $`10^{12}`$ G. The whole energy of the AS is transmitted to the charged particles, each of which can obtain energies $`10^{20}`$ eV. These particles are emitted into a cone with very small solid angle. They form an extremely short pulse whose width being less than millisecond. Actually, when we suppose that the relative velocity of the AS is equal to light velocity, it decays approximately within a period of $`10^4\text{sec}\mathrm{\hspace{0.17em}\hspace{0.17em}10}^5`$sec; in the period it passes the region where the magnetic field increases from $`0.5\times 10^{12}`$ G to $`10^{12}`$ G. These leptons may be converted into baryons and photons through the interactions with themselves, interstellar medium or ejection of progenitor of the neutron star. Thus when the magnetic field is sufficiently strong, ultra high energy cosmic rays can be produced.
We comment that the velocity of AS trapped gravitationally to a neutron star, is approximately given by the light velocity just when it collides with the neutron star. This is because an AS can be trapped to a neutron star when the AS approaches it within a distance $`10^{11}`$ cm as we will show below, and then, the AS collides with the neutron star losing its potential energy and angular momentum by emitting gravitational waves. Thus, the velocity of the AS reaches approximately the light velocity when it collides with the neutron star.
Furthermore, we can see from eqs(2) and (4) that the critical electric field depends on the factor of $`B_{12}\times m_9`$. On the other hand, the energy $`\mathrm{\Delta }E`$ obtained by accelerated charged particles depends only on the factor of $`B_{12}`$. Therefore, we find that if $`m_a`$ is smaller than $`10^9`$ eV, the maximal energy of cosmic rays can be larger than $`10^{20}`$ eV when $`B>10^{12}`$ G.
It is easy to see that the decay rate is negligibly small for the case of weak magnetic field $`10^{10}`$ G. Thus the AS collides directly with such a neutron star and dissipates the whole energy in the outercrust of the neutron star. Actually, the rate of the energy dissipation in the crust has been estimated previously and given approximately by $`10^{46}\text{erg/s}\text{cm}^3`$, while the energy density of the AS is given by $`10^{38}\text{erg}/\text{cm}^3`$. Thus even if the velocity of the AS is equal to the light velocity, it dissipates its whole energy in the crust: it never reaches the core of the neutron star. This estimation has been performed by noting that the energy dissipation of the AS arises due to the dissipation of an electric current ( $`=\sigma _cE_a`$ ) induced in the crust with conductivity $`\sigma _c`$; the value of $`10^{26}`$/s has been used for $`\sigma _a`$.
Anyway, this collision generates gamma ray bursts with energies $`10^{53}`$ erg. The ejection could be emitted into a cone with small solid angle as a jet. This is because particles ( mainly irons ) of the neutron star are accelerated and emitted in the direction of the strong electric field $`10^{13}B_{10}`$ eV/cm. The fact that the whole energy of AB is dissipated only in the outercrust, implies that the ejection may be particles composing the crust. Thus a fraction of the baryon contamination in the jet is less than $`10^5M_{}`$ as required observationally.
In the above case, the AS dissipates its whole energy in the first collision. On the other hand, an AS may collide several times with a neutron star when its mass is much smaller than the critical mass $`10^1M_{}`$. The mass has been assumed as a typical mass of ABs. We see from the general formula of $`R_a`$ that the radius of the AS with smaller mass than the typical one is larger than the typical radius $`10^5`$ cm. For example, if its mass is given by $`10^2M_{}`$, the radius is about $`10^6`$ cm. This is comparable to the radius of the neutron star. In addition, an electric field induced in such an AS is weaker than the typical one eq(2). This implies that the rate of the energy dissipation in a neutron star is smaller than the one quoted above. In such a case the collisions might occur several times. We expect that these collisions leads to GRBs with pulses of longer duration and softer gamma rays. Observationally the only energies of GRBs with long duration have been measured. We predict that the energies of GRBs with short duration are much larger than those of GRBs with longer duration.
We now wish to estimate an energy release rate per unit volume and per unit time; the energy released as ultra high energy cosmic rays. In order to do so we assume that the dark matter is composed mainly of axion stars. The estimation, however, involves several ambiguities associated with number density of neutron stars, energy density of dark matter or velocity of ASs in the Universe etc. Therefore the estimation does not lead to a conclusive result although our result is consistent with the observations.
First we note that luminosity density around our galaxy is observationally given by $`2\times 10^8\text{h}L_{}/\text{Mpc}^3`$ where h is Hubble constant with the unit of $`100`$ km/s Mpc and $`L_{}`$ represents solar luminosity. We take a mass density $`\rho `$ corresponding to this luminosity density such as $`\rho 10^9M_{}\text{h}^2/\text{Mpc}^3`$. Then, we suppose that the number density of neutron stars is given by $`10^2\rho `$. This comes from the fact that the rate of appearance of supernovae is about $`0.11`$ per $`10`$ years in our galaxy and the age of the galaxy is about $`10^{10}`$ years. Probably, the rate of the appearance could be larger in early stage of our galaxy than the one at present. Thus, in our galaxy $`10^9`$ neutron stars might be involved corresponding to the number $`10^{11}`$ of stars in the galaxy. Hence we guess that the number of the neutron stars are equal to about $`10^2`$ times the number of the stars. All of these neutron stars are assumed to possess strong magnetic field $`>10^{12}`$ G. Furthermore, to estimate the rate of the energy release we need to know the average density $`ϵ`$ of the dark matter. Here we use a value of $`0.5\times 10^{24}\text{g}/\text{cm}^3`$, which represents a local density of our halo. Using these parameters we can estimate the rate owing to the collision between the AS and the neutron star if the cross section of the collision is found. The collision cross section for a neutron star to trap an AS is estimated in the following. Namely an AS is trapped by the neutron star when the AS approaches the neutrons star within a distance $`L_c`$ in which its kinetic energy $`M_av^2/2`$ is equal to its potential energy $`\mathrm{\hspace{0.17em}1.5}\times M_{}M_aG/L_c`$ around the neutron star. $`G`$ is gravitational constant. Here the mass of the neutron star and the relative velocity $`v`$ are assumed to be $`1.5\times M_{}`$ and $`3\times 10^7`$ cm/s, respectively. Thus, the cross section is found such as $`L_c^2\pi `$ with $`L_c6\times 10^{11}\text{cm}`$. After being trapped, the AS collides with the neutron star by losing its potential energy and angular momentum owing to the emission of gravitational and electromagnetic waves; electromagnetic radiations arise due to the oscillating current $`J_a`$. The merger is similar to neutron star - neutron star merger. Time scale from birth to merger is much less than the age of the Universe when the distance $`L_c`$ between two stars at the birth is given such that $`L_c10^{11}`$ cm.
Therefore, we obtain the rate of the collision per Mpc<sup>3</sup> and per year,
$$ϵ\times 10^2\rho \times L_c^2\pi \times v\times 1\text{year}/10^1M_{}3\times 10^9\text{h}^2/\text{Mpc}^3\text{year}.$$
(6)
Since the energy of $`10^{53}`$ erg is released in the collision, we find that the rate of the energy release is given by
$$3\times 10^{44}\text{h}^2\text{erg}/\text{Mpc}^3\text{year}.$$
(7)
which agrees well with the observed one . Taking account of the fact, however, that there are several ambiguities in the parameters used above and in the observations, we understand that our model can explains roughly the observations.
Finally we discuss two possibilities of the observation of the axion stars. One is associated with gravitational lensing and the other one associated with monochromatic radiations from the ASs.
Since we assume that the halo of our galaxy is composed of ASs and their mass is $`10^1M_{}`$, the ASs are plausible candidates for gravitational microlenses, i.e. MACHO. Since baryonic candidates like white dwarfs, neutron stars, etc. have serious problems, nonbaryonic candidates are favored. The problems are associated with chemical abundance of carbon and nitrogen in the Universe: If these baryonic stars are MACHOs, an overabundance of the elements is produced far in excess of what is observed in our galaxy. Hence, the ASs are theoretically the most fascinating candidates for MACHOs since they are also candidates for the generators of both UHECRs and GRBs. If the most favorable mass of the MACHO is $`0.5M_{}`$, we need to choose $`m_a0.2\times 10^9`$ eV since the mass of the ASs is given by $`10^1M_{}/m_9`$. Smaller axion mass leads to stronger electric field as well as larger mass of AS. Thus, it yields higher energies of the cosmic rays, $`5\times 10^{20}`$ eV and of GBRs, $`5\times 10^{53}`$ erg than the ones we have claimed above. Accordingly, the determination of the mass of MACHOs gives the upper limit of both the energies of the ultra high energy cosmic rays and the energies released in the gamma ray bursts in our mechanism.
We point out another way of the observation of the axion stars. Since the electric field induced in ASs is oscillating, electromagnetic radiations are emitted from corresponding oscillating electric current $`J_a`$ as mentioned above; the frequence of these monochromatic radiations is $`2.4\times 10^5m_9`$ Hz. We expect that such radiations can be detected in advance of UHECRs; they are emitted by the AS revolving around a neutron star before the AS decaying into charged particles. It is easy to estimate their luminocity $`6.7\times 10^{41}B_{12}^2\text{erg/s}`$. If the emission arises in a distance $`10`$ Mpc from the earth, we obtain the flux at the earth, $`10^9\text{Jy}B_{12}^2/m_9`$; we have assumed that the velocity of the AS revolving is $`0.1\times `$ light velocity. Although the possibility of observing the radiations is very intriguing, it might be difficult to detect the radiations with such a low frequency because they could be absorbed by interstellar ionized gases.
The author wishs to express his thank to all the staffs in Theory Group, Tanashi Branch, High Energy Accelerator Research Organization for their warm hospitality extended to him. This work is supported by the Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture of Japan No.10640284 |
warning/0003/cs0003015.html | ar5iv | text | # On the semantics of merging
## Introduction
To be able to operate in its environment it is necessary for an intelligent agent to have a consistent view of the world. This demand is often complicated by the fact that such agents receive conflicting pieces of information from different sources. The process of combining possibly inconsistent pieces of information, known as *merging*, has many applications and has started to receive more attention recently (?????????). In this paper we propose a framework for the modelling of merging operations. The proposal has its roots in the work of Spohn (??). Unlike most approaches we adopt a description of merging on the level of *epistemic states* instead of *knowledge bases*.
First we give a brief introduction to the merging of knowledge bases, focussing on the work of Konieczny and Pino-Pérez (?). This is followed by a description of our framework for the merging of epistemic states. Then we construct a number of merging operations and show how they measure up to proposed properties of merging. Finally, we discuss links between merging and the *infobases* of Meyer (?).
We assume a finitely generated propositional language $`L`$ closed under the usual propositional connectives, and with a classical model-theoretic semantics. $`U`$ is the set of interpretations of $`L`$ and $`M(\alpha )`$ is the set of models of $`\alpha L`$. Classical entailment is denoted by $``$. We use $``$ to denote the concatenation of lists. We let $`x^n`$ denote the list consisting of $`n`$ versions of $`x`$. The length of a list $`l`$ is denoted by $`\left|l\right|`$.
## Merging knowledge bases
In the spirit of the work of Katsuno and Mendelzon (?), approaches to the merging of knowledge bases usually represent the beliefs of an agent as a single wff $`\varphi `$ of $`L`$, known as a *knowledge base*, where $`\varphi `$ represents the set of all wffs entailed by $`\varphi `$. The goal is to construct, from a finite list of such knowledge bases, an appropriate consistent knowledge base in some rational fashion. Konieczny and Pino-Pérez (?) have proposed a general framework for the merging of knowledge bases. A *knowledge list* $`e`$ is a finite list of consistent knowledge bases $`[\varphi _1,\mathrm{},\varphi _{\left|e\right|}]`$. Two knowledge lists $`e_1`$ and $`e_2`$ are *element-equivalent*, written as $`e_1e_2`$, iff for every element $`\varphi _1`$ of $`e_1`$ there is a unique element $`\varphi _2`$ (position-wise) of $`e_2`$ such that $`\varphi _1\varphi _2`$ and for every element $`\varphi _2`$ of $`e_2`$ there is a unique element $`\varphi _1`$ (position-wise) of $`e_1`$ such that $`\varphi _2\varphi _1`$. A *KP-merging operation* $`\delta `$ is a function from the set of all knowledge lists to the set of all knowledge bases satisfying the following postulates (the KP-postulates):
$`\delta (e)`$
If $`_{i=1}^{\left|e\right|}\varphi _i`$ then $`\delta (e)=_{i=1}^{\left|e\right|}\varphi _i`$
If $`e_1e_2`$ then $`\delta (e_1)\delta (e_2)`$
If $`\varphi _1\varphi _2`$ then $`\delta ([\varphi _1][\varphi _2])\varphi _1`$
$`\delta (e_1)\delta (e_2)\delta (e_1e_2)`$
If $`\delta (e_1)\delta (e_2)`$ then $`\delta (e_1e_2)\delta (e_1)\delta (e_2)`$
Konieczny and Pino-Pérez also distinguish between two subclasses of merging operations. An *arbitration* operation tries to take as many differing opinions as possible into account, while the intuition associated with *majority* operations is that the opinion of the majority should prevail. They initially propose the following postulates for arbitration and majority operations.
$`n\delta (e\varphi ^n)=\delta (e[\varphi ])`$
$`n\delta (e\varphi ^n)\varphi `$
It turns out that there is no KP-merging operation satisfying (arb). Unlike Konieczny and Pino-Pérez we are of the opinion that it is not (arb) that is at fault, but some of the KP-postulates. Below we argue against the inclusion of (KP4) and (KP6) as postulates that need to be satisfied by all merging operations.
## Merging epistemic states
In this section we discuss merging on the level of epistemic states. We see an epistemic state as providing a plausibility ranking of the interpretations of $`L`$; the lower the number assigned to an interpretation, the more plausible it is deemed to be.
###### Definition 0.1
An epistemic state $`\mathrm{\Phi }`$ is a function from $`U`$ to the set of natural numbers. Given an epistemic state $`\mathrm{\Phi }`$, the knowledge base associated with $`\mathrm{\Phi }`$, denoted by $`\varphi _\mathrm{\Phi }`$, is some $`\alpha L`$ such that $`M(\alpha )=\{u\mathrm{\Phi }(u)=0\}`$. $`\mathrm{}`$
This representation of an epistemic state and its associated knowledge base can be traced back to the work of Spohn (??). It should be clear that an epistemic state with an inconsistent associated knowledge base still contains useful information.
An *epistemic list* $`E=[\mathrm{\Phi }_1^E,\mathrm{},\mathrm{\Phi }_{\left|E\right|}^E]`$ is a finite list of epistemic states.
It is instructive to view an epistemic list pictorially as in figure 1. While such a pictorial view is only useful in representing epistemic lists containing two elements, it serves as a good foundation for understanding the principles underlying the merging of epistemic states in general.
For any epistemic state $`\mathrm{\Phi }`$, let
$$\mathrm{min}(\mathrm{\Phi })=\mathrm{min}\{\mathrm{\Phi }(u)uU\}$$
let
$$\mathrm{max}(\mathrm{\Phi })=\mathrm{max}\{\mathrm{\Phi }(u)uU\}$$
and for an epistemic list $`E`$, let
$$\mathrm{max}(E)=\mathrm{max}\{\mathrm{max}(\mathrm{\Phi }_i^E)1i\left|E\right|\}.$$
For an epistemic list $`E`$ and $`uU`$ we let $`\mathrm{min}^E(u)`$ be equal to
$$\mathrm{min}\{\mathrm{\Phi }_i^E(u)1i\left|E\right|\}$$
and we let $`\mathrm{max}^E(u)`$ be equal to
$$\mathrm{max}\{\mathrm{\Phi }_i^E(u)1i\left|E\right|\}.$$
We denote by $`seq(E)`$ the set of all sequences of length $`\left|E\right|`$ of natural numbers, ranging from $`0`$ to $`\mathrm{max}(E)`$. We denote by $`seq_{}(E)`$ the subset of $`seq(E)`$ of all sequences that are ordered non-decreasingly, and by $`seq_{}(E)`$ the subset of $`seq(E)`$ of all sequences that are ordered non-increasingly. For $`uU`$, we let $`s^E(u)`$ be the sequence containing the natural numbers $`\mathrm{\Phi }_1^E(u),\mathrm{},\mathrm{\Phi }_{\left|E\right|}^E(u)`$ in that order, we let $`s_{}^E(u)`$ be the sequence $`s^E(u)`$ ordered non-decreasingly, and we let $`s_{}^E(u)`$ be the sequence $`s^E(u)`$ ordered non-increasingly. Clearly $`s^E(u)seq(E)`$, $`s_{}^E(u)seq_{}(E)`$ and $`s_{}^E(u)seq_{}(E)`$. Given any set $`seq`$ of finite sequences of natural numbers and a total preorder $``$ on $`seq`$, we define the function $`\mathrm{\Omega }_{}^{seq}:seq\{0,\mathrm{},\left|seq\right|1\}`$ by assigning natural numbers to the elements of $`seq`$ in the order imposed by $``$, starting by assigning $`0`$ to the elements lowest down in $``$. We denote the *lexicographic* ordering on $`seq`$ by $`_{lex}`$.
A *merging operation on epistemic states* $`\mathrm{\Delta }`$ is a function from the set of all non-empty epistemic lists to the set of all epistemic states. We propose the following basic properties for the merging of epistemic states:
$`u`$ s.t. $`\mathrm{\Delta }(E)(u)=0`$
If $`\mathrm{\Phi }_i^E(u)=\mathrm{\Phi }_j^E(u)i,j`$ such that $`1i,j\left|E\right|`$ and $`s_{}^E(u)_{lex}s_{}^E(v)`$ then $`\mathrm{\Delta }(E)(u)<\mathrm{\Delta }(E)(v)`$
If $`\mathrm{\Phi }_i^E(u)\mathrm{\Phi }_i^E(v)i`$ such that $`1i\left|E\right|`$ then $`\mathrm{\Delta }(E)(u)\mathrm{\Delta }(E)(v)`$
If $`\mathrm{\Delta }(E)(u)\mathrm{\Delta }(E)(v)`$ then $`\mathrm{\Phi }_i^E(u)\mathrm{\Phi }_i^E(v)`$ for some $`i`$ such that $`1i\left|E\right|`$
(E1) is a restatement of (KP1) and (E2) generalises (KP2). (E3) states that if all epistemic states in $`E`$ agree that $`u`$ is at least as plausible as $`v`$, then so should the resulting epistemic state. (E4) expects justification for regarding an interpretation $`u`$ as at least as plausible as $`v`$: there has to be at least one epistemic state in $`E`$ which regards $`u`$ as at least as plausible as $`v`$. The following fundamental principle for the merging of epistemic states follows easily from (E3):
If $`\mathrm{\Phi }_i^E(u)=\mathrm{\Phi }_i^E(v)i`$ such that $`1i\left|E\right|`$ then $`\mathrm{\Delta }(E)(u)=\mathrm{\Delta }(E)(v)`$
(Unit) requires interpretations that are treated identically by all epistemic states in an epistemic list to be treated identically in the epistemic state resulting from a merging operation.
Two epistemic lists $`E_1`$ and $`E_2`$ are *element-equivalent*, written as $`E_1E_2`$, iff for every element $`\mathrm{\Phi }_1`$ of $`E_1`$ there is a unique element $`\mathrm{\Phi }_2`$ (position-wise) of $`E_2`$ such that $`\mathrm{\Phi }_1=\mathrm{\Phi }_2`$ and for every element $`\mathrm{\Phi }_2`$ of $`E_2`$ there is a unique element $`\mathrm{\Phi }_1`$ (position-wise) of $`E_1`$ such that $`\mathrm{\Phi }_2=\mathrm{\Phi }_1`$. The following property is a generalisation of (KP3). It requires merging to be commutative.
$`E_1E_2`$ implies $`\mathrm{\Delta }(E_1)=\mathrm{\Delta }(E_2)`$
We do not think that (Comm) should hold for all merging operations. Instead, (Comm) should be seen as a postulate picking out an interesting subclass of merging operations.
For a finite list of epistemic lists $`=[E_1,\mathrm{},E_{\left|\right|}]`$, let $`\mathrm{\Delta }()`$ denote the epistemic list $`[\mathrm{\Delta }(E_1),\mathrm{},\mathrm{\Delta }(E_{\left|\right|})]`$. We consider the following properties:
If $`\mathrm{\Delta }(E_i)(u)\mathrm{\Delta }(E_i)(v)i`$ such that $`1i\left|\right|`$ then $`\mathrm{\Delta }(_{i=1}^{\left|\right|}E_i)(u)\mathrm{\Delta }(_{i=1}^{\left|\right|}E_i)(v)`$
If $`\mathrm{\Delta }(_{i=1}^{\left|\right|}E_i)(u)\mathrm{\Delta }(_{i=1}^{\left|\right|}E_i)(v)`$ then for some $`i`$ such that $`1i\left|\right|`$, $`\mathrm{\Delta }(E_i)(u)\mathrm{\Delta }(E_i)(v)`$ for some $`1i\left|\right|`$
(E5) generalises (E3) and (E6) generalises (E4). In fact, (E5) also implies (KP5).
The arbitration postulate (arb) and the majority postulate (maj) can be generalised as follows:
$`n\mathrm{\Delta }(E[\mathrm{\Phi }])(u)=\mathrm{\Delta }(E\mathrm{\Phi }^n)(u)`$
$`n`$ s.t. $`u,vU,\mathrm{\Phi }(u)\mathrm{\Phi }(v)`$ if $`\mathrm{\Delta }(E\mathrm{\Phi }^n)(u)\mathrm{\Delta }(E\mathrm{\Phi }^n)(v)`$
We have not provided a generalised version of (KP4). The reason is that we do not regard it as a suitable postulate for merging. Our basic argument is that the models of a knowledge base associated with an epistemic state $`\mathrm{\Phi }_1`$ may sometimes be given such an implausible ranking by an epistemic state $`\mathrm{\Phi }_2`$ that it would seem reasonable to exclude all these models from the models of $`\varphi _{\mathrm{\Delta }([\mathrm{\Phi }_1][\mathrm{\Phi }_2])}`$. It is worthwhile noting that none of the merging operations we consider below satisfies (KP4). Similarly, we have not provided a generalised version of (KP6) since we regard it as too strong a condition to impose on all merging operations.<sup>1</sup><sup>1</sup>1(E6) can be regarded as a generalised version of a weaker form of (KP6), but (KP6) does not follow from (E6). Below we shall encounter a number of reasonable merging operations which do not satisy (KP6).
## Constructing merging operations
Konieczny and Pino-Pérez (?) discuss several merging operations on knowledge bases using Dalal’s measure of distance between interpretations (?). For any two interpretations $`u`$ and $`v`$, let $`dist(u,v)`$ denote the number of propositional atoms on which $`u`$ and $`v`$ differ. The distance $`Dist(\varphi ,u)`$ between a knowledge base $`\varphi `$ and an interpretation $`u`$ is defined as follows: $`Dist(\varphi ,u)=\mathrm{min}\{dist(u,v)vM(\varphi )\}`$. It is clear that this distance measure can be used to define an epistemic state $`\mathrm{\Phi }`$ as follows:
$$uU,\mathrm{\Phi }(u)=Dist(\varphi ,u).$$
It is easily seen that $`\mathrm{\Phi }(u)=0`$ iff $`uM(\varphi )`$ and therefore $`\varphi _\mathrm{\Phi }\varphi `$. Many of the merging operations on epistemic states that we propose below are appropriate generalisations of these merging operations on knowledge bases.
When reading through the remainder of this section, the reader should observe that the construction of every merging operation consists of two steps. In the first step natural numbers are assigned to interpretations. After the completion of this step it will often be the case that *none* of the interpretations have been assigned the value $`0`$. To ensure compliance with (E1) the second step performs an appropriate uniform subtraction of values which we shall refer to as *normalisation*.
### Arbitration
Inspired by an arbitration operation proposed by Liberatore and Schaerf (?) we propose the following two merging operations on epistemic states.
###### Definition 0.2
1. Let $`\mathrm{\Phi }_{ls}^E(u)=2\mathrm{min}^E(u)`$ if $`\mathrm{\Phi }_i^E(u)=\mathrm{\Phi }_j^E(u)`$ for $`1i,j\left|E\right|`$, and $`\mathrm{\Phi }_{ls}^E(u)=2\mathrm{min}^E(u)+1`$ otherwise. Then $`\mathrm{\Delta }_{ls}(E)(u)=\mathrm{\Phi }_{ls}^E(u)\mathrm{min}(\mathrm{\Phi }_{ls}^E)`$.
2. Let $`\mathrm{\Phi }_{Rls}^E(u)=\mathrm{\Omega }_{_{lex}}^{seq_{}(E)}(s_{}^E(u))`$. Then $`\mathrm{\Delta }_{Rls}(E)(u)=\mathrm{\Phi }_{Rls}^E(u)\mathrm{min}(\mathrm{\Phi }_{Rls}^E)`$.
$`\mathrm{}`$
Figure 2 contains a pictorial representation of $`\mathrm{\Delta }_{ls}`$ and figure 3 a pictorial representation of $`\mathrm{\Delta }_{Rls}`$. It can easily be shown that $`\mathrm{\Delta }_{Rls}`$ is a refined version of $`\mathrm{\Delta }_{ls}`$. Both satisfy (E1)-(E6) and (Comm), neither satisfies (Maj), and only $`\mathrm{\Delta }_{Rls}`$ satisfies (KP6). Moreover, $`\mathrm{\Delta }_{ls}`$ satisfies (Arb) but $`\mathrm{\Delta }_{Rls}`$ does not. So, while both are valid merging operations, $`\mathrm{\Delta }_{Rls}`$ should not be seen as an arbitation operation.
Next we consider two merging operations that are generalisations of the $`\delta _{\mathrm{max}}`$ and $`\delta _{Gmax}`$ operations of Konieczny and Pino-Pérez. The former was inspired by an example of Revesz’s model-fitting operations (?).
###### Definition 0.3
1. Let $`\mathrm{\Phi }_{\mathrm{max}}^E(u)=\mathrm{max}^E(u)`$. Then $`\mathrm{\Delta }_{\mathrm{max}}(E)(u)=\mathrm{\Phi }_{\mathrm{max}}^E(u)\mathrm{min}(\mathrm{\Phi }_{\mathrm{max}}^E)`$.
2. Let $`\mathrm{\Phi }_{Gmax}^E(u)=\mathrm{\Omega }_{_{lex}}^{seq_{}(E)}(s_{}^E(u))`$. Then $`\mathrm{\Delta }_{Gmax}(E)(u)=\mathrm{\Phi }_{Gmax}^E(u)\mathrm{min}(\mathrm{\Phi }_{Gmax}^E)`$.
$`\mathrm{}`$
Figure 4 contains a pictorial representation of $`\mathrm{\Delta }_{\mathrm{max}}`$ and figure 5 a pictorial representation of $`\mathrm{\Delta }_{Gmax}`$. Both satisfy (E1)-(E6), neither satisfies (Maj), and only $`\mathrm{\Delta }_{Gmax}`$ satisfies (KP6). Moreover, $`\mathrm{\Delta }_{\mathrm{max}}`$ satisfies (Arb), but $`\mathrm{\Delta }_{Gmax}`$ does not. So, analogous to the case above, both are valid merging operations but $`\mathrm{\Delta }_{Gmax}`$ should not be seen as an arbitation operation. The fact that we do not regard $`\mathrm{\Delta }_{Gmax}`$ as an arbitration operation is in conflict with the view of Konieczny and Pino-Pérez who regard $`\delta _{Gmax}`$ as an arbitration operation on knowledge bases even though it does not satisfy (arb). Conversely, Konieczny and Pino-Pérez do not regard $`\delta _{\mathrm{max}}`$ as a merging operation on knowledge bases since it fails to satisfy (KP6). But we regard it as a valid arbitration operation since it satisfies the postulates (E1)-(E6), (Comm) and (Arb).
### Consensus
In this section we consider the idea of a *consensus* operation, where agreement on the ranking of interprerations, instead of the ranking itself, is of overriding importance.
###### Definition 0.4
For $`sseq(E)`$, let
$$d^E(s)=\underset{i=1}{\overset{\left|E\right|}{}}\underset{j=i+1}{\overset{\left|E\right|}{}}\left|s_is_j\right|$$
where $`s_i`$ denotes the $`i`$th element of $`s`$.
1. Define the total preorder $``$ on $`seq(E)`$ as follows: $`st`$ iff $`d^E(s)d^E(t)`$. Let $`\mathrm{\Phi }_{cons}^E(u)=\mathrm{\Omega }_{}^{seq(E)}(s^E(u))`$. Then $`\mathrm{\Delta }_{cons}(E)(u)=\mathrm{\Phi }_{cons}^E(u)\mathrm{min}(\mathrm{\Phi }_{cons}^E)`$.
2. Define the total preorder $``$ on $`seq_{}(E)`$ as follows: $`st`$ iff $`d^E(s)<d^E(t)`$ or ($`d^E(s)=d^E(t)`$ and $`s_{lex}t`$). Now, let $`\mathrm{\Phi }_{Rcons}^E(u)=\mathrm{\Omega }_{}^{seq_{}(E)}(s_{}^E(u))`$. Then $`\mathrm{\Delta }_{Rcons}(E)(u)=\mathrm{\Phi }_{Rcons}^E(u)\mathrm{min}(\mathrm{\Phi }_{Rcons}^E)`$.
$`\mathrm{}`$
Figure 6 contains a pictorial representation of $`\mathrm{\Delta }_{cons}`$ and figure 7 a pictorial representation of $`\mathrm{\Delta }_{Rcons}`$. We do not regard these two operations as suitable candidates for merging, primarily because both fail to satisfy (E3) and (E4). Both satisfy (Unit), though. The problem with these consensus operations seems to be that they place too strong an emphasis on agreement and do not take the ranking of interpretations seriously enough.
### Majority
We consider the following two majority operations.
###### Definition 0.5
For $`sseq(E)`$, let
$$sum^E(s)=\underset{i=1}{\overset{\left|E\right|}{}}s_i$$
where $`s_i`$ is the $`i`$th element of $`s`$.
1. Let $`\mathrm{\Phi }_\mathrm{\Sigma }^E(u)=sum^E(s^E(u))`$. Then $`\mathrm{\Delta }_\mathrm{\Sigma }(E)(u)=\mathrm{\Phi }_\mathrm{\Sigma }^E(u)\mathrm{min}(\mathrm{\Phi }_\mathrm{\Sigma }^E)`$.
2. Define the total preorder $``$ on $`seq(E)`$ as follows:
$`st`$ iff $`sum^E(s)<sum^E(t)`$ or
($`sum^E(s)=sum^E(t)`$ and $`d^E(s)d^E(t)`$). Now, let $`\mathrm{\Phi }_{R\mathrm{\Sigma }}^E(u)=\mathrm{\Omega }_{}^{seq(E)}(s^E(u))`$. Then $`\mathrm{\Delta }_{R\mathrm{\Sigma }}(E)(u)=\mathrm{\Phi }_{R\mathrm{\Sigma }}^E(u)\mathrm{min}(\mathrm{\Phi }_{R\mathrm{\Sigma }}^E)`$.
$`\mathrm{}`$
Figure 8 contains a pictorial representation of $`\mathrm{\Delta }_\mathrm{\Sigma }`$ and figure 9 a pictorial representation of $`\mathrm{\Delta }_{R\mathrm{\Sigma }}`$. $`\mathrm{\Delta }_\mathrm{\Sigma }`$ is an appropriate generalisation of an example by Lin and Mendelzon (?). It was independently proposed by Revesz (?) as an example of weighted model fitting. The idea is simply to obtain the new plausibility ranking of an interpretation by summing the plausibility rankings given by the different epistemic states. $`\mathrm{\Delta }_{R\mathrm{\Sigma }}`$ is $`\mathrm{\Delta }_\mathrm{\Sigma }`$ refined by using consensus. Both $`\mathrm{\Delta }_\mathrm{\Sigma }`$ and $`\mathrm{\Delta }_{R\mathrm{\Sigma }}`$ satisfy (E1)-(E4), (Comm) and (Maj), and neither satisfies (Arb). But while $`\mathrm{\Delta }_\mathrm{\Sigma }`$ satisfies (E5)-(E6) and (KP5)-(KP6) as well, $`\mathrm{\Delta }_{R\mathrm{\Sigma }}`$ does not.
### Non-commutative merging
Thus far we have restricted ourselves to the construction of *commutative* merging operations – i.e., satisfying (Comm) – but a complete description of merging ought to take into account constructions such as that of Nayak (?), in which the merging of two epistemic states is obtained by a lexicographic refinement of one by the other. We present here a generalised version of Nayak’s proposal. For this case the epistemic states in an epistemic list are assumed to be ranked according to reliability. That is, given an epistemic list $`E=[\mathrm{\Phi }_1^E,\mathrm{},\mathrm{\Phi }_{\left|E\right|}^E]`$, $`\mathrm{\Phi }_i^E`$ is at least as reliable as $`\mathrm{\Phi }_j^E`$ iff $`ij`$.
###### Definition 0.6
Let $`\mathrm{\Phi }_{lex}^E(u)=\mathrm{\Omega }_{_{lex}}^{seq(E)}(s^E(u))`$. Then $`\mathrm{\Delta }_{lex}(E)(u)=\mathrm{\Phi }_{lex}^E(u)\mathrm{min}(\mathrm{\Phi }_{lex}^E)`$. $`\mathrm{}`$
$`\mathrm{\Delta }_{lex}`$ does not satisfy (Comm), but it satisfies (E1)-(E6), as well as (KP5)-(KP6). By exploiting the non-commutativity of $`\mathrm{\Delta }_{lex}`$, both (Arb) and (Maj) can be phrased in a way to ensure that $`\mathrm{\Delta }_{lex}`$ fails to satisfy them.
## Merging and infobases
Our description of merging uses a representation of epistemic states as functions assigning a plausibility ranking to the interpretations of $`L`$, but where do these plausibility rankings come from? One way in which to generate them is by using the *infobases* of Meyer (?). An infobase is a finite list of wffs. Intuitively it is a structured representation of the beliefs of an agent with a foundational flavour. It is assumed that every wff in an infobase is obtained independently. Meyer uses an infobase to define a total preorder on $`U`$, which is then used to perform belief change. However, we can also use an infobase to define an epistemic state. The idea is to consider the number of times that an interpretation occurs as a model of one of the wffs in an infobase: the more it occurs, the higher its plausibility ranking.
###### Definition 0.7
For $`uU`$, define the $`\mathrm{𝐼𝐵}`$-number $`u_{\mathrm{𝐼𝐵}}`$ of $`u`$ as the number of elements $`\alpha `$ in an infobase $`\mathrm{𝐼𝐵}`$ such that $`\alpha `$ and $`uM(\alpha )`$, and let
$$\mathrm{max}(\mathrm{𝐼𝐵})=\mathrm{max}\{u_{\mathrm{𝐼𝐵}}uU\}.$$
Now we define the epistemic state associated with $`\mathrm{𝐼𝐵}`$ as follows: for $`uU,\mathrm{\Phi }^{\mathrm{𝐼𝐵}}(u)=\mathrm{max}(\mathrm{𝐼𝐵})u_{\mathrm{𝐼𝐵}}`$. $`\mathrm{}`$
Observe that the knowledge base associated with an epistemic state $`\mathrm{\Phi }^{\mathrm{𝐼𝐵}}`$ is always consistent, regardless of whether the wffs in $`\mathrm{𝐼𝐵}`$ are jointly consistent. We show that infobases seem to provide a natural setting in which to apply merging.
Firstly, define an *infobase list* $`\mathrm{𝐸𝐵}=[\mathrm{𝐼𝐵}_1,\mathrm{},\mathrm{𝐼𝐵}_{\left|\mathrm{𝐸𝐵}\right|}]`$ as a finite non-empty list of infobases and let $`E^{\mathrm{𝐸𝐵}}`$ denote the epistemic list $`[\mathrm{\Phi }^{\mathrm{𝐼𝐵}_1},\mathrm{},\mathrm{\Phi }^{\mathrm{𝐼𝐵}_{\left|E\right|}}]`$ of epistemic states associated with the infobases occurring in $`\mathrm{𝐸𝐵}`$. Then it can be verified that $`\mathrm{\Delta }_\mathrm{\Sigma }(E^{\mathrm{𝐸𝐵}})=\mathrm{\Phi }^{\mathrm{𝐼𝐵}}`$ where $`\mathrm{𝐼𝐵}=_{i=1}^{\left|\mathrm{𝐸𝐵}\right|}\mathrm{𝐼𝐵}_i`$.
Secondly, Konieczny and Pino-Pérez (?) give a convincing example to show that we may sometimes want to include, as models of $`\delta (e)`$, interpretations other than the models of the knowledge bases in $`e`$. Below is a scaled down version of their example.
###### Example 0.8
We want to speculate on the stock exchange and we ask two equally reliable financial experts about two shares. Let the atom $`p`$ denote the fact that share 1 will rise and $`q`$ the fact that share 2 will rise. The first expert says that both shares will rise: $`\varphi _1=pq`$, while the second one believes that both shares will fall: $`\varphi _2=\neg p\neg q`$. Intuitively it seems reasonable to conclude that both experts are right (and wrong) about exactly one share, although we don’t know which share in either case. That is, we require the result of the merging of these two knowledge bases to be such that $`M(\delta ([\varphi _1][\varphi _2]))=\{10,01\}`$.<sup>2</sup><sup>2</sup>2We represent interpretations as sequences consisting of 0s (representing falsity) and 1s (representing truth), where the first digit in a sequence represents the truth value of $`p`$ and the second one the truth value of $`q`$. Observe that $`M(\delta ([\varphi _1][\varphi _2]))M(\varphi _1)M(\varphi _2)`$. $`\mathrm{}`$
An analysis of this example shows that both experts are assumed to make an implicit assumption of independence of the performance of the shares. Thus the beliefs of the first expert is best expressed as the infobase $`\mathrm{𝐼𝐵}_1=[p,q]`$ and the beliefs of the second expert as the infobase $`\mathrm{𝐼𝐵}_2=[\neg p,\neg q]`$. The epistemic states obtained from these two infobases are: $`\mathrm{\Phi }^{\mathrm{𝐼𝐵}_1}(11)=0,\mathrm{\Phi }^{\mathrm{𝐼𝐵}_1}(10)=\mathrm{\Phi }^{\mathrm{𝐼𝐵}_1}(01)=1,\mathrm{\Phi }^{\mathrm{𝐼𝐵}_1}(00)=2`$, and $`\mathrm{\Phi }^{\mathrm{𝐼𝐵}_2}(00)=0,\mathrm{\Phi }^{\mathrm{𝐼𝐵}_2}(10)=\mathrm{\Phi }^{\mathrm{𝐼𝐵}_2}(01)=1,\mathrm{\Phi }^{\mathrm{𝐼𝐵}_2}(11)=2`$. It can be verified that $`\mathrm{\Delta }_{\mathrm{max}}(E^{\mathrm{𝐸𝐵}})=\mathrm{\Delta }_{Gmax}(E^{\mathrm{𝐸𝐵}})=\mathrm{\Delta }_{R\mathrm{\Sigma }}(E^{\mathrm{𝐸𝐵}})=\mathrm{\Phi }`$, where $`\mathrm{𝐸𝐵}=[\mathrm{𝐼𝐵}_1,\mathrm{𝐼𝐵}_2]`$, $`\mathrm{\Phi }(10)=\mathrm{\Phi }(01)=0`$ and $`\mathrm{\Phi }(11)=\mathrm{\Phi }(00)=1`$. So $`\mathrm{\Delta }_{R\mathrm{\Sigma }}`$, $`\mathrm{\Delta }_{\mathrm{max}}`$ and $`\mathrm{\Delta }_{Gmax}`$ yield the results corresponding to our intuition for this example.
## Conclusion
The merging operations we have constructed provide evidence that (E1)-(E4) may be regarded as basic postulates for merging operations on epistemic states. Furthermore, we regard (Arb) as an appropriate postulate for the subclass of arbitration operations, (Maj) for the subclass of majority operations, and (Comm) for the subclass of commutative merging operations. The status of (E5) and (E6) is less clear. While all but one of the valid merging operations we have considered satisfy both, the fact that $`\mathrm{\Delta }_{R\mathrm{\Sigma }}`$ does not, suggests that they are not as universally applicable as (E1)-(E4). Perhaps they should be seen as picking out particular subclasses of merging operations in the way that (Arb), (Maj) and (Comm) do. |
warning/0003/astro-ph0003104.html | ar5iv | text | # The Role of Heating and Enrichment in Galaxy Formation
## 1. Introduction
It has long been recognized that the X-ray luminosity-temperature ($`L_xT`$) relation of clusters does not obey the simple scaling laws that would hold if clusters were formed from the collapse of unheated primordial gas, and thus the gas within clusters is likely to have been heated before the formation of the clusters themselves (Kaiser 1991). Subsequent investigations have determined that preheating is also necessary to explain the $`L_xT`$ relation as a function of cluster mass (Cavaliere, Menci, & Tozzi 1999).
The high metallicity of cluster gas and the claimed over-abundance of alpha elements (Gibson and Matteucci 1996; Lowenstein & Mushotzky 1996) points to preheating by SNII driven winds (Renzini et al. 1993; Trentham 1994; Nath & Chiba 1995). High enrichment is not restricted to the most massive clusters, but appears to be widespread, extending to groups of galaxies (Buote 2000). Empirically the epoch of this enrichment is now known to occur at $`z>0.5`$, and approximately solar enriched cluster gas has been detected at redshifts as large as $`z1`$ (Hattori et al. 1997). Hence it is natural to suppose that pre-enrichment and preheating are the consequence of very early and vigorous massive star-formation.
Cluster gas is evidently a sink for enriched and heated material, and the most likely culprit for this enrichment is dwarf galaxies. Theoretical work has shown that supernovae and OB winds in these low-mass objects should lead to the production of energetic outflows with temperatures on the order of $`10^6`$ K (Larson 1974; Dekel & Silk 1986; Vader 1986). This behavior has been clearly identified in studies of both local starbursting galaxies (Axon & Taylor 1978; Marlowe et al. 1995; Heckman 1997; Hunter et al. 1998; Martin 1998) and spectroscopy of high-$`z`$ galaxies ( Franx et al. 1997; Pettini et al. 1998; Frye & Broadhurst 1998; Warren et al. 1998). Whether these outflows lead to a catastrophic loss of the interstellar gas however, is likely to depend on a number of factors (De Young & Heckman 1994), and is a subject of current investigations (Murakami & Babul 1999; Mac Low & Ferrara 1999; Strickland & Stevens 1999).
In the generally investigated hierarchical models of structure formation such as the Cold Dark Matter model (CDM), the existence of an era of widespread enrichment by outflows from dwarf galaxies is in fact quite natural, because low-mass galaxies are expected to form in large numbers and at early times (e.g., White & Frenk 1991). Since these early galaxies will be found preferentially in the large-scale overdense regions that later form clusters, dwarf outflows are the obvious candidates for pre-cluster heating.
The existence of large numbers of small galaxies at high redshifts is also favored by observations, to help understand the steep number counts and low luminosities of faint galaxies (Broadhurst, Ellis, & Glazebrook 1992), the sizes of faint galaxies in Hubble Deep Field images (Bouwens, Broadhurst, & Silk 1998a,b), and the small sizes of the distant Lyman-break galaxies (Steidel et al. 1999). Locally, however, the space density of dwarf galaxies relative to massive galaxies is far less than predicted on the basis of the steep Press-Schechter slope for the faint end of the mass function (Ferguson & Binggeli 1994), prompting theoretical studies of the disruption of dwarves by tidal forces from neighboring objects (Moore et al. 1998), external UV radiation (Kepner, Babul, & Spergel 1997; Norman & Spaans 1997; Corbelli, Galli, & Palla 1997), and catastrophic mass loss during outflows (Larson 1974; Dekel & Silk 1986; Vader 1986).
While many of these mechanisms may have had an impact on the formation of dwarf galaxies, relatively little attention has been directed towards the influence of outflows on neighboring galaxies. As the earliest galaxies to form were highly clustered (Kaiser 1984) and typical outflow temperatures and velocities were much larger than the virial temperatures and velocities of these galaxies, it is likely that dwarf galaxies were strongly influenced by their neighbors (Scannapieco, Ferrara, & Broadhurst 2000). Similarly, the sources responsible for pre-enrichment of the intracluster medium may well have enriched larger protogalaxies, with important consequences for their metallicity histories and cooling times.
In this work, we conduct an idealized investigation of the impact of dwarf outflows on the history of galaxy formation. While the consequences of homogeneous heating on galaxy formation have been examined in the past (Blanchard, Valls-Gabaud, & Mamon 1992), the inhomogeneous nature of this process and the associated pre-enrichment of galaxies have not been addressed. Our aim is not, however, to construct a complete model of galaxy formation, and thus we do not track processes such as the formation of second-generation stars in galaxies, production of dust, transfer of angular momentum, or the structure of the interstellar medium. Rather we focus on the properties of the inhomogeneously heated and enriched intergalactic medium (IGM) out of which galaxies coalesced, the likely consequences of galaxy formation in this environment, and to what degree these issues must be accounted for within more detailed simulations of galaxy formation.
The structure of this work is as follows. In §2 we review the observational evidence that leads us to consider a model in which widespread dwarf outflows shocked and enriched the medium out of which larger galaxies formed. In §3 we describe a simple numerical Monte Carlo code that we use to asses the overall features of such a model. The results of our simulations are given in §4 in which we examine which aspects of galaxy formation are most sensitive to the presence of outflows. In §5 we discuss the limitations of our modeling, and conclusions are listed in §6.
## 2. Observational Evidence
In this section we review the observational evidence that points towards widespread early enrichment by dwarf galaxies. We show that the chemical and thermal properties of nearby galaxy clusters are both suggestive of such pre-enrichment, as is the weak evolution of these properties with redshift. We show that such a picture is consistent with our knowledge of the intergalactic medium as well, and helps to explain the relatively high metallicity seen in many Ly absorption systems. The presence of outflows has been identified in emission line studies of high-$`z`$ galaxies, and several nearby dwarf galaxies have been caught “red-handed”, surrounded by clouds of enriched gas with temperatures greatly exceeding virial. Finally, we review some unsolved questions in galaxy formation that may depend on preheating and enrichment and serve to motivate our exploratory numerical studies.
### 2.1. Galactic Outflows and the Intra Cluster Medium
Many cluster properties can be understood in the context of self-similar models (Kaiser 1986, 1990) which are exact for power-law initial fluctuation spectra and a reasonable approximation for more realistic spectra such as that of the Cold Dark Matter (CDM) model. While the optical properties of galaxies in clusters are consistent with these predictions, the slope of the X-ray luminosity-temperature relationship of the hot intra-cluster gas is too steep and slowly evolving to be understood in this context.
A convincing explanation of $`L_xT`$ evolution has been made by Kaiser (1991) who speculated that at an epoch before the formation of present-day clusters, the gas which formed the intracluster medium (ICM) was preheated, injecting sufficient energy to expel gas from the small potential wells that were non-linear at early times. By assuming that this heated gas later cooled adiabatically onto the large cluster-size potential wells, he was able to simply reproduce the observed X-ray properties of clusters. Numerous investigations of preheating and $`L_xT`$ evolution have also reached similar conclusions (David, Jones, & Forman, 1996; Mushotzky & Scharf 1997; Eke, Navarro, Frenk 1998) as have studies of the $`L_xT`$ relation in groups of galaxies (Cavaliere, Menci, & Tozzi 1999).
The origin of this preheating is most naturally due to Supernovae Type-II (SNeII) activity in the dwarf starburst population. Several studies of expanding HI gas in nearby dwarf galaxies show clear evidence of dense expanding shells with velocities above $`15`$ km/s (Marlowe et al. 1995; Heckman 1997; Hunter et al. 1998; Martin 1998). $`\mathrm{Ly}\alpha `$ studies confirm that similar outflows are present in higher $`z`$ galaxies at redshifts of order $`3`$ (Pettini et al. 1998) and even higher (Frye & Broadhurst 1998; Warren et al. 1998 ), where the overall space density of small starbursting galaxies is larger.
A natural prediction of such a picture would not only be the heating of the ICM by the dwarf population, but the enrichment of this gas by metals expelled by the SNeII powering the galactic winds. In fact, the metallicity of the intracluster medium is observed to be quite high and constant ($`0.3Z_{}`$) over a large range of cluster masses (Renzini 1997). This roughly constant value is evidence that the clusters have neither lost or gained a large amount of baryons during their evolution, as massive cluster outflows or accretion of pristine gas would lead to a large scatter. Furthermore, this metallicity is roughly the same today as observed at $`z0.3`$ and perhaps even at $`z1`$ (Hattori et al. 1997), suggesting that IGM enrichment took place early on in the lifetime of the clusters (Renzini 1999).
Several authors have investigated the possibility that this enrichment is due to the present cluster members and that most of the X-ray emitting hot gas in clusters is due to outflows from galaxies that populate them today. Okazaki et al. (1993) considered gas ejection from bright, elliptical galaxies in clusters, concluding that these galaxies could contribute no more than 10% of the ICM. Trentham (1994) suggested that the ICM instead may be primarily due outflows from precursors to the observed populations of dwarf galaxies in clusters, but Nath and Chiba (1995) concluded that the metallicities generated in this type of scenario were sufficient to explain only clusters with low-metallicity gas.
Perhaps most convincingly Gibson and Matteucci (1996) have studied the possibility that the elliptical galaxies observed in clusters today were responsible for the majority of ICM enrichment. By developing models of galactic winds consistent with the observed properties of cluster ellipticals they concluded that even in their “maximal models” in which all the gas returned by dying stars is ejected into the ICM, neither the giant elliptical, the dwarf spheroidals, or both these populations combined can be responsible for more that 33/38% of the ICM. And yet this gas is hot and metal rich.
A final piece of evidence as to the origin of ICM metals comes from studies of relative abundances. While for many years, studies of metals in clusters were limited to measurements of iron abundances, the X-ray Advanced Satellite for Cosmology and Astrophysics (ASCA) (Tanaka, Inoue,& Holt 1994) provided the first opportunity to study other metals. This allowed Loewenstein & Mushotzky (1996) to observe excesses of alpha-elements such as O, Mg, and Si, which indicate that most of the ICM enrichment was due to SNeII rather than Type Ia (SNeIa). This result was later contested by Ishimaru & Arimoto (1997), who claimed that SNeIa enrichment could be responsible for 50% or more of ICM metals, based on the solar ‘meteoritic’ metallicity rather than the solar coronal gas. Gibson, Loewenstein, & Mushotzky (1997) pointed out, however, that this result was specifically linked to the SNeII yield used in their analysis. While they were unable to rule out Ishimaru & Arimoto’s claim within theoretical uncertainties, they showed that more recent SN models that treat convection and mass-loss, reduce the fractional contribution ICM by SNeIa to less than 5%. Renzini (1997) on the other hand while confirming the importance of early SneII, has argued that later enrichment by significant numbers of SNeIa are required to produce the observed high level of Fe relative to alpha elements.
In summary then, metallicities studies of clusters indicate that the ICM was enriched early on, that this enrichment is too high to be due to purely to the observed cluster members, and that the relative abundances of metals are suggestive of widespread SNeII enrichment; all indicative of enrichment by high-redshift dwarf galaxies.
### 2.2. Galactic Outflows and the Intergalactic Medium
The first galaxies to form would have necessarily had a huge impact not only in regions that would later form clusters, but on the Intergalactic Medium (IGM) as a whole. In standard schemes for hierarchical structure formation, the first baryonic objects formed at redshifts $`40`$ and with masses of order $`10^5M_{}`$, somewhat smaller than the dwarf galaxies observed today (Haiman, Thoul, & Loeb 1996). While several authors have shown that the total flux of UV photons from stars within these objects is sufficient to reionize the universe (Couchman & Rees 1986; Fukugita & Kawakasi 1994; Shapiro, Giroux, & Babul 1994), these objects are likely to suppress their own formation long before this occurs. As molecular hydrogen is easily photo-dissociated by 11.2-13.6 eV photons, to which the universe is otherwise transparent, the emission from the first stars quickly destroys all avenues for cooling by molecular line emission. This quickly raises the minimum virial temperature necessary to cool effectively to approximately 10,000 K, suppressing the further formation of objects with masses $`10^8M_{}`$ (Haiman Rees, & Loeb 1997; Ciardi et al. 1999).
Thus reionization, if achieved by galaxies, was relatively late ($`z15`$) and associated with objects of similar size or larger than dwarf galaxies. This high mass scale, implies that the IGM phase transition must have occurred in a clumpy and inhomogeneous manner, with important implications for small-scale microwave background anisotropies (see eg. Aghanim et al. 1996; Miralda-Escude, Haehnelt, & Rees 2000; Scannapieco 2000).
Early enrichment and heating is similarly expected to be inhomogeneous, following the spatial distribution of the first galaxies which are restricted to rare overdense regions. Hence measurements of a low average metallicity ($`0.01Z_{}`$) of the Ly-forest at high redshift (Songalia & Cowie 1996; Savaglio 1997; Songalia 1997) may be viewed not so much as evidence for a low metallicity intergalactic medium, as an inhomogeneous one. High-redshift regions of higher density that later develop into clusters would therefore represent the most polluted volumes, where one would naturally expect a significantly higher mean metallicity.
Direct evidence that the IGM is enriched by galaxy outflows may be inferred from detailed observations of local dwarf spheroidal galaxies. For example, ASCA observations (della Ceca et al. 1996) show that the star-forming dwarf galaxy NGC 1569 is surrounded by a hot ($`8\times 10^6`$ K) halo of gas whose temperature greatly exceeds the virial temperature of the galaxy and whose line strengths are consistent with $`0.25Z_{}`$. Similarly, X-ray observations of the dwarf irregular NGC 4449 by Bomans, Chu, & Hopp (1997) show that it is embedded in a supergiant shell of $`2\times 10^6`$K gas with a metallicity of $`0.3Z_{}.`$
### 2.3. Galactic Outflows and Galaxy Formation
The presence of an inhomogeneous IGM with strong temperature and metallicity fluctuations creates quite a different environment for galaxy formation than the primordial conditions often assumed. This becomes clear when one contrasts the virialization and cooling of primordial gas to that of gas pre-enriched and heated by dwarf outflows. In Figure 1 we replot a classic comparison first made in Dekel and Silk (1986). Collapse redshifts of $`1\sigma `$ to $`2\sigma `$ dark matter halos (dotted lines) are calculated from linear theory for both a flat CDM ($`\mathrm{\Omega }_0=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, $`\mathrm{\Omega }_b=0.07`$, $`\sigma _8=0.6`$, and $`\mathrm{\Gamma }=0.44`$) and a $`\mathrm{\Lambda }`$CDM model ($`\mathrm{\Omega }_0=0.35`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.65`$, $`\mathrm{\Omega }_b=0.06`$, $`\sigma _8=0.87`$ and $`\mathrm{\Gamma }=0.18`$), assuming spherical collapse. The solid lines show the lowest redshift that a sphere of gas can virialize and still have time to cool and form a galaxy by $`z=0`$, as calculated by the cooling models described in §3.2. In each pair the upper line corresponds to primordial gas, and the lower line to gas that has been enriched to $`0.1Z_{}`$. Here we see that even modest enrichment by outflows has the potential to greatly accelerate galaxy formation on the $`10^{12}M_{}/h`$ scale due to the long cooling times of large clouds without the additional avenue for cooling afforded by line emission by metals.
Also on this plot, we show the mass of halos associated with a virial temperature of $`5\times 10^5`$K, typical of galaxy outflows (dashed line). This serves as an estimate of the smallest mass of galaxies that can form in areas impacted by outflows without being disrupted by shocks. This suggests that while high-mass galaxy formation may be enhanced by dwarf outflows, the formation of dwarf galaxies has the potential to be suppressed.
The features shown on this plot invite comparison with the properties observed in elliptical galaxies, which are preferentially found in the most enriched regions in the universe, and whose mass functions are biased to large values relative to the rest of the galaxy population. It is also interesting to note that various analytical and N-body studies of cold dark matter models (Kauffmann, White, & Guiderdoni 1993; Klypin et al. 1999; Moore et al. 1999) have shown that $`50`$ satellites with circular velocities $`20`$ km/s should be found within 600 kpc of the Galaxy, while only 11 are observed. Again, this lack of local dwarf galaxies is suggestive of galaxy formation in an intergalactic medium that has been impacted by outflows.
Closer to home, the abundance distribution of long-lived stars in the solar neighborhood shows far too few low-metallicity stars compared with a simple “closed box” model of galactic chemical evolution. This is the long standing G-dwarf problem, first pointed out by van de Bergh (1962) and Schmidt (1963). The sudden drop in the number of G-dwarfs with metallicities below $`0.1Z_{}`$ has lead to a number of models in which an initial production spike or “prompt initial enrichment” adds metals to the gas out of which the majority of stars form (e.g., Truran & Cameron 1971; Ostriker & Thuan 1975; Köppen & Arimoto 1990). This ad hoc floor to the metallicity then allows good fits to the Galactic stellar data.
Similar pre-enrichment may also be necessary to explain the metallicity distribution of stars outside our own galaxy. Integrated spectra of elliptical galaxies have been studied with care (Worthey, Dorman, & Jones 1996) concluding rather puzzlingly that such galaxies do not have a closed-box history but again show the need for a “floor” of around $`0.1Z_{}`$ (Thomas, Greggio, & Bender 1999).
## 3. Cosmological Monte Carlo
The great preponderance of observational clues and pointers suggests that we should seriously consider a scenario in which the formation of modern-day galaxies was preceded by an era of IGM enrichment by a high-redshift population of dwarf starbursting galaxies. Note that from a theoretical point a view, such a model is not ad hoc, but rather a consequence of hierarchal structure formation. As smaller objects form early and are able to suppresses further small-scale formation, it is only natural that larger galaxies would have formed in a second wave of collapse within this pre-enriched medium.
The density dependence of the effect of outflows on galaxy formation requires a 3-D calculation. To date only the average properties of the IGM as a function of time have been explored in one dimensional calculations by Blanchard, Valls-Gabaud, & Mamon (1992) and by Nath and Trentham (1997). Such averages are appropriate to the IGM as a whole, but underestimate the effect on galaxy formation. Thus for example, the outflow model proposed by Nath and Trentham (1997) produces metallicities $`.01Z_{}`$, while the observed metallicities of cluster gas and elliptical galaxies are comparable to solar abundance (Renzini 1993).
In order to determine the impact of outflows on galaxy formation, we have developed a simple cosmological Monte Carlo code, in which we realize a small volume of the universe and study the linear evolution of objects in the mass range $`2\times 10^8\mathrm{\Omega }_0/h\mathrm{M}_{}M1\times 10^{13}\mathrm{\Omega }_0/h\mathrm{M}_{}`$ within it, with reasonable extensions to deal with collapse. Our philosophy is not to pretend that we can possibly reproduce in detail the complicated processes of dark matter halo collapse, galaxy formation, and gas infall and expulsion, but rather to explore a model that captures the essential features and predicts with some confidence rough magnitudes and trends that can be compared with observations.
### 3.1. Collapse of Dark Matter Halos
We choose for our simulation a cubic comoving volume of $`(12\mathrm{M}\mathrm{p}\mathrm{c}/h^1)^3`$ divided into $`512^3`$ cells and with periodic boundary conditions. On this mesh we construct a cosmological linear overdensity field $`\delta (𝐱,z)=\rho (𝐱,z)/\overline{\rho }(z)`$ where $`\overline{\rho }(z)`$ is global average density of the universe, $`𝐱`$ is a position in comoving coordinates, and $`z`$ is the redshift. Transforming into Fourier space the density field can be reexpressed as $`\stackrel{~}{\delta }(𝐤,z)d^3𝐱\mathrm{exp}(i𝐤𝐱)\delta (𝐱,z)`$. From linear theory, the evolution of the density field as a function of time is given simply by $`\stackrel{~}{\delta }(𝐤,z)=\stackrel{~}{\delta }_0(𝐤)D(z)/D_0`$ where $`D(z)`$ is the dimensionless growth factor, $`\stackrel{~}{\delta }_0(𝐤)\stackrel{~}{\delta }(𝐤,0)`$, and $`D_0D(0).`$
In the CDM model, the Gaussian random field $`\stackrel{~}{\delta }(𝐤)`$ can by constructed by randomly choosing modes such that the variance is given by $`\stackrel{~}{\delta }(𝐤)\stackrel{~}{\delta }(𝐤^{})=\stackrel{~}{\delta }(𝐤)\stackrel{~}{\delta }^{}(𝐤^{})=(2\pi )^3\stackrel{~}{\delta }(𝐤+𝐤^{})P(𝐤)`$. The power spectrum is then
$$P(k)=2\pi ^2\delta _H^2k^nT^2(k_p\mathrm{Mpc}/h\mathrm{\Gamma }),$$
(1)
where $`T(q)`$ is the CDM transfer function given by Eq. (G3) of Bardeen et al. (1996), and we choose $`a_0H_0=c`$, such that $`k_p=k/a_0=kH_0/c`$ is the physical wave number. The shape parameter $`\mathrm{\Gamma }`$ is defined as $`\mathrm{\Gamma }\mathrm{\Omega }_0h\mathrm{exp}(\mathrm{\Omega }_b\mathrm{\Omega }_b/\mathrm{\Omega }_0)`$ (Sugiyama 1995) and is observed be to $`\mathrm{\Gamma }=0.23_{0.034}^{+0.042}`$ given a spectral index of $`n=1`$ (Vianna & Liddle 1996). We fix the normalization factor, $`\delta _H`$, against the number abundance of clusters. The variance of the mass inclosed in a sphere of radius $`R`$ is
$$\sigma ^2(R)=\frac{1}{2\pi ^2}_0^{\mathrm{}}k^2𝑑kP(k)W^2(kR).$$
(2)
If we choose a spherical top-hat window function defined by $`W(x)3\left[\frac{\mathrm{sin}(x)}{x^3}\frac{\mathrm{cos}(x)}{x^2}\right]`$ then we can normalize our fluctuation spectrum by setting $`\sigma _8\sigma (8\mathrm{Mpc}h^1)=(0.6\pm 0.1)\mathrm{\Omega }_0^{C(\mathrm{\Omega }_0)},`$ where $`C(\mathrm{\Omega }_0)=.36+0.31\mathrm{\Omega }_00.23\mathrm{\Omega }_0^2`$ in the open case and $`C(\mathrm{\Omega }_0)=.59+0.16\mathrm{\Omega }_00.06\mathrm{\Omega }_0^2`$ if we take ($`\mathrm{\Omega }_0+\mathrm{\Omega }_\mathrm{\Lambda }=1`$) (Vianna & Liddle 1996).
Having constructed $`\stackrel{~}{\delta }_0(𝐤)`$ in this manner, we then convolve it with window functions with lengths scales corresponding to ten masses arranged logarithmically from $`M_12\times 10^8\mathrm{\Omega }_0/hM_{}`$ to $`M_{10}10^{13}\mathrm{\Omega }_0/hM_{}`$, to obtain $`\delta _0^{M_1}(𝐱)\delta _0(𝐱,R_{2\times 10^8\mathrm{\Omega }_0/hM_{}})=\frac{d^3𝐤}{(2\pi )^3}\mathrm{exp}(i𝐤𝐱)\stackrel{~}{\delta }(𝐤)W(R_{2\times 10^8\mathrm{\Omega }_0/hM_{}}k)`$, $`\delta _0^{M_2}(𝐱)\delta _0(𝐱,R_{6.6\times 10^8\mathrm{\Omega }_0/hM_{}})`$, … $`\delta _0^{M_{10}}(𝐱)`$, the linear density field at $`z=0`$ smoothed at each mass scale. Here the minimum mass of our simulations is set by molecular cooling constraints discussed further below.
From the spherical collapse model we can identify the linear overdensity $`\delta _{\mathrm{sc}}(z)`$ at which the true density field has virialized. For flat models this value is $`1.69`$ at all times while in open cases it is a weakly decreasing function of redshift as fitted in Appendix A of Kitayama & Suto (1996). For each peak in the linear overdensity field such that $`\delta _0^{M_i}(𝐱,R)>\delta _{\mathrm{sc}}(z=0)`$, we identify a collapse redshift $`z_{\mathrm{sc}}`$ such that
$$\delta _0^{M_i}(𝐱)=\delta _c(z_{\mathrm{sc}})\frac{D_0}{D(z_{\mathrm{sc}})}.$$
(3)
This approach has the limitation that it becomes inaccurate as $`\sigma (R)D(z)/D_0`$ approaches $`\delta _{\mathrm{sc}}`$. Thus while the spherical collapse model works well at determining the number and spatial clustering of rare objects (Lacy & Cole 1993; Mo, Jing, & White 1996, 1997) it fails to agree with numerical simulations for more common low-mass halos (Lacey & Cole 1994; Sheth & Tormen 1999). Furthermore, the object-by-object identification of linear peaks with collapsed objects has been shown to be unreliable, even in cases in which the statistical properties are in good agreement with N-body simulations (Bond et al. 1991).
Sheth, Mo, & Tormen (1999) have shown that these predictions can be improved by accounting for the ellipticity of collapsing clouds caused by tidal forces from nearby collapsing peaks. By including a correction factor for the critical linear overdensity, they have been able not only to improve statistical predictions over the spherical model, but do well on an object-by-object basis. As our Monte-Carlo approach depends on this object-by-object identification, we can greatly improve the accuracy of our simulations by comparing each peak with $`\delta _{\mathrm{ec}}(z)`$, such that
$`\delta _0^{M_i}=`$ $`\delta _{\mathrm{ec}}(z_{\mathrm{ec}}){\displaystyle \frac{D_0}{D(z_{\mathrm{ec}})}}`$
$`=`$ $`\delta _{\mathrm{sc}}(z_{\mathrm{ec}}){\displaystyle \frac{D_0}{D(z_{\mathrm{ec}})}}\left[1+\beta \left({\displaystyle \frac{\sigma ^2(R)}{\delta _{\mathrm{sc}}^2(z_{\mathrm{ec}})}}{\displaystyle \frac{D(z_{\mathrm{ec}})^2}{D_0^2}}\right)^\gamma \right],`$ (4)
where $`\beta =0.47`$ and $`\gamma =0.615`$.
Having identified the collapsed peaks by either algorithm, we then arrange them in order of decreasing collapse redshift, $`z_\mathrm{c}=z_{\mathrm{sc}}`$ or $`z_{\mathrm{ec}}`$, to obtain a list of candidate points and redshifts for collapsed halos at each of the different mass scales: $`𝒳_{\mathrm{candidate}}^{M_i}\{𝐱_1^{M_i},𝐱_2^{M_i},\mathrm{}\}\mathrm{and}𝒵_{\mathrm{c},\mathrm{candidate}}^{M_i}\{z_{\mathrm{c},1}^{M_i},z_{\mathrm{c},2}^{M_i},\mathrm{}\}.`$ From this list we exclude unphysical points corresponding to all halos that collapse within a previously collapsed halo at the same or greater mass scale. That is we remove all points such that
$$M_iM_j,𝐱_k^{M_i}𝐱_l^{M_j}<R_{M_j},\mathrm{and}z_{\mathrm{c},k}^{M_i}<z_{\mathrm{c},l}^{M_j},$$
(5)
where the distance is calculated accounting for the periodicity of the simulation volume.
With these points excluded, we then have the full history of the collapsing dark matter halos at each of the various mass scales: $`𝒳^{M_i}\{𝐱_1^{M_i},𝐱_2^{M_i},\mathrm{}\}`$, and $`𝒵_\mathrm{c}^{M_i}\{z_{\mathrm{c},1}^{M_i},z_{\mathrm{c},2}^{M_i},\mathrm{}\}`$. The merger history of these objects is also simply calculated by searching for all points such that
$$M_i<M_j,𝐱_k^{M_i}𝐱_l^{M_j}<R_{M_j},\mathrm{and}z_{\mathrm{c},k}^{M_i}>z_{\mathrm{c},l}^{M_j}.$$
(6)
Each such case indicates that the halo centered at $`𝐱_k^{M_i}`$ was absorbed into the larger halo centered at $`𝐱_l^{M_j}`$ at a redshift of $`z_{\mathrm{c},l}^{M_j}.`$
### 3.2. Comparison with Analytical Results
At each redshift, the number density of collapsed halos in the simulation can be compared to analytical predictions. For the spherical collapse model this gives (Press & Schechter 1974)
$`{\displaystyle \frac{dn_{\mathrm{sc}}(M,z)}{dM}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\rho (z)}{M}}{\displaystyle \frac{\delta _{\mathrm{sc}}(z)D_0}{\sigma (M)^2D(z)}}`$ (7)
$`\mathrm{exp}\left({\displaystyle \frac{\delta _{\mathrm{sc}}^2(z)D_0^2}{2\sigma ^2(M)D(z)^2}}\right){\displaystyle \frac{d\sigma (M)}{dM}}.`$
Strictly speaking, Eq. (7) corresponds to the number of peaks as determined by a sharp $`k`$-space filter, rather than the number determined by the top-hot filter used in our code. Equating these quantities has become somewhat of a common practice however, and has been shown to be a good estimate.
In Figure 2 we compare the number density of objects in our simulation as function of redshift with the number predicted by Eq. (7). As both the Press-Schechter integral, and our peak finding algorithm become inaccurate at late times, we exclude all peaks that collapse with a redshift lower that a minimum value, $`z_{\mathrm{min}}(M_i)`$ such that $`\sigma (M_i)D(z_{\mathrm{min}}(M_i))/D_0=1.3\delta _{\mathrm{sc}}.`$ We consider the $`\mathrm{\Lambda }`$CDM model shown in Figure 1, ($`\mathrm{\Omega }_0=0.35`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.65`$, $`\mathrm{\Omega }_b=0.06`$, $`\sigma _8=0.87`$ and $`\mathrm{\Gamma }=0.18`$, $`h=0.65`$).
Here we see that the number of identified objects agrees well at most mass scales and redshifts. At the smaller scales, there is some discrepancy in the collapse time of objects, with the code slightly over-predicting the number of objects formed at very early times. At larger-range scales, the Monte-Carlo approach does the best, matching the Press-Schechter predictions at all redshifts. Finally, at the highest mass scales, the low number of objects introduces significant statistical noise in our comparison.
The spatial distribution of collapsed halos can also be compared to analytical predictions. If we define the correlation function at a smoothing scale $`R`$ as the Fourier transform of the power spectrum smoothed by an appropriate top-hat window function, $`\xi _R(r)\frac{1}{2\pi ^2}_0^{\mathrm{}}k^2𝑑kP(k)W^2(kR)\frac{\mathrm{sin}(kr)}{kr}`$, then the correlation of collapsed peaks of scale $`R`$ can be approximated by (Kaiser 1984)
$`1+\xi _{R,\nu }(r)=`$ $`\left({\displaystyle \frac{2}{\pi }}\right)^{1/2}\left[\mathrm{erfc}(\nu /2^{1/2})\right]^2\times `$ (8)
$`{\displaystyle _\nu ^{\mathrm{}}}𝑑ye^{y^2/2}\mathrm{erfc}\left[{\displaystyle \frac{\nu y\xi _R(r)/\xi _R(0)}{\sqrt{22\xi _R^2(r)/\xi _R^2(0)}}}\right],`$
where erfc is the complimentary error function and $`\nu \delta _cD_0/\sigma (R)D(z).`$ Note that this expression arises from considering the fraction of Gaussian-distributed objects above a certain threshold, rather than from the more sophisticated excursion set approach (Bond, Efstathiou, & Kaiser 1991; Lacy & Cole 1993) used to derive Eq. (7).
In Figure 3 we compare the correlation function as given in Eq. (8) to that derived from the distances between the collapsed peaks at the five lowest mass-scales in our simulation, with redshifts chosen such that $`\nu =1.5`$ at each length scale. Eq. (8) is intended to represent the correlations between collapsed objects at or above a particular mass scale corresponding to the smoothing length scale. In practice, however, it makes little difference whether we compute $`\xi _{R_{M_i}}(r)`$ considering distances between all objects with masses above $`M_i`$ or only between those in the mass bin corresponding to $`M_i`$ itself. As comparisons only within a particular bin are both easier to calculate and understand physically, however, it is this quantity that we plot in Figure 3.
Here we see that the numerical and analytical expressions agree at lengths scales $`R^{M_i}`$, but diverge at smaller distances. As the simulation simply excludes all points closer than $`R_{M_i}`$ there is a sharp fall of in $`\xi _{R_{\mathrm{M}_\mathrm{i}}}(r)`$ below this distance. The analytical expression, however, does not exclude pairs of objects that are contained within a larger virialized object (as would be excluded in an excursion set calculation) and increases dramatically at small $`r`$. Thus these quantities are derived in a manner such that they must be discrepant at small distances, and we can have some confidence that the overall distribution of objects in our simulations is reasonable.
While the spherical collapse model is a logical testing ground for our code and can be simply compared to analytical predictions, a more accurate simulation results by applying the elliptical collapse criteria for halo formation, Eq. (4). This allows not only for a more accurate identification of collapsed peaks on an object-by-object basis, but delays the collapse of the more common peaks, and thus the minimum accurate redshift at each mass scale is decreased to $`z_{\mathrm{min}}(M_i)`$ such that $`\sigma (M_i)D(z_{\mathrm{min}}(M_i))/D_0=1.3\delta _{\mathrm{ec}}`$.
In this model, the Press Schechter relation, Eq. (7), is modified to
$`{\displaystyle \frac{dn_{\mathrm{ec}}(M,z)}{dM}}=`$ $`A\left[1+\left({\displaystyle \frac{\sigma (M)D(z)}{\delta _{\mathrm{sc}}D_0}}\right)^{2q}\right]\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\rho (z)}{M}}`$ (9)
$`{\displaystyle \frac{\delta _{\mathrm{sc}}(z)}{\sigma (M)^2D(z)}}\mathrm{exp}\left({\displaystyle \frac{\delta _{\mathrm{sc}}^2(z)D_0^2}{2\sigma ^2(M)D(z)^2}}\right){\displaystyle \frac{d\sigma (M)}{dM}},`$
where $`A=0.322`$ and $`q=0.3.`$ Note that this function is essentially equivalent to the spherical collapse predictions for rare, high-mass peaks, and thus the $`\sigma _8`$ normalization from the number density of galaxy clusters is unchanged in this model. This expression has been checked against the number density of objects in our simulation and again is in good agreement.
### 3.3. Initial Gas Infall
While our simple nonlinear collapse model captures the overall history of the dark matter distribution, the evolution of the gas is more involved. While gas traces dark matter at early times, the formation of bound objects and subsequent star formation requires that the gas be able to cool to $`T0`$.
If gas collapses and virializes along with a dark matter perturbation, and we assume an isothermal distribution, then it will be heated to a temperature of (Eke, Cole, & Frenk 1996)
$$T_{\mathrm{vir}}=\frac{90\mathrm{K}}{\beta }\left(\frac{6.8}{5X+3}\right)M_6^{2/3}(1+z_c)\left(\frac{\mathrm{\Omega }_0\mathrm{\Delta }_c(z_c)}{\mathrm{\Omega }(z_c)18\pi ^2}\right)^{1/3},$$
(10)
where $`\beta `$ is the ratio of specific galaxy kinetic energy to specific gas thermal energy, $`X`$ is the hydrogen mass fraction which we take to be $`0.76`$, $`M_6M/(10^6M_{}/h)`$, and $`\mathrm{\Delta }_c`$ is the ratio of the mean halo density to the critical density at the redshift of collapse, a constant ($`18\pi ^2`$) for the $`\mathrm{\Omega }=1`$ case and a otherwise a weak function of $`z`$ fitted in Kitayama & Suto (1996). Navarro, Frenk, & White (1995) have shown that this relation with $`\beta =1.07\pm 0.05`$ is an accurate approximation to the results of N-body/hydrodynamic simulations in the case of an $`\mathrm{\Omega }=1`$ universe. Here we follow Eke, Cole & Frenk (1996) and fix $`\beta =1`$ for all cosmological models.
Once the gas has fallen into a potential well, we adopt a simple model to calculate the time scale at which the gas will cool and form stars. Neglecting the gravitational energy associated with further gas infall we have
$$n_e^2(r)\mathrm{\Lambda }(T(r),Z(r))=\frac{3}{2}n(r)k_B\frac{dT(r)}{dt},$$
(11)
where $`n_e(r)`$ is the number density of electrons at a radius $`r`$, $`n_{\mathrm{tot}}(r)`$ is the total number density, $`T(r)`$ is the temperature of the gas, $`k_B`$ is the Boltzmann constant, and $`\mathrm{\Lambda }`$ is the radiative cooling function, which is strongly dependent on the metallicity $`Z(r)`$ of the gas. Here we take $`\mathrm{\Lambda }`$ as tabulated by Sutherland and Dopita (1993) for equilibrium configurations, although strictly speaking these estimates are not exact due to the difference in abundance ratios for alpha elements in SNeII outflows as compared to solar proportions. Considering only helium and hydrogen the mean molecular weight is $`\mu m_p=\rho _g/n=m_p4/(85Y)`$ and the ratio of the electron and total number density are related by $`\eta \frac{n_e}{n_{\mathrm{tot}}}=\frac{4(Y^11)+2}{8(Y^11)+3}`$, where $`Y`$ is the helium fraction by mass, here taken to be $`0.25.`$
For a general configuration of gas, Eq. (11) is difficult to solve, and we therefore adopt a simple model for each collapsing halo as an isothermal sphere ($`\rho (r)r^2`$) with constant metallicity. We then follow the simple heuristic model of White & Frenk (1991) in which all the gas within some “cooling radius” $`r_{\mathrm{cool}}`$ cools instantaneously and all the gas outside this radius stays at the virial temperature of the halo, with $`r_{\mathrm{cool}}`$ moving outward with time. For an isothermal sphere this gives
$`{\displaystyle \frac{dM_{\mathrm{cool}}}{dt}}=`$ $`4\pi \rho _g(r_{\mathrm{cool}})r_{\mathrm{cool}}^2{\displaystyle \frac{dr_{\mathrm{cool}}}{dt}}`$
$`=`$ $`12f_{\mathrm{hot}}^{3/2}\left({\displaystyle \frac{T_{\mathrm{vir}}}{\mathrm{K}}}\right)\left({\displaystyle \frac{\mathrm{\Lambda }(T_{\mathrm{vir}},Z)}{10^{23}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^3}}{\displaystyle \frac{\mathrm{yr}}{t}}\right)^{1/2}M_{}\mathrm{yr}^1,`$ (12)
where $`f_{\mathrm{hot}}`$ is the fraction of the halo mass in the form of hot gas. While derived from a simple model, this expression is in good agreement with similarity solutions for cooling flows given by Bertschinger (1989) (which are $`28\%`$ smaller than this expression), and the one-dimensional simulations of Forcada-Miró & White (1996) (which are $`15\%`$ smaller). We therefore follow Somerville (1997) in multiplying the right hand side of this equation by an overall factor $`f_0`$ which we take to be $`0.8`$. Finally, we approximate $`f_{\mathrm{hot}}M_{\mathrm{baryon}}/M_{\mathrm{halo}}=\mathrm{\Omega }_b/\mathrm{\Omega }_0`$ at all times, and take the time at which $`M_{\mathrm{cool}}=M_{\mathrm{baryon}}`$ to be time at which a galaxy is formed.
We show this formation time as a function of mass for four different metallicities at $`z=1`$ in Figure 4. From this diagram we see that while cooling times are negligible for dwarf galaxies, cooling times for primordial-abundance halos with masses above $`10^{12}M_{}`$ are sufficiently long to prevent galaxy formation at low redshifts. This is dependent on metallicity, as we saw in Figure 1. In our simulation we therefore neglect cooling times for halos with masses $`10^{10}M_{}`$, while using the appropriate pre-enriched cooling times for more massive objects.
Note also the sharp rise in cooling times at temperatures below $`10,000`$ K, due to the lack of atomic transitions with energies below a few eV. Cooling at lower temperatures can only occur through molecular transitions, and thus the collapse of $`T_{\mathrm{vir}}10^4`$ primordial gas clouds depends strongly on the existence of molecular hydrogen. A detailed investigation of the production and dissociation of $`\mathrm{H}_2`$ in the early universe has been conducted by Haiman, Rees, & Loeb (1997). As $`\mathrm{H}_2`$ is easily photo-dissociated by 11.2-13.6 eV photons, to which the universe is otherwise transparent even before reionization, they find that a UV flux of $`10^{22}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1\mathrm{Hz}^1\mathrm{sr}^1`$, is capable of dissociating all $`\mathrm{H}_2`$ in collapsing halos (see also Ciardi et al. 1999). As this is more than two orders of magnitudes smaller that the reionizing flux, we assume in our simulations that a small population of early stars quickly depleted the primordial gas of molecular hydrogen. Thus we exclude the formation of all galaxies with virial temperatures below $`10,000`$ K, which sets the lower halo mass value to $`5\times 10^7\mathrm{M}_{}h^1`$, justifying the choice of spatial resolution for our simulations.
### 3.4. Heating and Enrichment of the IGM
Once a dwarf galaxy has been formed, we assume that a galactic outflow is immediately generated due to SNeII from the short-lived massive stellar population. In order to get a rough estimate of the mass of metals injected into the ICM, the radius over which they are mixed into the surrounding gas, and the heating of the ICM, we model them according to the simple analytical model developed in Tegmark, Silk, & Evrard (1993, hereafter TSE). Here the authors treat the outflow as an expanding shell that sweeps up most of the baryonic IGM and loses only a small fraction $`f_m=0.1`$ to the interior, such that the mass of the shell is
$$m(t)\frac{4\pi }{3}R(t)^3\rho _b(1f_m)$$
(13)
where $`R(t)`$ is the radius of the outflow. Acceleration of the shell is due to the internal pressure and deceleration from gravitational breaking, both estimated in the thin shell approximation (Ostriker & McKee 1988). With these approximations $`\ddot{R}`$ becomes
$$\ddot{R}=\frac{8\pi pG}{\mathrm{\Omega }_bH^2R}\frac{3}{R}(\dot{R}HR)^2\mathrm{\Omega }_0\frac{H^2R}{2}\frac{GM}{R^2}.$$
(14)
Note that the gravitational term here includes a correction for the presence of a halo of mass $`M`$ out of which the outflow expands. The pressure of the hot gas within the shell is given by $`p=\frac{2E_T}{3V}`$ where $`E_T`$ is the thermal energy of the interior, and energy conservation gives $`\frac{dE_T}{dt}=Lp\frac{dV}{dt}`$, where $`L`$ is the luminosity, incorporating all heating and cooling of the interior plasma.
The evolution of the blast wave is completely determined by the assumed luminosity history of the blast wave. Following TSE we consider five contributions to this luminosity: the cooling by Compton drag against the microwave background ($`L_{\mathrm{comp}}`$), the cooling due to bremsstrahlung and other two body interactions ($`L_{\mathrm{ne}^2}`$) as in Eq. (11), energy injection from supernovae ($`L_{\mathrm{sn}}`$), cooling by ionization of the IGM ($`L_{\mathrm{ion}}`$), and heating from collisions between the shell and the IGM ($`L_{\mathrm{diss}}`$); such that,
$$L=L_{\mathrm{comp}}L_{\mathrm{ne}^2}+L_{\mathrm{sn}}L_{\mathrm{ion}}+L_{\mathrm{diss}}.$$
(15)
The first of these we fix according to Eq. (4) of TSE, while $`L_{\mathrm{ne}^2}`$ is negligible in this regime. The other parameters require some calibration to the observed properties of outflows from dwarf ellipticals and our overall model of the IGM at the time of dwarf formation. Let us consider these individually.
Following TSE we assume that the energy in the blast wave is proportional to the number of SNeII and winds for O and B stars and that this in turn is proportional to the baryonic mass which has collapsed. Energy injection occurs over a relatively short time during which the dwarf galaxy is undergoing a starburst phase. We assume that a fraction of $`ϵ_{\mathrm{SF}}`$ of the baryons in each galaxy collapse to form stars, that one supernova is formed for every 100 $`M_{}`$ of stars in the halo, and that each supernova provides a typical energy output of $`10^{51}`$ ergs, with an equal contribution coming from winds from massive stars. If a fraction $`ϵ_{\mathrm{wind}}`$ of this input goes into the galactic outflow then the total energy from the supernovae is $`E=ϵ_{\mathrm{wind}}ϵ_{\mathrm{sf}}2.0\times 10^{55}\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}M_6.`$ Assuming that this energy is released at a constant rate during a period of $`t_{\mathrm{burn}}=1.0\times 10^8`$ years, we find
$$L_{\mathrm{sn}}=ϵ_{\mathrm{wind}}ϵ_{\mathrm{sf}}1.6\times 10^6L_{}/hM_6\times \frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}$$
(16)
where the factor of $`h`$ is a result of the units we have chosen for $`M_6`$. Dekel and Silk (1986) have proposed that the observed differences in surface brightness and metallicity between the observed classes of diffuse dwarfs and normal galaxies can be understood in terms of outflows occurring in halos with virial velocities below critical value on of the order of 100 km/s, which corresponds to approximately $`2.5\times 10^9ϵ_{\mathrm{SF}}(1+z)^{3/2}\mathrm{\Omega }_b/\mathrm{\Omega }_0`$ supernovae. Here we take a conservative approach and exclude all bursts greater than $`2\times 10^7ϵ_{\mathrm{SF}}\mathrm{\Omega }_b/\mathrm{\Omega }_0,`$ which is typically a few hundred thousand supernovae and corresponds to a maximum luminosity of
$$L_{\mathrm{sn},\mathrm{max}}=ϵ_{\mathrm{wind}}ϵ_{\mathrm{sf}}3.2\times 10^9L_{}\times \frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}$$
(17)
achieved by objects above $`2\times 10^9M_{}.`$
Another potential source of drag on the expanding shell is energy losses to ionizing the IGM, given by $`L_{\mathrm{ion}}=f_{\mathrm{neutral}}f_\mathrm{m}n_bE_o\mathrm{\hspace{0.17em}4}\pi R^2[\dot{R}HR],`$ where $`n_b`$ is the number density of baryons, $`E_013.6`$ eV, and $`f_{\mathrm{ion}}`$ is the ionization fraction of the IGM into which the shell is expanding. As they were trying to develop a model for IGM reionization by winds, TSE assume $`f_{\mathrm{neutral}}=1`$ at all times. Here we choose for our fiducial model the case in which blast waves expand into a medium that has completely been pre-ionized ($`f_{\mathrm{neutral}}=0`$) possibly by ionizing photons from active galactic nuclei, or from ionization fronts from the massive stars that precede the winds.
Finally $`L_{\mathrm{diss}}`$ accounts for the heating of the interior plasma due to collisions between the shell and the IGM. Taking this to be some fixed fraction $`f_{\mathrm{diss}}`$ of the kinetic energy lost by the shell gives
$$L_{\mathrm{diss}}=f_{\mathrm{diss}}\frac{3m}{2R}(\dot{R}HR)^3.$$
(18)
In TSE the authors consider the extreme cases in which $`f_{\mathrm{diss}}=0`$ and $`f_{\mathrm{diss}}=1`$, showing that their results are relatively insensitive to this parameter. Here we simply fix $`f_{\mathrm{diss}}=0.5`$ throughout.
With these approximations we can use the solution given in TSE to calculate the full expansion history of the galactic outflow, with their $`L_{}`$ and $`E_{}`$ appropriately modified to account for the gravitational potential of the collapsed halo as given by our Eq. (16). We step forward in time (backwards in redshift) in increments of $`\mathrm{\Delta }z=0.01`$ starting at a maximum redshift of $`z_{\mathrm{max}}=25`$, using a fourth-order Runge-Kutta technique. Figure 5 shows the first galaxy to undergo an outflow, a $`1.1\times 10^8M_{}`$ dwarf expanding into an ionized medium at $`z=25`$ in one such simulation.
As our simulations make no attempt to include ionization and do not account for the presence of ionizing radiation from stars within the dwarf galaxy or an external source of UV radiation, we expect our solution to only be accurate while the temperature of the material is greater than 10, 000 K ($`k_bT=13.6`$ eV). After this point, we assume as a reasonable approximation that the halo simply expands with the Hubble flow and remains at a fixed temperature of 10,000 K.
In this figure, we also include for comparison the results of one such blast wave expanding into a neutral medium. This closely traces the ionized solution up until the temperature of the outflow drops to $`10,000`$ K.
Finally, we model the metallicity of the expanding halo by assuming that each supernova typically ejects 2$`M_{}`$ of metals into the intergalactic medium. This number is consistent with the average stellar yields in SNeII simulations as compiled in Nagataki & Sato (1997), although there are significant theoretical uncertainties between the various simulations. The total mass of metal injected into the interstellar medium of the star forming galaxies is then
$$M_{6,z}=0.02ϵ_{\mathrm{sf}}(\mathrm{\Omega }_b/\mathrm{\Omega }_0)M_6.$$
(19)
Our model has two important free parameters, $`ϵ_{\mathrm{sf}}`$ and $`ϵ_{\mathrm{wind}}`$. While the amount of energy in the winds is directly dependent on the product of these two parameters, their ratio is somewhat more unconstrained. Our approach here then will be to fix the star formation efficiency at a somewhat standard value of $`ϵ_{\mathrm{sf}}=0.1`$ and allow the fraction of the SN energy channeled into the wind to vary between 5 and 20 percent to quantify model uncertainties.
With this simple model we can construct from $`𝒳^{M_i}`$ and $`𝒵_\mathrm{c}^{M_i}`$ the full histories of the blast waves expanding into the intergalactic medium: the center of each expanding bubble ($`𝒳^{\mathrm{blast}}\{𝐱_1^{\mathrm{blast}},𝐱_2^{\mathrm{blast}},\mathrm{}\}`$), its comoving radius as a function of redshift ($`^{\mathrm{blast}}(z)\{R_1^{\mathrm{blast}}(z),R_2^{\mathrm{blast}}(z),\mathrm{}\}`$), temperature as a function of redshift ($`𝒯^{\mathrm{blast}}(z)\{T_1^{\mathrm{blast}}(z),T_2^{\mathrm{blast}}(z),\mathrm{}\}`$), and mass of ejected metals ($`_Z^{\mathrm{blast}}\{M_{Z,1}^{\mathrm{blast}},M_{Z,2}^{\mathrm{blast}},\mathrm{}\})`$. While $`^{\mathrm{blast}}(z)`$ and $`𝒯^{\mathrm{blast}}(z)`$ are updated in redshifts intervals of $`\mathrm{\Delta }z=0.01`$, the history is only stored in intervals of $`\mathrm{\Delta }z=0.05`$ to save memory. Note that the jagged appearance at late times of the comoving size of the outflow in Figure 5 is caused by this discrepancy in times steps, rather than by a numerical instability. Note also that as outflows only occur in objects with masses $`10^{10}`$ for which cooling times are negligible, we can compute the entire outflow history of our Monte Carlo volume without considering the impact of metal enrichment on the gas cooling times. Suppression of galaxy formation due to outflows from nearby objects, however, must be considered for outflow-scale second generation objects.
### 3.5. Formation of ‘Second Generation’ Objects
As we saw in §2, galactic outflows affect subsequent galaxy formation both by shocking the IGM as well as enriching it with metals. Shocks can suppress the formation of nearby galaxies by two main mechanisms, heating the halo gas or striping it from the dark matter. In the first case the gas associated with a collapsing halo is heated to above the viral temperature of the halo but remains bound to the collapsing dark matter. The thermal pressure of the gas then overcomes the dark matter potential and the gas expands out of the halo, preventing galaxy formation. In the second case suppression is caused by stripping of the baryons by a shock from a nearby source. The momentum of the moving shock is sufficient to carry with it the gas associated with the halo of an unvirialized perturbation, thus emptying the dark matter halo of its gas and preventing a galaxy from forming.
In a companion letter (Scannapieco, Ferrara, & Broadhurst 2000), we evaluate both these scenarios of halo suppression and find the stripping mechanism to be dominant, because most shock-heated clouds can radiatively cool within a sound crossing time. As in the case of outflow generation, momentum transfer between the outflows and the collapsing halos is a complicated process, dependent on the overall density profile of the collapsing halo and the possibility of smaller collapsed objects within it. Again, our approach in this exploratory study will be to adopt simple criteria to estimate on average when the momentum in a shell is sufficient to remove material from a collapsing halo.
We allow for baryonic stripping to occur only if a halo has not yet virialized, and assume that stripping occurs in all cases in which a shock moves through the center of the pre-virialized halo with sufficient momentum to accelerate the gas to the escape velocity. Thus we exclude all objects for which
$$fM_sv_sM_cv_e,$$
(20)
where $`f=\mathrm{}^2/4r_s^2`$ is the solid angle of the shell that is subtended by the collapsing halo, $`\mathrm{}`$ is the radius of the collapsing region when the shock moves through its center, $`M_s=\frac{4\pi }{3}\mathrm{\Omega }_b\mathrm{\Omega }(z)/\mathrm{\Omega }_0\rho _cr_s^3`$ is the mass of the material swept up by the shock and $`M_c`$ is the baryonic mass of the collapsing halo. We can estimate the comoving radius of the collapsing halo as
$$\mathrm{}=R_{M_j}(1+\delta _{\mathrm{NL}})^{1/3},$$
(21)
where $`\delta _{\mathrm{NL}}`$ is the nonlinear overdensity of the region, which is given in the spherical collapse model as
$$1+\delta _{\mathrm{NL}}=\frac{9}{2}\frac{(\theta \mathrm{sin}\theta )^2}{(1\mathrm{cos}\theta )^3},$$
(22)
where the collapse parameter $`\theta `$ is given by $`(\theta \mathrm{sin}\theta )^{2/3}\pi ^{2/3}=D(z_{\mathrm{cross}})/D(z_c)`$ where $`z_{\mathrm{cross}}`$ is the redshift at which the shock moves through the center of the halo. Note that in the elliptical collapse case, $`z_c`$ is computed according to Eq. (4), slowing down the condensation of the cloud somewhat. Note also that we treat halos with subcondensations that have already virialized in the same way. In this case the shock would lose energy impinging on the condensed sub-objects but would also have less uncondensed gas to accelerate to escape velocity. We assume for our purposes here that these effects roughly compensate, although more detailed hydrodynamical studies would help to refine this comparison.
Finally, we calculate the metallicity of each collapsing sphere by volume averaging the contributions of each of the expanding blast waves that pass within the collapse radius $`r_{\mathrm{col}}`$. Each halo is assigned a collapse mass in metals $`M_{Z,k}^{M_i}`$, which is taken to be zero initialy and modified by each blast front that passes within the collapse radius. For each such occurrence, the mass in metals is updated to
$$M_{Z,k}^{M_i}M_{Z,k}^{M_i}+\frac{V_{\mathrm{overlap}}}{\frac{4\pi }{3}R_{M_i}^3}M_Z^{\mathrm{blast}}$$
(23)
where $`V_{\mathrm{overlap}}`$ is the volume of intersection of the two spheres. By dividing this mass by the total baryonic mass of the galaxy we can compute the initial metallicity of the object. This value is then used to compute the collapse time of the object by using the appropriate cooling function, $`\mathrm{\Lambda }(T,Z)`$.
## 4. Results
In this section we summarize the results of our simulations, both in relation to the suppression of low-mass galaxy formation by baryonic stripping of pre-virialized halos, as well as the pre-enrichment of high mass galaxies by metals carried in galactic outflows. We show that the tendency for small galaxies to disrupt the formation of their neighbors helps to explain the small number of Milky Way satellites relative to Cold Dark Matter model predictions, and leads to a lack of small-scale galaxy mergers, suggesting a bell-shaped luminosity function for elliptical galaxies. Pre-enrichment helps to explain the high G-dwarf metallicities found in our own and other galaxies, but the shorter cooling times of metal enriched clouds seems to have little effect on the overall number of large galaxies formed. We examine the effect of varying model parameters on these features and discuss the compatibility between the star formations rates and IGM heating in our models and constraints from optical observations and measurements of thermal distortions in the cosmic microwave background.
### 4.1. Shocks and Low-Mass Galaxy Formation
Having constructed a distribution of halos with collapse times and spatial orientation in good agreement with analytical expressions, we now consider the impact of outflows and enrichment on galaxy formation within this distribution. We first restrict our attention to the $`\mathrm{\Lambda }`$CDM model of Figure 1, as the currently favored cosmological model, and consider the effects of varying the cosmology and other model parameters in §4.4. In Figure 6 we show the collapse and formation of outflows in our fiducial model with $`ϵ_{\mathrm{wind}}=0.1`$. Here we see that regions effected by galactic outflows are highly correlated with regions in which new objects are collapsing, illustrating the necessity of an inhomogeneous approach.
In Figure 7 we show the volume fraction in the expanding shells as a function of $`z`$. This quantity is not estimated on a cell-by-cell basis, but rather approximated by its value in the case in which the positions of the bubbles are uncorrelated,
$$\mathrm{F}(z)1\mathrm{exp}\left(\frac{\frac{4\pi }{3}_{i=1}^NR_i^{\mathrm{blast}}(z)^3}{(12\mathrm{M}\mathrm{p}\mathrm{c}/h)^3}\right).$$
(24)
Note that this quantity overestimates the volume within outflows as the correlations between outflow positions can be significant, and are even greater in a full N-body approach than in our simple Monte Carlo.
Also shown on this plot is the volume-averaged temperature within the shells and the overall volume-averaged temperature of the simulation. Here we see that at redshifts much greater than $``$ 5, the areas near dwarf starburst galaxies were significantly hotter than the IGM as a whole. This discrepancy is somewhat reduced at later times, when the outflows begin to fill a significant fraction of the total volume of the universe. Even at these late times, however, the universe remains inhomogeneous, with new starbursting dwarf galaxies forming within regions only mildly impacted by the earliest outflowing objects. Thus the root mean squared (RMS) fluctuation in the temperature remains significant at redshifts $``$ 1.
This high degree of inhomogeneous shocking naturally has a large impact on galaxy formation at masses below $`10^{10}M_{}`$. In Figure 8 we plot the number of collapsed halos as a function of redshift, along with the number of unsuppressed objects as described in §3.4. While the earliest dwarf galaxies form along with the collapsing halos, by $`z12`$ outflows become important. As $`M10^9M_{}`$ halos collapse after this redshift, galaxy formation at this scale is highly suppressed, leading to a “mass desert” between the major population of polluting dwarf galaxies and the larger galaxies that are seen today. It is also clear from these figures that galaxy formation is likely to have occurred in two stages, with the dwarf outflow population forming mostly at the highest redshifts, and larger galaxies forming only later when halos of sufficient mass began to collapse.
Also in this figure, we plot the total number density of collapsed halos with masses between $`1.1\times 10^8M_{}`$ and $`1.4\times 10^{11}M_{}`$, along with the total number of galaxies in this mass range. Notice that only $`30\%`$ of the halos are populated, consistent with the factor of $`4`$ suppression needed to reconcile the number of predicted and observed Milky-Way satellites (Kauffmann, White, & Guiderdoni 1993; Klypin et al. 1999; Moore et al. 1999). Our scenario thus provides a natural mechanism for the formation of “dark halo” satellites around our galaxy, which may be associated with the abundant High-Velocity Clouds as discussed by Blitz et al. (1999).
### 4.2. Enrichment and High-Mass Galaxy Formation
In Figure 9, we plot the volume-averaged metallicity within the outflows and the simulation overall, as well as the RMS variation between the shells. Here we see that like the temperature, the mean metallicity within the outflows is the greatest at early times, and decreases to a few times the mean metallicity of the universe at a redshift of $`5`$. The mean IGM metallicity from $`z=1`$ to $`z=5`$ ranges from about $`\mathrm{log}_{10}[Z/Z_{}]1.8`$ to $`\mathrm{log}_{10}[Z/Z_{}]2.6`$ to in good agreement with observed metallicities in Ly absorption systems (e.g. Cowie & Songaila 1998, McDonald et al 1999, Ellison et al. 1999), as well as previous studies of enrichment by dwarfs (Nath & Trentham 1997). The RMS variations in the metallicities within the outflows are even more severe than the temperature variations however, and remain several times larger than the mean metallicity at all times.
This highly inhomogeneous distribution of metals is consistent with observed metallicity inhomogeneities as discussed in §2.2 and suggests that the mean metallicity of the gas out of which most galaxies formed may be significantly higher than that of the IGM. Figure 9 also shows the initial metallicity of the IGM newly formed galaxies. Note that this does not include the contribution from the metals in the progenitor galaxies themselves. Here we see that indeed the mean metallicity of collapsed objects is an increasing function of mass which is expected to be $`\mathrm{log}_{10}[Z/Z_{}]1`$ for Milky-Way sized objects, and higher than the mean IGM metallicity at for all mass scales. This value is comparable to the $`0.1Z_{}`$ pre-enrichment necessary to explain the lack of low-metallicity G-dwarfs in the Solar neighborhood. The overall relatively high initial metallicity is also consistent with observations that suggest the need for a “floor” of around $`0.1Z_{}`$ in the metallicity distribution of stars in other massive galaxies. Note that we make no attempt to trace subsequent star formation and interstellar medium (ISM) enrichment, and our values only provide a lower limit on the lowest metallicity stars in disk galaxies that formed largely from the IGM gas. Note also that these metallicities are slowly increasing with mass, suggesting an alternative explanation for the mean metallicity luminosity relationship, which is usually understood as ISM enrichment within galaxies and their progenitors (see eg. Vader 1986; Ferguson & Binggeli 1994; Kauffmann & Charlot 1998).
As collapsing halos in which galaxy formation is suppressed are found closest to outflowing galaxies, where both shock velocities and metallicities are high, it is likely that whatever residual gas remains in these objects would be of similar of even higher metallicity than that of unsupressed objects. While the calculation of these metallicities is beyond the scope of our simulations, it is interesting to note that some of the high velocity clouds may also have metallicities that actually exceed those of many dwarf galaxies although may still be somewhat less than suprasolar values predicted by a “galactic-fountain” model (Sembach et. al 1999; Wakker et al. 1999).
In Figure 10 we show the effect of cooling on galaxy formation in large-mass halos. Here we see that while halo collapse and galaxy formation are almost simultaneous at masses $`5\times 10^{10}M_{}`$, long cooling times suppress all galaxy formation on scales $`5\times 10^{12}M_{}`$. The increased initial metallicities help to accelerate galaxy formation in objects between these mass limits somewhat, although the overall final number densities remain unchanged.
### 4.3. Properties of Elliptical Galaxies
As elliptical galaxies tend to be found in the most dense and enriched regions of space, it is natural to expect that outflows would have had the largest impact on these objects. In order to study this connection, we follow the conventional wisdom that ellipticals correspond to mergers of (see e.g., Barnes 1992; Hernquist 1993) large progenitors.
We therefore define $`M_{\mathrm{merger},k}^{M_i}`$ as the total contribution to the mass of halo $`k`$ of mass $`M_i`$ due to mergers of unsupressed galaxies of mass scales $`M_{i1}`$ and $`M_{i2}`$, that is
$`M_{\mathrm{merger},k}^{M_i}`$ $`{\displaystyle \underset{l}{}}\{\begin{array}{cc}M_{i1}\hfill & \text{if }𝐱_k^{M_i}𝐱_l^{M_{i1}}R_{M_i}\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}`$ (25)
$`+{\displaystyle \underset{l}{}}\{\begin{array}{cc}M_{i2}\hfill & \text{if }𝐱_k^{M_i}𝐱_l^{M_{i2}}R_{M_i}\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}`$
where we consider only objects that have cooled and have not been swept away by shocks while forming. We then identify as ellipticals all collapsed objects with $`M_{\mathrm{merger},k}^{M_i}0.5M_i`$, in which half of the total mass comes from a merger of large progenitors.
We fix this threshold in order to approximate the observed field elliptical fraction of $`15\%`$ (Baugh, Cole, & Frenk 1996), and ask the question of what mass scales correspond to these objects.
In Figure 11 we plot the total space density of galaxies in our simulations, along with the space density of galaxies identified as ellipticals. In this figure we see that the mass distribution of ellipticals is quite different than that of the overall galaxy population. As the outflows from dwarf galaxies suppress the formation of nearby dwarfs, there are very few mergers at small masses, with virial temperatures much smaller than that typical of outflows. Thus while the number density of halos continues to increase with decreasing mass, the number density of ellipticals is the same at $`1.4\times 10^{11}M_{}`$ and $`4.4\times 10^{10}M_{}`$, and decreases dramatically for even lower masses.
Using the observed Faber-Jackson relation between the luminosity of elliptical galaxies and the velocity dispersion, we can recast this as an absolute r-band magnitude of $`M_r=1810\mathrm{log}_{10}\sigma _{100}`$, (Oegerle & Hoessel 1991) where $`\sigma _{100}`$ is the velocity dispersion in units of $`100`$ km/s. This gives the luminosity function of elliptical galaxies shown in the bottom of Figure 11. While these results are necessarily crude, and are bounded at the high-luminosity end by the finite size of out simulations, they nevertheless naturally reproduce the observed bell-shaped luminosity function of elliptical galaxies (see eg., Bromley et al. 1998).
### 4.4. Effects of Model Uncertainties
We have only two free parameters in our simplified modeling, both relating to the gas outflow, and all other modeling based on a simple spherical or elliptical collapse model, involving the usual assumptions. Our conclusions are most strongly dependent on the product of the two free parameters $`ϵ_{\mathrm{wind}}`$ and $`ϵ_{\mathrm{sf}}`$, which together determine the total energy being channeled into galactic outflows. In our fiducial model, we took $`ϵ_{\mathrm{sf}}`$ = 0.1 and $`ϵ_{\mathrm{wind}}ϵ_{\mathrm{sf}}=.01`$. These canonical values are motivated by the current observational work on dwarf galaxy outflow energetics, but a large spread is naturally expected and thus we have some freedom in changing this efficiency. We therefore examined two extreme cases with $`ϵ_{\mathrm{wind}}=.05`$ and $`ϵ_{\mathrm{wind}}=.2`$ to asses the robustness of the features described above. These cases bound the shaded regions of the results plotted in Figures, 7, 8, 11, 12.
In Figure 7 we see the impact of model uncertainties on the total volume fraction and temperatures of outflows in our simulations. The impact of these differences on the suppression of dwarf galaxies is shown in the shaded regions in Figure 8. Here we see that the existence of two-stage evolution with a “mass desert” in the $`10^910^{10}M_{}`$ range persists for all values of these parameters. The total number of unsuppressed dwarf galaxies ranges from $`10\%30\%`$ reinforcing the idea that outflow suppression may be important in reproducing the small number of Milky-Way satellites.
The effect of varying $`ϵ_{\mathrm{sf}}ϵ_{\mathrm{wind}}`$ on the distribution of elliptical galaxies is shown in Figure 11. While the uncertainties in this case are somewhat more severe than in Figure 8, a large fall-off in the number of objects with masses below $`5\times 10^{10}M_{}`$ is present for all cases.
We show the metallicities of the IGM component of the forming galaxies in each of these cases in the upper panel of Figure 12. In both cases the mean metallicity of collapsed objects is an increasing function of mass which is higher than the mean IGM metallicity. The predicted Milky-Way scale IGM metallicity is again always $`\mathrm{log}_{10}[Z/Z_{}]1`$ and thus the G-dwarf problem is easily understood.
In the lower panel of this figure, we compute the star formation rate in dwarf galaxies to assuage any concern that our code might be pumping an inordinate amount of energy into the IGM. Note that our models do not include subsequent star-formation in disk galaxies and thus this figure represents a lower bound to the number of stars formed at each redshift. Nevertheless, our values are well within observational constraints (see eg., Adelberger & Steidel 2000).
We can also compare our models to the COBE constraints on the optical depth to the surface of last scatter as well as the the distortion of the cosmic microwave background (CMB) spectrum due to the Sunyaev-Zel’dovich effect. We can estimate the optical depth in our model as
$$\tau =5.9\times 10^3\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}_0^{z_{\mathrm{max}}}𝑑z\frac{dx}{dz}\mathrm{\Omega }(z)(1+z)^2\mathrm{F}(z),$$
(26)
where $`x`$ is again the comoving distance defined such that $`a_0H_0=c`$. This analysis gives $`\tau =0.05`$ for the fiducial model, and varies from $`\tau =0.04`$ to $`\tau =0.06`$ in the low and high energy outflow cases, which are all well within the observational limit of $`\tau 0.5`$ (Griffiths, Barbosa, & Liddle 1999).
The degree of CMB spectral distortions is given by the Compton-$`y`$ parameter which is the convolution of the optical depth with the electron temperature along the line of sight (Zel’dovich & Sunyaev 1969; Sunyaev & Zel’dovich 1972). Thus
$$y=1.0\times 10^8\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}_0^{z_{\mathrm{max}}}𝑑z\frac{dx}{dz}T_5(z)\mathrm{\Omega }(z)(1+z)^2\mathrm{F}(z),$$
(27)
where $`T_5(z)`$ is the mean temperature within the outflows as a function of redshift in units of $`10^5`$K. In this case $`y`$ is $`1.6\times 10^6`$ in the fiducial model, and varies from $`1.0\times 10^6`$ to $`2.3\times 10^6`$ for the extreme cases. These values are again within the observational constraint of $`y1.5\times 10^5`$ (Fixsen et al. 1996). Nevertheless, our modeling suggests that the spectral distortions due to galactic outflows should be observable with the next generation of cosmic microwave background experiments.
Varying the second parameter in our model, $`ϵ_{\mathrm{sf}}`$, while keeping $`ϵ_{\mathrm{sf}}ϵ_{\mathrm{wind}}`$ fixed acts as an overall shift in the star formation rate and metallicities in Figures 9 and 12 and can be estimated simply by eye. Including the ionization drag term in Eq. (11) and allowing the shocks to act as the ionizing source of the IGM has little effect on our results as suggested by Figure 5 and verified by a full simulation. Thus we see that all the major features of the fiducial model persist over a wide range of parameters.
We have also examined the cosmological dependence of our results by examining a flat model with parameters as described in §2.3 and $`ϵ_{\mathrm{wind}}=0.1`$. In this case only $`50\%`$ suppression of galaxies occurs, again with the most impact on galaxies in the few times $`10^9M_{}`$ range. In this model elliptical formation falls off below the $`10^{11}M_{}`$ scale but more gradually than in the open case, and initial Milky-Way metallicities again about $`\mathrm{log}_{10}[Z/Z_{}]=1.0`$. Thus even in a model in which structure forms much more quickly than suggested by observations (see eg., Bahcall, Fan, & Cen 1997), the major features in our model persist.
## 5. Discussion
Our treatment of galactic outflows is intended as an exploratory study of these processes, and makes a number of approximations that should be made explicit. While a Monte Carlo approach uncovers many of the issues of inhomogeneity that can not properly be studied analytically, it fails to capture the nonlinear structure and clustering present in a full N-body treatment. At late times peaks will be grouped along sheets and filaments, and thus clustering is likely to be more severe than in our simulations and of a more complicated nature.
A number of uncertainties also arise from our simple treatment of IGM enrichment and shocking. One concern is the presence of radiative feedback from ionization fronts around dwarf galaxies. Our assumption that fronts precede the formation of outflows, and thus each shell can be treated as expanding into an ionized medium is probably a reasonable one, while our placement of a “temperature floor” at the ionization temperature of hydrogen at all times in our simulations is stronger approximation. Fortunately the extreme fragility of molecular hydrogen suggests that the formation of objects with virial temperatures below $``$10,000 K may have not been largely affected by reionization. While the inhomogeneous structure of reionization is likely to have had a great impact on the CMB fluctuations (Aghanim et al. 1996; Miralda-Escude, Haehnelt, & Rees 2000; Scannapieco 2000), the formation of $`10^7M_{}`$ galaxies is likely to have been halted by the dissociation of molecular hydrogen by the first stars, long before reionization took place.
A bigger concern may be the structures of outflows and the ejection fractions of interstellar gas and metals. This has been studied in detail by Mac Low & Ferrara (1999) and Ferrara & Tolstoy (2000), who conclude that efficient ejection of the interstellar medium or “blowaway” occurs only in halos with masses $`10^7M_{}`$. The issue of whether dwarfs retain a sizeable fraction of their gas, however, is to some degree decoupled from the formation of outflows. In halos in the mass range $`10^7M_{}M10^9M_{}`$, a “blowout” occurs in which the super-bubbles around groups of SNeII punch out of the galaxy, shocking the surrounding IGM and efficiently ejecting metals while failing to excavate the interstellar medium of the galaxy as a whole. Such a scheme in which outflows are formed primarily of the enriched IGM surrounding dwarf galaxies is equivalent to outflows of galactic gas for the purposes of our simulations. Additionally such a picture would help to reconcile observations of expanding high metallicity shells around dwarf galaxies as described in §2 with observations of multiple episodes of and HI gas in dwarf spheroidal galaxies (Smecker-Hane et al. 1994; Grebel 1998) some of which even suggest that many of these objects are gas-rich, but with extended HI envelopes (Blitz & Robishaw 2000).
That being said, the TSE model adopted for outflow evolution in our simulations is almost certainly oversimplified. The “blowout” scenario described by Mac Low & Ferrara (1999) ejects matter perpendicular to the thick disk, and is thus quite different from the spherical outflow modeled here. Even if the winds from dwarf galaxies can be reasonably approximated as spherical shells at early times, such outflows of low-density heated gas would necessarily become Rayleigh-Taylor unstable as they expanded into the denser IGM. Thus material is most likely to flow out in a number of heated clumps at large distances.
Also ignored in our model are shell-shell interactions which would contribute additional heating but slow expansion in regions which two shocks meet. Similarly, an additional source of IGM enrichment is the merger mechanism described in Gnedin & Ostriker (1997) and Gnedin (1998), in which a significant fraction of the interstellar medium of merging galaxies is ejected.
Finally, the criteria for halo suppression used in our simulation is over-simplified. While regions of space in our simulations that meet the stripping criterion, Eq. (20), are likely to have delayed or suppressed galaxy formation, the one-to-one correspondence assumed in our simulations is unlikely. In reality the formation of galaxies is likely to be a complicated function of the number and strength of the shocks moving through the regions, and the structure of these regions during and after halo collapse. This subject merits further investigation and would help to sharpen our conclusions.
## 6. Conclusions
While the impact of preheating and enrichment on the observed properties of galaxy clusters has long been recognized, the impact on the observed properties of galaxies themselves has been little explored. By accounting for the observed outflows from dwarf galaxies, we have been able to show how many of the unexplained properties of galaxies and the IGM can be naturally understood.
Firstly, the suppression of low-mass galaxy formation by outflows provides a natural explanation for the factor of $`4`$ discrepancy between the number of observed Milky-Way satellites and predictions from standard CDM models that do not include outflows. Suppression also provides a natural mechanism for the formation of “dark halos” which may be associated with the High-Velocity clouds.
Secondly, baryonic stripping results in a bell-shaped luminosity function of ellipticals. The lower the mass of a halo, the more likely it is to generate an outflow that strips material from a similar mass neighboring pre-virialized halo that would otherwise later form into a galaxy. This results in very few pairs of neighboring low-mass galaxies, and hence a relative deficit of major low mass-mergers.
Finally, our models of enrichment of protogalactic gas predict a trend of increasing metallicity with galaxy mass in good agreement with inferences from observations. The initial metallicity predicted for a Milky-Way mass galaxy is $`0.1Z_{}`$ providing a natural initial floor at the level required to solve the G-dwarf problem and the more general lack of low metallicity stars in well studied massive elliptical galaxies relative to “closed-box” models of chemical enrichment.
These results are persistent over a wide range of model parameters and cosmologies, and are not a result of fine-tuning parameters or invoking additional physics.
While galaxy outflows are already incorporated into modern studies of galaxy formation, this is done only as an internal modification, ignoring pre-enrichment. Keeping track of the effect of outflows on neighboring halos is essential in understanding the properties of galaxy clusters, and hence it is not surprising that these effects would play a major role in the formation of galaxies as well. While the details await further investigation, it is clear that any complete picture of galaxy formation must account for heating and enrichment.
We would like to thank Rychard Bouwens, Andrew Cumming, Julianne Dalcanton, Marc Davis, Andrea Ferrara, Ignacio Ferreras, Brenda Frye, John Huchra, Siang Peng Oh, Alvio Renzini, David M. Sherfesee, Joseph Silk, Jonathan C. Tan, Robert Thacker, and Simon White for helpful comments and discussions. This research was supported in part by the National Science Foundation under Grant PHY94-07194. TJB acknowledges NASA grant AR07522.01-96A. |
warning/0003/cond-mat0003412.html | ar5iv | text | # Symmetry Dependence of Localization in Quasi- 1- dimensional Disordered Wires
## Abstract
The crossover in energy level statistics of a quasi-1-dimensional disordered wire as a function of its length $`L`$ is used, in order to derive its averaged localization length, without magnetic field, in a magnetic field and for moderate spin orbit scattering strength. An analytical function of the magnetic field for the local level spacing is obtained, and found to be in excellent agreement with the magnetic field dependent activation energy, recently measured in low-mobility quasi-one-dimensional wires. This formula can be used to extract directly and accurately the localization length from magnetoresistance experiments. In general, the local level spacing is shown to be proportional to the excitation gap of a virtual particle, moving on a compact symmetric space. Pacs- numbers: 72.15.Rn,73.20.Fz,02.20.Qs
In disordered wires quantum interference results in localization of all states for arbitrary disorder strength, if the wire is infinite. It has been discovered that the localization length depends on the global symmetry of the wire : $`L_c=\beta \pi \mathrm{}\nu SD_0`$, where $`\beta =1,2,4`$, corresponding to no magnetic field, finite magnetic field, and strong spin- orbit scattering or magnetic impurities, respectively. $`\nu (E)`$ is the electronic density of states in the wire. $`D_0`$ is the classical diffusion constant of the electrons in the wire, and $`S`$ its crossection. This result was obtained by calculating the spatial decay of the density correlation function for wires whose thickness exceeds the mean free path $`l`$. Independently, it was obtained by calculating the transmission probability through thin, few channel wires. Recently, the doubling of the localization length was observed in sub-micron thin wires of doped GaAs by Khavin, Gershenson and Bogdanov, who found a continously decreasing activation energy when the magnetic field is increased, saturating indeed at one half of its field free value . However, the crossover as function of moderate symmetry breaking fields defied any attempt to study it with nonperturbative methods. The only results were based on a heuristic approach by Bouchaud, a kind of semiclassical analysis by Imry and Lerner and numerical studies. Recently, however, Kolesnikov and Efetov succeeded to tackle that complex problem and derived the density- density correlation function with the supersymmetry method in the crossover regime. The result does not scale with a single localization length.
Here, the problem of the crossover is addressed from a different perspective. The statistics of discrete energy levels of a finite coherent, disordered metal particle is an efficient way to characterize its properties . This can be studied by calculating a disorder averaged autocorrelation function between two energies at a distance $`\omega `$ in the energy level spectrum. It is an oscillatory function whose amplitude decays with a power law when the energy levels in the vicinity of the central energy $`E`$ are extended. A Gaussian decaying function is a strong indication that all states are localized. The autocorrelation function of spectral determinants (ASD) is defined by $`C(\omega )=\overline{C}(\omega )/\overline{C}(0),\overline{C}(\omega )=\text{det}(E+\omega /2H)\text{det}(E\omega /2H).`$ $`E`$ is a central energy. We consider the hamiltonian of disordered electrons
$$H=\left[𝐩q𝐀\right]^2/2m+V(𝐱),$$
(1)
where $`q`$ is the electron charge, $`c`$ the velocity of light. $`V(𝐱)`$ is taken to be a Gaussian distributed random function $`V(𝐱)=0`$, and $`V(𝐱)V(𝐱^{})=(\mathrm{\Delta }SL/2\pi \tau )\delta (𝐱𝐱^{}),`$ which models randomly distributed, uncorrelated impurities in the sample. $`\mathrm{\Delta }`$ is the mean level spacing, $`1/\tau `$ the elastic scattering rate. The vector potential is used in the gauge $`𝐀=(By,0,0)`$, where $`x`$ is the coordinate along the wire of length L, $`y`$ the one in the direction perpendicular both to the wire and the magnetic field $`𝐁`$, which is directed perpendicular to the wire. The angular brackets denote averaging over impurities.
The ASD appears as an intermediate step in the Grassman replica trick and as the compact sector of the supersymmetric field theory of disordered systems . It is a non-self-averaging quantity. It has been shown recently, that the ASD can distinguish between an uncorrelated spectrum of localized states and a correlated spectrum of extended states in the vicinity of $`E`$. This way, the metal- insulator crossover as a function of the length of a disordered wire has been explored with the ASD. A localization length $`\xi _c`$ in a moderate magnetic field has been derived as the crossover length scale and shown to coincide with $`L_c(\beta =2)`$. The ASD also yields qualitative information on the location of localized states in a quantum- Hall- system .
The impurity averaged ASD can be written as a partition function
$$\overline{C}(\omega )=\text{Tr}\mathrm{exp}(L\overline{H}\left[Q\right]),$$
(2)
where $`\overline{H}`$ is an effective Hamiltonian of matrices Q on a compact manifold, determined by the symmetries of the Hamiltonian $`H`$ of disordered electrons. Thus, the problem reduces to the one of finding the spectrum of the effective Hamiltonian $`\overline{H}`$.
There is a finite gap $`E_G`$ between the ground state energy and the energy of the next excited state of $`\overline{H}(\omega =0)`$. For a long wire, $`LE_G1`$, the ASD becomes, according to Eq. (2), $`C(\omega )=\mathrm{exp}(const.L\omega ^2/E_G)`$. This is typical for a a spectrum of localized states. In the other limit $`LE_G1`$, all modes of $`\overline{H}`$ do contribute to the trace in the partition function Eq. (2) with equal weight, yielding the correlation function of a spectrum of extended states. Thus, the crossover length $`\xi _c1/E_G`$ can be identified with an averaged localization length.
Being a product of two spectral determinants, $`4\alpha `$ -component Grassman fields are needed to get the functional integral representation of the ASD. Here, $`\alpha =1`$, when the Hamiltonian is independent of the spin of the electrons, and each level is doubly spin degenerate. There is one pair of Grassman fields for each determinant in the ASD and each pair is composed of a Grassman field and its time reversed one, as obtained by complex conjugation. $`\alpha =2`$ has to be considered, when the Hamiltonian does depend on spin, as for the case with moderately strong magnetic impurity or spin- orbit scattering. This necessitates the use of a vector of a spinor and the corresponding time reversed one. One notes a global invariance under rotations of these vectors.
Averaging over impurities according to Eq. (1), one gets an interacting theory of Grassman fields with interaction strength $`1/\tau `$. It can be decoupled with a Gaussian integral over $`4\alpha \times 4\alpha `$ matrices which preserve the global rotational invariance. The Grassman fields can now be integrated out exactly. Finding the saddle point of the integral over matrices and integrating out longitudinal (massive) modes, for $`\omega 1/\tau `$ and $`1/\tau \mathrm{\Delta }`$, the ASD reduces to a functional integral over transverse modes Q, being elements of respective symmetric spaces. For disordered wires whose thickness $`W`$ is larger than the mean free path $`l`$ but smaller than the magnetic length $`W<l_B`$, the action reduces to the one of the one-dimensional compact nonlinear sigma model ,
$$\overline{C}(\omega )=\underset{x}{}\text{d}Q(x)\mathrm{exp}(F[Q]),$$
(3)
$`F[Q]`$ $`=`$ $`\alpha {\displaystyle \frac{1}{16}}L_{CU}{\displaystyle _0^L}\text{d}x\text{Tr}\left[(_xQ(x))^2{\displaystyle \frac{q^2}{\mathrm{}^2}}\overline{y^2}B^2[Q,\tau _3]^2\right]`$ (4)
$`+`$ $`i\alpha {\displaystyle \frac{\pi }{4}}{\displaystyle \frac{\omega }{\mathrm{\Delta }}}{\displaystyle \frac{\text{d}x}{L}\text{Tr}\mathrm{\Lambda }_3Q(x)}.`$ (5)
where $`L_{CU}=L_C(\beta =2)=2\pi \mathrm{}\nu SD_0`$ is the localization length in the wire in a moderately strong magnetic field . The integral over the crossection $`S`$ of the wire is done, giving $`\overline{y^2}`$. Here, and in the following, $`\mathrm{\Lambda }_i`$ are the Pauli matrices in the subbasis of the left and the right spectral determinant, $`\tau _i`$ the ones in the subbasis spanned by time reversal and $`\sigma _i`$ the ones in the space spanned by the spinor, for $`i=1,2,3`$. The fluctuations of the matrices $`Q`$ are transverse, $`Q^2=1`$, and restricted by global rotational invariance to $`Q^+=Q`$. For $`\alpha =1`$ in addition $`Q^TC=CQ`$, where $`C=i\tau _2`$. This constrains $`Q`$ to be an element of the group defined on the compact symplectic symmetric space Sp(2)$`/`$( Sp(1) $`\times `$ Sp(1)). For a moderate magnetic field, $`Q`$ is reduced to a $`2\times 2`$\- matrix by the broken time reversal symmetry. This reduces the space of Q to $`U(2)/(U(1)\times U(1))`$. For $`\alpha =2`$ the matrix $`C`$ is, due to the time reversal of the spinor, substituted by $`i\sigma _2\tau _1`$. Both magnetic impurities and spin-orbit scattering reduce the Q matrix to unity in spin space. Thus, C has effectively the form $`\tau _1`$. The condition $`Q^TC=CQ`$ leads therefore to a new symmetry class, when the spin symmetry is broken but the time reversal symmetry remains intact. This is the case for moderately strong spin-orbit scattering. Then, $`Q`$ are $`4\times 4`$\- matrices on the orthogonal symmetric space $`O(4)/(O(2)\times O(2))`$ . With magnetic impurities both the spin and time reversal symmetry is broken, and the Q- matrices are in the unitary symmetric space $`U(2)/(U(1)\times U(1))`$ as for a moderate magnetic field and spin degenerate levels. The difference in the prefactor $`\alpha `$ in the action, Eq.(4), remains. One can extend this approach to other compact symmetric spaces with physical realizations, see Ref. for a complete classification. For $`\omega /\mathrm{\Delta }<L_{CU}/L`$, the spatial variation of $`Q`$ can be neglected and one retains the same ASD as for random matrices of orthogonal, unitary and symplectic symmetry, characterizing the energy levels of an ergodic particle without magnetic field, with magnetic field, and with spin- orbit- scattering, respectively. The confusing fact, that random orthogonal matrices result for the ASD in a functional integral over compact symplectic Q- matrices and vice versa, results from the sign change of the Grassman variables under time reversal.
We can derive the corresponding Hamiltonian $`\overline{H}`$ by means of the transfer matrix method, reducing the one-dimensional integral in Eq. (3) to a single functional integral. Thus, the ASD is obtained in the simple form of Eq. (2), with the effective Hamiltonian
$$\overline{H}(\omega =0)=\frac{1}{\alpha L_{CU}}(4\mathrm{\Delta }_Q\frac{1}{16}X^2Tr_Q[Q,\tau _3]^2).$$
(6)
$`\mathrm{\Delta }_Q`$ is that part of the Laplacian on the symmetric space, which couples to $`Tr[\mathrm{\Lambda }_3Q]`$. $`X=2\pi \alpha L_{CU}(\overline{y^2})^{(1/2)}B/\varphi _0`$, where $`\varphi _0=q/h`$ is the flux quantum.
The problem is now equivalent to a particle with “mass” $`(\alpha /8)L_{CU}(E)`$ moving on the symmetric space of $`Q`$ in a harmonic potential with “frequency” $`X/(\alpha L_{CU})`$, and in an external field $`i\alpha (\pi /4)\omega /(L\mathrm{\Delta })`$, in “time” $`x`$, the coordinate along the wire. To find the ASD as a function of $`\omega `$ and the length of the wire $`L`$, one can do a Fourier analysis in terms of the spectrum and eigenfunctions of the effective Hamiltonian at zero frequency, $`\overline{H}(\omega =0)`$ . For $`L\xi _c`$, only the gap between the ground and the first excited state of $`\overline{H}(\omega =0)`$, determines the ASD.
Without magnetic field, $`B=0`$, the Laplacian is
$`\mathrm{\Delta }_Q`$ $`=`$ $`_{\lambda _C}(1\lambda _C^2)_{\lambda _C}+2{\displaystyle \frac{1\lambda _C^2}{\lambda _C}}_{\lambda _C}`$ (7)
$`+`$ $`{\displaystyle \frac{1}{\lambda _C^2}}_{\lambda _D}(1\lambda _D^2)_{\lambda _D},`$ (8)
where $`\lambda _{C,D}[1,1]`$. Its ground state is $`1`$ and the first excited state is $`\lambda _C\lambda _D`$. Thus, the gap is
$$E_G(B=0)=16/L_{CU}.$$
(9)
For moderate magnetic field, with the condition $`L_{CU}(\overline{y^2})^{(1/2)}B\varphi _0=h/q`$, all degrees of freedom arising from time reversal invariance are frozen out, due to the term $`Tr_Q[Q,\tau _3]^2=16(\lambda _C^21)`$ which fixes $`\lambda _C^2=1`$. Then, the Laplacian reduces to
$`\mathrm{\Delta }_Q=_{\lambda _D}(1\lambda _D^2)_{\lambda _D}.`$ (10)
Its eigenfunctions are the Legendre polynomials. There is a gap above the isotropic ground state of magnitude
$$E_G(X1)=8/L_{CU}.$$
(11)
For moderate magnetic impurity scattering exceeding the local level spacing, $`1/\tau _s>\mathrm{\Delta }_C`$, the Laplacian is given by Eq.(8). Thus, due to $`\alpha =2`$, the gap is reduced to $`E_G(1/\tau _S>\mathrm{\Delta }_C)=4/L_{CU}`$. For moderately strong spin- orbit scattering $`1/\tau _{SO}>\mathrm{\Delta }_C`$, the Laplace operator is
$$\mathrm{\Delta }_Q=\underset{l=1,2}{}_{\lambda _l}(1\lambda _l^2)_{\lambda _l},$$
(12)
where $`\lambda _{1,2}[1,1]`$. The ground state is $`\psi _0=1`$, the first excited state is doubly degenerate, $`\psi _{11}=\lambda _1`$, $`\psi _{12}=\lambda _2`$. Thus, the gap is the same as for magnetic impurities,
$$E_G(1/\tau _{SO}>\mathrm{\Delta }_C)=4/L_{CU}.$$
(13)
An external magnetic field lifts this degeneracy but does not change the gap.
Using the crossover in energy level statistics as the definition of a localization length as above, we get in a quasi- 1 -dim. wire,
$$\xi _c=1/E_G(\beta )=(1/16)\beta L_{CU},$$
(14)
where $`\beta =1,2,4`$ corresponding to no magnetic field, finite magnetic field, and strong spin- orbit scattering or magnetic impurities, respectively. Comparing with the known equation for the localization length, $`L_c`$, we find that the dependence of the ratios $`\beta `$ on the symmetry are in perfect agreement with the result as obtained from the spatial decay of the density- density- correlation function, while it defers by the overall constant $`1/8`$.
This relation can be proven directly. The ASD at zero frequency $`\overline{C}(0)_L`$ of the wire of length $`L`$, becomes, when the wire is divided into two parts, $`\overline{C}(0)_{L/2}^2`$. For $`L\mathrm{}`$, we find that the relative difference is:
$$f(L)=\frac{\overline{C}(0)_{L/2}^2}{\overline{C}(0)_L}1=2\mathrm{exp}(LE_G/2),$$
(15)
exponentially decaying with the length $`L`$. $`f(L)`$ can be estimated, following an argument by Mott: When the two halves of the wire get connected, the Eigenstates of the two separate halves become hybridized and the Eigenenergy of a state $`\psi _n`$ is changed by $`\pm \mathrm{\Delta }_C\mathrm{exp}(2x_n/L_C)`$. $`x_n`$ is random, depending on the position of an eigenstate with closest energy in the other half of the wire. Thus, averaging over $`x_n`$ gives:
$$f(L)+\mathrm{exp}(4L/L_C).$$
(16)
Comparison with Eq. (15) yields indeed $`1/L_C=8E_G`$.
The Crossover Behaviour of the Localization Length
This direct relation of the ASD to the spatial overlap between wavefunctions allows us to extend this approach to get an analytical solution for the crossover behaviour of the localization length and the local level spacing as a magnetic field is turned on, and without strong spin- orbit scattering. While a heuristic approach and numerical studies seemed to indicate a continous increase of the localization length, the analytical result does show that both limiting localization lengths $`L_c(\beta =1)`$ and $`L_c(\beta =2)`$ are present in the crossover regime and that there is no single parameter scaling. This is explained by arguing that the far tails of the wavefunctions do cover a large enough area to have fully broken time reversal symmetry, decaying with the length scale $`L_c(\beta =2)`$ even if the magnetic field is too weak to affect the properties of the bulk of the wavefunction, which does decay at smaller length scales with the shorter localization length $`L_c(\beta =1)`$, corresponding to the time reversal symmetric case. The quantity studied there is the imprurity averaged correlation function of local wavefunction amplitudes and its momenta at a fixed energy $`ϵ`$: $`Y(ϵ)=<_\alpha \psi _\alpha (0)^2\psi _\alpha (r)^2\delta (ϵϵ_\alpha )>`$. It is averaged over a distribution of eigenfunctions in different impurity representations. Thus, each eigenfunction could decay exponentially with a single localization length, but with a distribution which has two maxima, at $`L_c(\beta =1)`$ and $`L_c(\beta =2)`$, whose weight is a function of the magnetic field in the crossover regime. It is interesting to ask if this property of the eigenfunctions does have a consequence on the energy level statistics as well, or, if that is still governed by a single parameter as the magnetic field is varied.
The effective Hamiltonian for moderate magnetic fields is given, without spin dependent scattering, $`\alpha =1`$, by:
$$\overline{H}=\frac{1}{L_{CU}}(4\mathrm{\Delta }_Q+X^2(1\lambda _C^2)),$$
(17)
where the Laplacian is Eq. and $`X=2\pi \varphi /\varphi _0`$, where $`\varphi `$ is the magnetic flux through an area $`A=L_{CU}(\overline{y^2})^(1/2)`$.
In the limit $`X0`$ the ground state and first excited state approach $`1,\lambda _C\lambda _D`$, respectively. In the limit $`X1`$, $`\lambda _C^2`$ becomes fixed to 1. Thus, the Ansatz $`\psi _0(\lambda _C)=\mathrm{exp}(A_0(1\lambda _C^2))`$, and $`\psi _1(\lambda _C,\lambda _D)=\lambda _C\lambda _D\mathrm{exp}(A_1(1\lambda _C^2))`$, solves $`\overline{H}\psi =\overline{E}\psi `$ to first order in $`z=X^2(1\lambda _C^2)`$ and yields the gap:
$$E_G(X)=4(2+\sqrt{49+X^2}\sqrt{25+X^2})/L_{CU}.$$
(18)
This solution is valid in both the limits $`X1`$ and $`X1`$, interpolating the region $`X1`$. The resulting ratio of local energy level spacings $`\mathrm{\Delta }_C(B)/\mathrm{\Delta }_C(0)=E_G(B)/E_G(0)`$, is shown in Fig. 5 to be in excellent qualitative agreement with experimental data for the magnetic field dependent activation energy, measured recently in transport experiments.
There is a quantitative discrepancy by a factor $`.6`$ between the best fit $`X=.036H/Oe`$, and $`X=2\pi \varphi /\varphi _0`$, $`\varphi =\mu _0HL_{CU}(\overline{y^2})^{(1/2)}`$. With the experimental parameters $`\alpha =1,L_{CO}=.61\mu m`$, width $`W=.2\mu m`$ of sample 5 in Ref. and $`\overline{y^2}=W^2/12`$ for a 2- dimensional wire, it yields $`X=.021H/Oe`$. We note that smooth confinement can give $`\overline{y^2}>W^2/12`$, thus explaining this discrepancy and the observed difference between $`W`$ as obtained from the sample resistance and estimated from the analysis of the weak localization magnetoresistance, which also depends on $`\overline{y^2}`$. When $`W^2B>\varphi _0`$ or $`H>H^{}=2.110^{11}Oe[m^2/W^2]=525Oe`$, the derivation has to be extended to the 2- dimensional nonlinear sigma- model, and Eq. 18 seizes to be valid. Then, $`L_C`$ increases as function of magnetic field more strongly.
The asymptotic behaviour, of $`\delta L_C(B)B^2`$ for small and $`1/B`$ at large magnetic fields of Eq. (18) does agree with the suggestion by Bouchaud.
The two length scale physics of Ref. has thus no consequence for the energy level statistics as studied with the ASD. Since we could also show a direct relation between this spectral statistics and the exponential decay of localized eigenfunctions, it is suggestive to be the non self averaging property of the ASD which washes out the two scale physics.
The author gratefully acknowledges, usefull comments by Isa Zarekeshev, discussions with Bernhard Kramer and support from DFG. He thanks Yuri Khavin for usefull communications and Martin Zirnbauer for pointing out a serious confusion in the notation. |
warning/0003/cond-mat0003191.html | ar5iv | text | # Cu-O network dependent core-hole screening in low-dimensional cuprates: A theoretical analysis
## I Introduction
The influence of dimensionality and lattice geometry on the electronic structure of crystals is a classical topic in solid state physics. However, in the presence of strong electronic correlations like those observed in the high-$`T_c`$ superconducting cuprates, the theoretical analysis of electronic properties is especially difficult. Recently, the electronic structure of various undoped cuprates has been studied experimentally using core-level X-ray photoemission spectroscopy (XPS). By comparing these spectra to the results of theoretical model calculations we can try to improve our understanding of the electronic properties of the cuprates. In addition, a theoretical analysis allows us to check the validity of the microscopic models used for the calculations, and it provides values for the physical parameters which are contained in the model Hamiltonians.
The core-level photoemission process leads to a final state in which an electron is missing in the $`2p_{3/2}`$ core orbital of a Cu site. The resulting positive charge is screened by valence electrons, predominantly from Cu $`3d_{x^2y^2}`$ and O $`2p_{x,y}`$ orbitals. The contribution of these screening processes to the final state determines the detailed form of the experimental spectrum. Therefore, core-level spectra contain information about the dynamics of valence electrons.
Since the atomic ground state of the Cu-O system is a $`2p^63d^9`$ configuration, it is advantageous to use the hole picture. In this case the core-level photoemission process amounts to the creation of a core hole which has a strong repulsive Coulomb interaction with the valence hole in the $`3d_{x^2y^2}`$ orbital of the same Cu site.
Experimentally, the Cu $`2p`$ core-level spectrum of many formally divalent Cu compounds shows a dominant main line and a pronounced satellite structure. In analogy to an approach for metals, this satellite structure was assigned to a (“poorly screened”) final state (denoted $`3d^9`$) in which the valence hole largely remains on the Cu site. In this way, the structure of the satellite could be explained by multiplet splitting due to the remaining $`d`$-hole.
The main line was interpreted as originating from a final state in which the hole resides in the neighbouring ligands (denoted $`3d^{10}`$L). However, for several compounds the main line is found to be asymmetric, and to have a large width. For instance, three substructures were identified in the main line of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> and Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> (which will be analyzed in this paper). Substructures of this kind are the reason why the main line in the Cu $`2p_{3/2}`$ spectra cannot be attributed to a local $`3d^{10}`$L excitation only: This interpretation does not account for the asymmetry and the large width of the main line.
An alternative explanation was found using numerical exact diagonalization for a chain of three plaquettes (Cu<sub>3</sub>O<sub>10</sub>). It was shown that the valence hole could delocalize further in the crystal. This leads to a lowest eigenstate of $`3d^{10}`$ character, in which the hole is mainly pushed out onto the neighbouring CuO<sub>4</sub> units, forming a Zhang-Rice singlet. These delocalization processes in the hole picture are equivalent to the screening processes in the electron picture.
Obviously, this screening of the core hole depends on the geometry of the CuO network. This suggests the possibility of finite-size effects in cluster calculations. In fact, convergence of the spectra with respect to the system size was found only for chain clusters with lengths of seven plaquettes (Cu<sub>7</sub>O<sub>21</sub>). These calculations led to a satisfactory agreement with the experimental results for Sr<sub>2</sub>CuO<sub>3</sub>.
Another popular model for the description of core-level spectroscopy is the single-site Anderson impurity model. Recently, the influence of dimensionality on Cu $`2p_{3/2}`$ spectra of several cuprates was analyzed using this model. Overall, very good agreement with the experimental main lines was obtained. The satellite structures, however, were not discussed. In particular, no multiplet effects were included. Furthermore, the absolute energetical positions of the main lines of Sr<sub>2</sub>CuO<sub>3</sub> and Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> were not correctly reproduced. These shortcomings are avoided in our treatment.
In the present paper, we discuss the influence of of the dimensionality on the Cu $`2p_{3/2}`$ core level spectra of undoped cuprates by means of the Mori-Zwanzig projection technique. Both exchange splitting and delocalization properties are described within one framework using a multi-band Hubbard model. In addition, a clear description of the final states is obtained.
Our analysis allows to distinguish between effects of the Cu-O geometry and effects from material-specific properties. Examples for material-specific properties are the Cu-O distance or components which do not belong to the Cu-O network. Some of these properties are described by the values of model parameters like the Cu-O charge-transfer energy or the Cu-O hybridization strength. A phenomenological analysis cannot conclusively determine whether an effect is due to dimensionality or not. For instance, due to differences in material-specific properties a peak which is brought about by a similar process in two different compounds may not appear at the same binding energy in both spectra. Therefore, the investigation of microscopic models is necessary.
We compare our calculations with the experimental spectra of three materials which consist of different Cu-O networks: Bi<sub>2</sub>CuO<sub>4</sub> which contains separated CuO<sub>4</sub> plaquettes (“zero”-dimensional Cu-O<sub>4</sub> network), Sr<sub>2</sub>CuO<sub>3</sub> where the plaquettes form linear chains (one-dimensional Cu-O<sub>3</sub> network), and Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> which contains CuO<sub>2</sub> planes (two-dimensional Cu-O<sub>2</sub> network).
In the case of zero-dimensional edge-shared or separated plaquettes (like Bi<sub>2</sub>CuO<sub>4</sub>) the main line consists of a single feature which is explained by a local process in which the valence hole moves from the core-hole site to the surrounding O sites. For one-dimensional zigzag or linear chains (like Sr<sub>2</sub>CuO<sub>3</sub>) two features contribute to the main line. In the case of two-dimensional planar systems (like Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>) the main line is composed of at least three features. In addition to these effects there is an interesting trend in the intensity ratio $`I_s/I_m`$ between the satellite and the main line. For zero-dimensional systems this ratio is rather large, e.g. $`I_s/I_m0.58`$ for Bi<sub>2</sub>CuO<sub>4</sub>. On the other hand, in the case of one-dimensional networks the ratio is small, for instance $`I_s/I_m=0.37`$ for Sr<sub>2</sub>CuO<sub>3</sub>. The ratio for two-dimensional systems lies in-between these values, e.g. $`I_s/I_m=0.52`$ for Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>.
The paper is organized as follows. In Sec. II we introduce the model Hamiltonian and describe the method of calculation. Section III contains a discussion of the results. Our conclusions are summarized in Sec. IV.
## II Model and calculation
For the calculation of the Cu 2p<sub>3/2</sub> photoemission spectra we use a multi-band Hubbard Hamiltonian to describe the Cu-O network, while additional terms represent the interaction between the core hole and valence holes on the core-hole site. The full Hamiltonian in the hole picture is
$`H`$ $`=`$ $`\mathrm{\Delta }{\displaystyle \underset{j\sigma }{}}n_{j\sigma }^p+U_{dd}{\displaystyle \underset{i}{}}n_i^dn_i^d+U_{dc}{\displaystyle \underset{\sigma \xi }{}}n_{0\sigma }^dn_{0\xi }^c`$ (3)
$`+I_{dc}\text{S}_0^d\text{J}_0^c+t_{pd}{\displaystyle \underset{ij\sigma }{}}\varphi _{pd}^{ij}(p_{j\sigma }^{}d_{i\sigma }+h.c.)`$
$`+t_{pp}{\displaystyle \underset{jj^{}\sigma }{}}\varphi _{pp}^{jj^{}}p_{j\sigma }^{}p_{j^{}\sigma }\text{,}`$
where $`d_{i\sigma }^{}`$ ($`p_{j\sigma }^{}`$) create a hole with spin $`\sigma `$ in the $`i`$-th Cu $`3d`$ orbital ($`j`$-th O $`2p`$ orbital) and $`n_{i\sigma }^d`$ ($`n_{j\sigma }^p`$) are the corresponding occupation-number operators. The first and second term on the r.h.s. of Eq. (3) describe the charge-transfer energy $`\mathrm{\Delta }`$ and the on-site Coulomb repulsion $`U_{dd}`$ between Cu $`3d`$ valence holes. The third and fourth term represent the local Coulomb repulsion $`U_{dc}`$ and the effective exchange interaction $`I_{dc}`$ between the $`3d`$ valence holes and the $`2p_{3/2}`$ core hole at the core-hole site (which is taken to be site $`i=0`$). $`\text{S}_0^d`$ is the spin-$`1/2`$ operator of a Cu $`3d`$ hole, while $`\text{J}_0^c`$ and $`n_{0\xi }^c`$ are the pseudo-spin $`3/2`$ operator and the number operator of the $`2p_{3/2}`$ core hole (with $`\xi =\pm 3/2,\pm 1/2`$) at site $`i=0`$. Finally, the last two terms on the r.h.s. of Eq. (3) describe the hybridization of Cu $`3d`$ and O $`2p`$ orbitals (hopping strength $`t_{pd}`$) and of O $`2p`$ orbitals (hopping strength $`t_{pp}`$). The factors $`\varphi _{pd}^{ij}`$ and $`\varphi _{pp}^{jj^{}}`$ give the correct sign for the hopping processes and $`ij`$ denotes the summation over nearest neighbor pairs. Hamiltonian (3) describes delocalization processes and multiplet splitting within one framework.
The spectral intensity $`I\left(\omega \right)`$, as a function of binding energy $`\omega `$, is obtained from the hole-hole correlation function
$`I\left(\omega \right)`$ $`=`$ $`{\displaystyle \underset{\xi }{}}\mathrm{}\left[G_{00}^\xi \left(\omega +i0\right)\right]\text{,}`$ (4)
$`G_{00}^\xi \left(\omega +i0\right)`$ $`=`$ $`\mathrm{\Psi }|c_{0\xi }{\displaystyle \frac{1}{\omega +i0}}c_{0\xi }^{}|\mathrm{\Psi }\text{,}`$ (5)
where $`c_{0\xi }^{}`$ creates a core hole with pseudo-spin $`\xi `$ at site $`i=0`$. $``$ is the Liouville operator defined by $`A=[H,A]`$ for any operator $`A`$. $`|\mathrm{\Psi }`$ is the full ground state of $`H`$ before the core hole is created in the photoemission process.
As illustrated in Ref. , the full ground state $`|\mathrm{\Psi }`$ can be constructed by starting from the atomic, Néel-ordered ground state $`|\psi _0`$. In $`|\psi _0`$ all O sites are empty and every Cu site is singly occupied with alternating spin direction. Using an exponential transformation, $`|\psi _0`$ is approximately transformed into the full ground state $`|\mathrm{\Psi }`$
$$|\mathrm{\Psi }=\mathrm{exp}\left(\underset{i\alpha }{}\lambda _\alpha F_{i,\alpha }\right)|\mathrm{\Psi }_N\text{,}$$
(6)
with fluctuation operators $`F_{i,\alpha }`$ and fluctuation strengths $`\lambda _\alpha `$. The operators $`F_{i,\alpha }`$ describe delocalizations of a valence hole which was initially located at Cu site $`i`$. For instance, fluctuation $`F_{i,2d}`$ describes the creation of doubly occupied Cu sites due to the hopping of a hole from Cu site $`i`$, via nearest neighbor O site $`j`$, to the nearest neighbor Cu sites $`k`$
$`F_{i,2d}={\displaystyle \underset{jk\sigma }{}}n_{k\overline{\sigma }}^dd_{k\sigma }^{}\left(1n_{j\sigma }^p\right)d_{i\sigma }\text{.}`$
In the case of a single CuO<sub>4</sub> plaquette one fluctuation operator is sufficient to obtain the exact ground state. For a CuO<sub>3</sub> chain we use up to nine fluctuation operators thus allowing for delocalizations leading as far as to the next-nearest neighbor plaquette. Due to its higher symmetry, the ground state of a CuO<sub>2</sub> plane is well described by only five fluctuation operators. For a detailed description of the fluctuation operators $`F_{i,\alpha }`$ see Ref. . The fluctuation strengths $`\lambda _\alpha `$ are determined using the set of equations
$$0=\mathrm{\Psi }\left|[H,F_\alpha ^{}]\right|\mathrm{\Psi }\text{ , }\alpha =1,2,\mathrm{}\text{.}$$
(7)
which follows from the condition that $`|\mathrm{\Psi }`$ is an eigenstate of the full Hamiltonian (3). The parameters $`\lambda _\alpha `$ are found to decrease exponentially with increasing length of the fluctuation processes. This result is a retrospective justification for the neglect of fluctuations which lead beyond the range covered by the operators $`F_{i,\alpha }`$. Ground state (6) has charge properties which are size consistent and agree well with the results of Quantum Monte Carlo simulations.
From Eq.(4) the Cu $`2p_{3/2}`$ photoemission spectra are calculated using the Mori-Zwanzig projection technique. This method uses a set of operators $`D_\mu `$, the so-called dynamical variables. For these dynamical variables the following matrix equation holds
$$\underset{\gamma }{}\left(z\delta _{\mu \gamma }\omega _{\mu \gamma }\mathrm{\Sigma }_{\mu \gamma }\left(z\right)\right)G_{\gamma \nu }\left(z\right)=\chi _{\mu \nu }\text{,}$$
(8)
where $`z=\omega +i0`$, and where $`\delta _{\mu \gamma }`$ is the unity matrix. The correlation functions $`G_{\gamma \nu }\left(z\right)`$ are given by
$$G_{\gamma \nu }\left(z\right)=\mathrm{\Psi }\left|D_\gamma ^{}\frac{1}{z}D_\nu \right|\mathrm{\Psi }\text{.}$$
(9)
In Eq.(8) the susceptibility matrix $`\chi _{\mu \nu }`$, the frequency matrix $`\omega _{\mu \gamma }`$, and the self-energy matrix $`\mathrm{\Sigma }_{\mu \gamma }`$ are defined by
$`\chi _{\mu \nu }`$ $`=`$ $`\mathrm{\Psi }\left|D_\mu ^{}D_\nu \right|\mathrm{\Psi }\text{,}`$ (10)
$`\omega _{\mu \gamma }`$ $`=`$ $`{\displaystyle \underset{\eta }{}}\mathrm{\Psi }\left|D_\mu ^{}D_\eta \right|\mathrm{\Psi }\chi _{\eta \gamma }^1\text{,}`$ (11)
$`\mathrm{\Sigma }_{\mu \gamma }\left(z\right)`$ $`=`$ $`{\displaystyle \underset{\eta }{}}\mathrm{\Psi }\left|D_\mu ^{}Q{\displaystyle \frac{1}{zQQ}}QD_\eta \right|\mathrm{\Psi }\chi _{\eta \gamma }^1\text{.}`$ (12)
$`Q`$ is a projector on the subspace which is orthogonal to the space spanned by the dynamical variables $`D_\mu `$. It is defined by
$$Q=1\underset{\mu \nu }{}D_\mu |\mathrm{\Psi }\chi _{\mu \nu }^1\mathrm{\Psi }|D_\nu ^{}\text{.}$$
(13)
The set $`\left\{D_\mu \right\}`$ should contain the dynamical variable $`D_{0\xi }=c_{0\xi }^{}`$. Then, one of the correlation functions in Eq.(9) is the hole-hole correlation function $`G_{00}^\xi `$ of Eq.(5). In this case, solving Eq.(8) for $`G_{00}^\xi `$ one obtains the spectral intensity $`I\left(\omega \right)`$ from Eq.(4). One possible approach to obtain an approximate solution of Eq.(8) is to make the set $`\left\{D_\mu \right\}`$ of dynamical variables sufficiently large so that the self-energies $`\mathrm{\Sigma }_{\mu \gamma }`$ can be neglected. In this case an approximation for $`I\left(\omega \right)`$ results which can be systematically improved by further enlarging the set of dynamical variables until the results are converged. Due to the neglect of the self-energy matrix the calculated spectra show no broadening. Thus, a convolution with an artificial line width $`\mathrm{\Gamma }`$ is necessary for a comparison with the experiment.
Besides $`D_{0\xi }`$ the chosen set of dynamical variables includes the operator
$`D_{0\xi }^{^{}}={\displaystyle \underset{\sigma }{}}d_{0,\sigma }^{}d_{0\sigma }c_{0,\xi +2\sigma }^{}\text{,}`$
which describes the creation of a Cu core hole on site $`i=0`$ and a spin flip of the valence hole on the same site due to the exchange interaction with the core hole. Furthermore, we include the following dynamical variables
$`D_{\alpha \xi }`$ $`=`$ $`F_{0,\alpha }c_{0\xi }^{}\text{,}`$ (14)
$`D_{\alpha \xi }^{^{}}`$ $`=`$ $`F_{0,\alpha }{\displaystyle \underset{\sigma }{}}d_{0,\sigma }^{}d_{0\sigma }c_{0,\xi +2\sigma }^{}\text{,}`$ (15)
for all fluctuation operators $`F_{0,\alpha }`$ used in the ground-state calculation. These dynamical variables describe the creation of a core hole and the subsequent delocalization of the valence hole from the same site, without $`\left(D_{\alpha \xi }\right)`$ and with spin flip $`\left(D_{\alpha \xi }^{^{}}\right)`$. In the case of a single CuO<sub>4</sub> plaquette a set of $`14`$ dynamical variables suffices for the exact solution. For a CuO<sub>2</sub> plane $`33`$ variables lead to a well-converged spectrum, while $`40`$ variables are sufficient for a CuO<sub>3</sub> chain.
In Fig.1 the convergence of the results is exemplified for the case of a CuO<sub>3</sub> chain. Figure 1(a) shows the spectrum obtained with $`40`$ dynamical variables while for the spectrum in Fig. 1(b) the set of dynamical variables is almost twice as large ($`70`$ variables). The inclusion of the additional variables leads to a small redistribution of spectral weight around $`935`$ eV binding energy. This change is only visible when a small line width $`\mathrm{\Gamma }`$ is used for the convolution of the line spectrum.
## III Results
A typical set of values for the parameters in model (3) has been obtained for La<sub>2</sub>CuO<sub>4</sub> by band-structure calculation
$`\mathrm{\Delta }`$ $`=`$ $`3.5\text{ eV, }U_d=8.8\text{ eV,}`$ (16)
$`t_{pd}`$ $`=`$ $`1.3\text{ eV, }t_{pp}=0.65\text{ eV.}`$ (17)
For the comparison with the experimental spectra we start with this set and adapt some of the parameters to the specific systems which we investigate. The value of the exchange parameter $`I_{dc}=1.5`$ eV and the line width $`\mathrm{\Gamma }=1.8\text{eV}`$ of the Gaussian function to convolute the spectrum are obtained from a comparison of the exact solution for a single CuO<sub>4</sub> plaquette with the experimental result for Bi<sub>2</sub>CuO<sub>4</sub>, see Fig. 2. From this comparison we furthermore determine the absolute energetical position of all calculated spectra. Thereby we allow for general shifts of $`\pm 0.3\text{eV}`$ which is the experimental accuracy of the absolute energy values.
To fit the experimental spectra, we vary $`\mathrm{\Delta }`$ and $`U_{dc}`$ until the calculated satellite to main-peak intensity ratio $`I_s/I_m`$ coincides with the experimental one. It is not possible to use only $`\mathrm{\Delta }`$ as fit parameter because the separation between the satellite and the main line is mainly determined by the difference $`U_{dc}\mathrm{\Delta }`$. However, we can reduce the number of free parameters by keeping this difference constant.
For the spectra of Bi<sub>2</sub>CuO<sub>4</sub> and Sr<sub>2</sub>CuO<sub>3</sub> we keep $`U_{dc}\mathrm{\Delta }=5\text{eV}`$, while the smaller satellite-main line separation in the spectrum of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> is taken into account by keeping $`U_{dc}\mathrm{\Delta }=4.2\text{eV}`$. In this way, we obtain a good fit of the Bi<sub>2</sub>CuO<sub>4</sub> spectrum with $`\mathrm{\Delta }=3.5\text{eV}`$, see Fig. 2, while the spectrum of Sr<sub>2</sub>CuO<sub>3</sub> is well described by the result of the projection technique for a CuO<sub>3</sub> chain with $`\mathrm{\Delta }=2.7\text{eV}`$, see Fig. 3. These results have already been described elsewhere. Here, we therefore discuss only those aspects which are relevant for the influence of dimensionality and lattice geometry.
The lower parts of Figs. 2 and 3 show the delocalization properties of the final states which lead to the most important lines in the calculated spectra. The final state associated with the satellite lines, i.e. state (a) in Figs. 2 and 3, is highly localized. Since the valence hole remains mainly on the core-hole site, the influence of geometry is negligible. This is consistent with the experimental observation that the satellite peaks depend much less on the geometry than the main lines.
Final state (b) which causes the main line in the CuO<sub>4</sub> plaquette spectrum (Fig. 2) leads to a shoulder in the main line of the CuO<sub>3</sub> chain spectrum (Fig. 3) and is also rather localized, with most of the valence hole density concentrated on the O sites surrounding the core-hole site. However, in the case of the CuO<sub>3</sub> chain the hole density at the O sites in chain direction is smaller than at the O sites perpendicular to it. A comparable effect has been observed in exact diagonalization studies. The most important effect of the network geometry is the emergence of peak (c) in the CuO<sub>3</sub> chain spectrum, see Fig. 3. This peak, which dominates the main line, is associated with a delocalization of the valence hole to the neighboring plaquettes which may be interpreted as the formation of a Zhang-Rice singlet.
In the case of the CuO<sub>2</sub> plane it turns out that a fit using $`\mathrm{\Delta }`$ as described above does not lead to a satisfying agreement with the spectrum of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>, see Fig. 4. For the standard value $`\mathrm{\Delta }=3.5\text{eV}`$ (dashed line in Fig. 4) the calculated ratio $`I_s/I_m=0.4`$ is too small compared to the experimental value $`I_s/I_m=0.52`$. With increasing $`\mathrm{\Delta }`$ (i.e. with decreasing charge fluctuations) the relative intensity of both the satellite and the shoulder structure around $`935\text{eV}`$ binding energy increases. For $`\mathrm{\Delta }=4.25\text{eV}`$ (solid line in Fig. 4) the calculated ratio $`I_s/I_m`$ is equal to the experimental one. However, the form of the main line is not reproduced correctly. The intensity of the shoulder structure is overestimated while a line seems to be missing around $`934\text{eV}`$ binding energy.
In order to obtain a better fit we have varied $`t_{pd}`$ and $`U_{dc}`$ instead of $`\mathrm{\Delta }`$ and $`U_{dc}`$. As shown in Fig. 5 this leads to a better agreement with the spectrum of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. Keeping $`\mathrm{\Delta }=3.5\text{eV}`$, the experimental ratio $`I_s/I_m=0.52`$ is reproduced for $`t_{pd}=1.15\text{eV}`$ and $`U_{dc}=8.1\text{eV}`$. The delocalization properties of the most important final states are shown in the lower part of Fig.5. As opposed to the zero- and one-dimensional cases shown in Figs. 2 and 3, the valence hole from the core-hole site may now delocalize in two dimensions. Nevertheless the final states obtained for a CuO<sub>2</sub> plane are rather similar to those of the lower dimensional systems. For state (a) this is easy to explain by the local nature of the satellite peak. The fact that the delocalization in states (b) and (c) is not significantly larger in two dimensions than in one dimension is, on the other hand, somewhat surprising. Key element in the explanation of this effect are the four Cu sites which are the diagonal nearest neighbors of the core-hole site (cf. the lower part of Fig. 5). Due to antiferromagnetic correlations these diagonal Cu sites are predominantly occupied by valence holes which have the same spin direction as the valence hole on the core-hole site. Therefore, due to the Pauli principle the holes from the diagonal Cu sites suppress fluctuations of the valence hole from the core-hole site. The increase of charge fluctuations due to the higher dimensionality is largely compensated by this suppression. An analogous effect has already been found to be important for the ground-state charge properties of Hamiltonian (3).
It is interesting to compare the results of the projection technique with exact diagonalization calculations (without multiplet splitting) of a Cu<sub>5</sub>O<sub>16</sub> cluster from Refs. . This cluster contains five plaquettes in a cross-like configuration where the central Cu site is the core-hole site. Notice that the Cu<sub>5</sub>O<sub>16</sub> system does not contain the diagonal Cu sites which, as discussed above, suppress fluctuations from the central Cu site. Therefore one expects this system to display features of artificially strong delocalization like, e.g., a reduced ratio $`I_s/I_m`$. In fact, while the diagonalization shows a similar shoulder structure as the projection technique, cf. Fig. 4, the intensity of this shoulder, its separation from the lowest-energy line as well as the ratio $`I_s/I_m`$ are smaller.
The slightly reduced value $`t_{pd}=1.15\text{eV}`$ for the Cu-O hopping strength which leads to the spectrum shown in Fig. 5 may be explained by the larger Cu-O distance in Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> compared to La<sub>2</sub>CuO<sub>4</sub> Nevertheless, the agreement of the theoretical result with the experimental spectrum is still not satisfactory. The calculated main line is dominated by two features, and at least one excitation seems to be missing in the region around $`934\text{eV}`$ binding energy. Notice, that in the theoretical main line of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> from Ref. an excitation is missing as well.
These results suggest the conclusion that Hamiltonian (3) still does not include all degrees of freedom which are necessary for a detailed description of the main line in the Cu 2p<sub>3/2</sub> spectrum of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. Since all effects of the Cu-O network dimensionality are already included in model (3) we expect the missing third main-line feature to be a material-specific effect. Orbitals which are not yet taken into account are, for example, non-planar orbitals in the CuO<sub>2</sub> system, like the Cu $`3d_{z^2r^2}`$ orbital. However, in view of the good agreement with the experiments shown in Figs. 2 and 3, these orbitals do not seem necessary for the description of Bi<sub>2</sub>CuO<sub>4</sub> and Sr<sub>2</sub>CuO<sub>3</sub>.
It may also be possible that sites which do not belong to the CuO<sub>2</sub> plane (like the Cl apex site) contribute to the screening in Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. Notice that the Cu $`2p`$ main line in copper oxides without apex-oxygen site (CuO, Nd<sub>2</sub>CuO<sub>4</sub>) is narrower than in compounds with one apex oxygen (Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub>) or two apex oxygens (La<sub>2</sub>CuO<sub>4</sub>) per copper site. However, both the large spatial distance between the Cl and the CuO<sub>2</sub> plane and the large energy difference between the Cl-$`3p`$ and the Cu-$`3d`$-O-$`2p`$ line in the valence photoelectron spectrum suggest that screening from Cl sites should be small. Another possible candidate for the explanation of the missing third main-line feature are non-bonding oxygen $`2p`$ orbitals. Recently, a feature in the optical spectrum of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> has been attributed to these orbitals.
The influence of Cu-O network dimensionality on the Cu $`2p_{3/2}`$ spectra is illustrated in Fig. 6 where we show results of the projection technique for zero-, one- and two-dimensional Cu-O networks. Since the same parameter set (17) has been used for all geometries, all changes in the spectra are exclusively due to dimensionality effects. The most important effect in the spectra is that for the one- and two-dimensional system an additional excitation appears at lower binding energies. As discussed above the final states of this excitation are delocalized and rather similar for both the one- and the two-dimensional structures, see states (c) in Figs. 3 and 5. Overall, there is only a quantitative change from one to two dimensions, in contrast to the qualitative change observed between zero and higher dimensions. This conclusion is in principal agreement with results of exact diagonalization calculations. The peak around $`934\text{eV}`$ binding energy which dominates the main line in the case of zero dimensions becomes a shoulder structure which decreases in intensity and shifts towards higher binding energies as the dimensionality increases. Nevertheless, the final state associated with this peak preserves its main properties (a large valence-hole density at the O sites around the core-hole site) with changing dimensionality, see states (a) in Figs. 2, 3, and 5. As the dimensionality increases the delocalization in the Cu-O network increases. Therefore one observes a monotonic decrease in the ratio $`I_s/I_m`$ for increasing dimensions. Since this trend is not observed experimentally, the actual value of $`I_s/I_m`$ has to depend mainly on material-specific properties. In our calculation these properties are reflected by the values of the model parameters, cf. Figs. 2 to 5.
## IV Summary
Summing up, we have calculated the Cu $`2p_{3/2}`$ core-level spectra of zero- one- and two-dimensional Cu-O networks using a model Hamiltonian which describes exchange splitting and delocalization within one framework. The spectral intensity has been obtained using the Mori-Zwanzig projection technique. This method leads to the exact solution for a single CuO<sub>4</sub> plaquette, and we observe excellent convergence in the case of an infinite CuO<sub>3</sub> chain and an infinite CuO<sub>2</sub> plane. The delocalization properties of the final states obtained in the calculation are easy to interpret. The results have been compared to experimental spectra by using either $`\mathrm{\Delta }`$ or $`t_{pd}`$ as a fit parameter. While there is a good agreement between theory and experiment in the case of Bi<sub>2</sub>CuO<sub>4</sub> and Sr<sub>2</sub>CuO<sub>3</sub>, one excitation seems to be missing when compared with the spectrum of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. An analysis of the influence of dimensionality effects, as compared to effects due to material-specific properties, indicates that this missing feature may be due to orbitals which are not contained in the model. It is also found that dimensionality plays a minor role for the satellite-main line intensity ratio $`I_s/I_m`$ which is mainly determined by the values of the model parameters (especially the value of $`\mathrm{\Delta }`$).
###### Acknowledgements.
Discussions with J. Fink, M. S. Golden, A. Goldoni, R. E. Hetzel, F. Parmigiani, and L. Sangaletti are gratefully acknowledged. This work is supported by DFG through the research program of the SFB 463, Dresden. |
warning/0003/math0003069.html | ar5iv | text | # Untitled Document
Introduction à la Conjecture d’Alexandru
Rappelons quelques résultats de Bernstein, Gelfand, Gelfand, Delorme, Beilinson, Guinzburg et Soergel. Soient g une algèbre de Lie semisimple complexe, b une sous-algèbre de Borel et h une sous-algèbre de Cartan contenue dans b. Soient $`𝒪`$ la catégorie associée à ces données par BGG et $`𝒪_\rho `$ la sous-catégorie pleine de $`𝒪`$ dont les objets ont le caractère infinitésimal généralisé du module trivial. Notons $`\rho `$ la demi-somme des racines positives et $`W`$ le groupe de Weyl, muni de sa fonction longueur $`\mathrm{}`$ et de son ordre de Bruhat. À $`wW`$ attachons le module de Verma $`M_w`$ de plus haut poids $`w\rho \rho `$ ; rappelons que $`M_w`$ a un unique sous-module maximal ; notons $`L_w`$ le quotient correspondant. Soit $`P_w`$ un revètement projectif de $`L_w`$ ; posons $`P:=_wP_w`$, $`A:=(End_\text{g}P)^{op}`$ ; notons $`A`$-df la catégorie des $`A`$-modules de dimension finie et $`E`$ l’équivalence $`Hom_\text{g}(P,)`$ de $`𝒪_\rho `$ sur $`A`$-df. Par abus notons encore $`M_w`$ et $`L_w`$ les images de ces objets par $`E`$, et désignons par $`𝑴_w`$ et $`𝑳_w`$ leurs classes respectives dans le groupe de Grothendieck. Notons $`e_wA`$ la projection sur $`P_w`$.
Théorème 1. On a $`M_wAe_w/{\displaystyle \underset{xw}{}}Ae_xAe_w=Ae_w/{\displaystyle \underset{x>w}{}}Ae_xAe_w.`$
Théorème 2. On a $`End_A(M_w)=`$.
Considérons les polynômes de Delorme $`a_{x,y}:=SPExt_A^{}(M_x,L_y)`$$`SP`$ signifie $``$série de Poincaré$``$.
Théorème 3. On a $`𝑳_y=_xa_{x,y}(1)𝑴_x.`$
Théorème 4. Il existe des polynômes $`P_{x,y}`$ tels que
(1) $`a_{x,y}=t^{\mathrm{}(y)\mathrm{}(x)}P_{x,y}(t^2),`$
(2) $`P_{x,y}0xyP_{x,y}(0)=1,`$
(3) $`P_{x,x}=1,`$
(4) $`degP_{x,y}<{\displaystyle \frac{\mathrm{}(y)\mathrm{}(x)}{2}}\text{ si }x<y`$.
Théorème 5. On a $`SPExt_A^{}(L_x,L_y)=_za_{z,x}a_{z,y}`$.
Voici des analogues conjecturaux des ces énoncés pour les modules de Harish-Chandra.
Soient $`G`$ un groupe de Lie semi-simple connexe à centre fini, $`K`$ un sous-groupe compact maximal, $``$ la catégorie des modules de Harish-Chandra associée à ces données et $`_\rho `$ la sous-catégorie pleine de $``$ dont les objets ont le caractère infinitésimal généralisé du module trivial. Notons $`r`$ le rang (réel) de $`G`$ et $`Z`$ l’algèbre $`[[z_1,\mathrm{},z_r]]`$. Soit $`A`$ une $`Z`$-algèbre telle que $`_\rho A`$-df, $`A`$ est de type fini sur $`Z`$, $`A`$ est commutative modulo son radical $`R`$ et $`A`$ est $`R`$-adiquement complète. (De telles algèbres existent et sont isomorphes en tant que $``$-algèbres.) Choisissons une sous-algèbre $`A_0`$ de $`A`$ relevant $`A/R`$ et notons $`\{e_i|iI\}`$ l’ensemble (fini) des idempotents minimaux de $`A_0`$. Soit $`L_i`$ le $`A`$-module simple associé à $`iI`$, soit $`\mathrm{}(i)`$ la dimension projective de $`L_i`$ et $``$ le plus petit ordre sur $`I`$ satisfaisant $`ij`$ chaque fois que
$$\mathrm{}(j)=\mathrm{}(i)+1\text{et}Ext_A^1(L_j,L_i)0.$$
Utilisons librement les analogues évidents des notations introduites dans le cadre de la catégorie $`𝒪`$.
Conjecture 1. On a $`Ae_i/_{ji}Ae_jAe_i=Ae_i/_{j>i}Ae_jAe_i.`$
Notons ce module $`M_i`$ et posons
$$\overline{M}_i:=M_i/\text{rad}(End_AM_i)M_i.$$
Cet objet ne coïncide pas toujours avec le module de Langlands correspondant.
Conjecture 2. On a $`End_A(\overline{M}_i)=`$.
Considérons les polynômes de Delorme $`a_{ij}:=SPExt_A^{}(M_i,L_j)`$.
Conjecture 3. On a $`𝑳_j=_ia_{ij}(1)\overline{𝑴}_i.`$
Conjecture 4. Il existe des polynômes $`p_{ij}`$ satisfaisant (1), …, (4).
Conjecture 5. Il existe des polynômes $`d_k`$ tels que
$$SPExt_A^{}(L_i,L_j)=\underset{k}{}d_ka_{ki}a_{kj}.$$
Le principal inconvénient de cette approche des modules de Harish-Chandra est que, contrairement à ce qui se passe pour les modules de BGG, rien de tout cela n’est calculable ! Voici un remède à la fois partiel et conjectural à ce mal. Supposons que $`G`$ et $`K`$ ont même rang. Dans la classification de Langlands $`L_i`$ apparaît comme l’unique quotient simple d’un module induit à partir d’un sous-groupe parabolique $`P_i`$ ; soit $`\text{p}_i=\text{m}_i\text{a}_i\text{n}_i`$ la décomposition de Langlands de Lie$`(P_i)`$ ; posons
$$\stackrel{~}{d}_i:=\left(1t^2\right)^{dim\text{a}_i};$$
soit $`\stackrel{~}{\mathrm{}}(i)`$ la dimension de la $`K_{}`$-orbite attachée à $`i`$ et $`(\stackrel{~}{p}_{ij})`$ la famille des polynômes de Kazhdan-Lusztig-Vogan ; posons
$$\stackrel{~}{a}_{ij}(t)=t^{\stackrel{~}{\mathrm{}}(j)\stackrel{~}{\mathrm{}}(i)}\stackrel{~}{p}_{ij}(t^2).$$
Conjecture 5’. On a $`SPExt_A^{}(L_i,L_j)=_k\stackrel{~}{d}_k\stackrel{~}{a}_{ki}\stackrel{~}{a}_{kj}`$.
This text and others are available at http://www.iecn.u-nancy.fr/$``$gaillard
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warning/0003/hep-th0003234.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Recent works on theories of open strings and D$`p`$-branes with a constant nonvanishing Neveu-Schwarz 2-form $`B_{ij}`$ have suggested that the noncommutativity which appears is an underlying and very general property of such theories . Since Hopf algebras (HAs) often lie at the root of noncommutative systems, we were motivated to look for a HA structure for these theories, and showed that the noncommutative $``$-product was in fact a specific case of a more general multiplication defined in terms of the R-matrix $`R`$ of a quasitriangular HA $``$; furthermore, when $`=`$, where $``$ was the HA of functions on $`^{p+1}`$, as was the case for the aforementioned noncommutative string theories, we found an explicit form for $`R`$ which covers both the commutative ($`B_{ij}=0`$) and noncommutative ($`B_{ij}0`$) cases.
However, it was not immediately apparent how we could introduce derivations on the algebra endowed with this multiplication, $`\widehat{}`$. One way to think of derivations on a HA $``$ is as elements of the dually paired HA $`^{}`$, with the action of the latter on the former given in terms of the HA properties of both. The problem was that $`\widehat{}`$ was shown not to be a HA, and therefore neither was the dually paired space $`\widehat{}^{}`$. This precluded the interpretation of the latter as local derivatives on $`\widehat{}`$, so it was not immediately obvious how one might define a gauge theory on the noncommutative space, since we needed a derivative in order to construct the (noncommutative) field strength tensor $`\widehat{F}_{ij}`$ from the gauge field $`\widehat{A}_i`$, i.e$`\widehat{F}_{ij}=_{[i}\widehat{A}_{j]}i\widehat{A}_{[i}\widehat{A}_{j]}`$ . The question was, what could we use for $`_i`$? We speculated that we might have to replace local derivatives by difference operators, but this guess seemed to be contradicted by the fact that regular derivatives were used consistently in .
In this follow-up note to , we explain why the usual derivatives are in fact the correct ones when dealing with noncommutative string theory: We show that even though $`\widehat{}`$ is not a HA, there exists a HA which has a well-defined action on $`\widehat{}`$ and plays the role of the space of derivations. This holds for arbitrary $``$, and for the specific case where $`=`$, this HA is $`^{}`$ and the action is the same as that for the usual partial derivative. This is done in Section 2.
The ability to relate the commutative and noncommutative theories via the R-matrix, however, turns out to be a bit of a fluke, being true only if $`^{}`$ is cocommutative. While this is certainly true of the space of derivations $`^{}`$, if we want to be as general as possible, we must relax this condition. In Section 3, we demonstrate how this can be done by using the Drinfel’d twist , which allows us to find a generalisation of the $``$-product and the space of derivations on the algebra constructed with this $``$. (Related but more mathematical treatments of this construction may be found in , and very recently , which covers much of the same in a broader context.)
However, using the Drinfel’d twist gives exactly the same derivatives and noncommutative product as if we had used an R-matrix approach, so why pick one over the other? In Section 4, we present two arguments why we think the former is more appropriate: First, the R-matrix construction connects the commutative and noncommutative cases only when $``$ is commutative (i.e$`^{}`$ is cocommutative), whereas the Drinfel’d twist includes both cases, and therefore does not require us to make any a priori assumptions about the algebraic structure of $``$. Secondly, an R-matrix must satisfy two coproduct conditions, while the analogous element in the Drinfel’d twisting only has to fulfill one, and there may be a way of naturally implementing the latter using the Ward identity which must arise out of gauge-fixing the form of $`B_{ij}`$. These reasons lead us to think that the Drinfel’d twist plays a fundamental role in noncommutative string theory, specifically in helping to determine the form of the low-energy effective action of the theory in terms of the commutative coordinates, e.g. Born-Infeld.
Throughout this letter we use terms and notations described in our previous paper ; the reader is referred therein for the details.
## 2 The Leibniz Rule and the $``$-product
We begin with a HA $``$ and its dually paired HA $`^{}`$. The (left) action of $`x^{}`$ on $`f`$ is given by $`xf:=f_{(1)}x,f_{(2)}`$, and satisfies the Leibniz rule $`x(fg)=\left(x_{(1)}f\right)\left(x_{(2)}g\right)`$. The elements of $`^{}`$ (with the exception of the unit 1 and its multiples) thus may be thought of as derivations on $``$.
In the case where $`^{}`$ is quasitriangular with R-matrix $`R`$, we can define a new multiplication between $`f,g`$ as
$$fg:=f_{(1)}g_{(1)}R_{21},f_{(2)}g_{(2)}.$$
(2.1)
$`\widehat{}`$ is taken to be the algebra equivalent to $``$ as a vector space and with the multiplication $``$. It is not a HA, so neither is the dually paired coalgebra $`\widehat{}^{}`$. We therefore cannot think of $`\widehat{}^{}`$ as derivations of $`\widehat{}`$.
However, let’s go ahead and compute the action of $`x^{}`$ on the product $`fg`$:
$`x(fg)`$ $`=`$ $`x\left(f_{(1)}g_{(1)}\right)R_{21},f_{(2)}g_{(2)}`$ (2.2)
$`=`$ $`f_{(1)}g_{(1)}x,f_{(2)}g_{(2)}R_{21},f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(1)}\mathrm{\Delta }(x)R_{21},f_{(2)}g_{(2)}f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(1)}\mathrm{\Delta }(x)R_{21},f_{(2)}g_{(2)}.`$
Using $`\mathrm{\Delta }(x)=R_{21}\left(\tau \mathrm{\Delta }(x)\right)R_{21}^1`$, we obtain
$`x(fg)`$ $`=`$ $`f_{(1)}g_{(1)}R_{21}\left(\tau \mathrm{\Delta }\right)(x),f_{(2)}g_{(2)}`$ (2.3)
$`=`$ $`f_{(1)}g_{(1)}R_{21}\left(\tau \mathrm{\Delta }\right)(x),f_{(2)}g_{(2)}f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(2)}R_{21},f_{(2)}g_{(2)}x_{(2)}x_{(1)},f_{(3)}g_{(3)}`$
$`=`$ $`\left(f_{(1)}g_{(2)}\right)x_{(2)},f_{(2)}x_{(1)},g_{(2)}`$
$`=`$ $`\left(x_{(2)}f\right)\left(x_{(1)}g\right).`$
So we see that the Leibniz rule is ‘reversed’: The first piece of the coproduct of $`x`$ acts on the second function, and vice versa.
In , we reviewed the construction of the HA dually paired to $``$, denoted $`^{}`$. Recall that the coproduct and antipode for $`x^{}`$ were defined by
$`\mathrm{\Delta }(x),fg`$ $`:=`$ $`x,fg,`$
$`S(x),f`$ $`:=`$ $`x,S(f).`$ (2.4)
However, this is not the only HA which may be constructed to the vector space dual to $``$: A different HA, called the opposite dual and denoted $`^{\mathrm{op}}`$This is the same well-known HA which plays a key role in the construction of the Drinfel’d double ., can be defined by keeping all the relations between $``$ and $`^{}`$ except the above two, which are replaced by
$`\mathrm{\Delta }^{}(x),fg`$ $`:=`$ $`x,gf,`$
$`S^{}(x),f`$ $`:=`$ $`x,S^1(f).`$ (2.5)
We see that the coproduct on $`^{\mathrm{op}}`$ is the one on $`^{}`$ with the two spaces flipped, i.e$`\mathrm{\Delta }^{}=\tau \mathrm{\Delta }`$. From (2.3) we can see that although $`^{}`$ does not have a action on $`\widehat{}`$, $`^{\mathrm{op}}`$ does, because
$$x(fg)=\left(x_{(1)^{}}f\right)\left(x_{(2)^{}}g\right),$$
(2.6)
where we have used the notation $`\mathrm{\Delta }^{}(x):=x_{(1)^{}}x_{(2)^{}}`$. Hence, we conclude that $`^{\mathrm{op}}`$, not $`^{}`$ or $`\widehat{}^{}`$, is the space of derivations on $`\widehat{}`$. And since in general $``$ is not cocommutative, the Leibniz rule on $`\widehat{}`$ does not have the same form as that on $``$. However, if $`\mathrm{\Delta }^{}=\mathrm{\Delta }`$, i.e$`^{}`$ is cocommutative, then $`^{\mathrm{op}}`$ and $`^{}`$ are the same, and the space of derivations is the same for $`\widehat{}`$ as for $``$.
If we now look at noncommutative string theory, the algebra $``$ is the function algebra $``$ spanned by monomials in the coordinate maps $`x^i`$ taking the D$`p`$-brane into $`^{p+1}`$. The $``$-product is introduced by using the R-matrix
$$R:=e^{\frac{i}{2}\theta ^{ij}_i_j},$$
(2.7)
where $`\theta ^{ij}`$ is related to $`B_{ij}`$ and the open string metric $`g_{ij}`$ by
$$\theta ^{ij}:=\left(2\pi \alpha ^{}\right)^2\left(\frac{1}{g+2\pi \alpha ^{}B}B\frac{1}{g2\pi \alpha ^{}B}\right)^{ij}.$$
(2.8)
This $`R`$ gives the (noncommutative) product between functions $`f`$ and $`g`$ as
$$f(x)g(x)=e^{\frac{i}{2}\theta ^{ij}\frac{^2}{\xi ^i\zeta ^j}}f(x+\xi )g(x+\zeta )|_{\xi =\zeta =0}.$$
(2.9)
The dually paired HA $`^{}`$ is the space spanned by monomials of the partial derivatives $`_i`$. The coproduct is generated by $`\mathrm{\Delta }\left(_i\right)=_i1+1_i`$, and with the action on $``$ being the usual derivative, the Leibniz rule is the familiar
$$_i\left(f(x)g(x)\right)=\left(_if(x)\right)g(x)+f(x)\left(_ig(x)\right)$$
(2.10)
(the $``$ signifying the action has been suppressed in the above two equations). But since this HA is cocommutative, this is also the Leibniz rule for the action of $`^{\mathrm{op}}`$ on $``$. Hence, $`^{}=^{\mathrm{op}}`$, and this is the reason that one can use the familar derivatives even when the space is noncommutative, as was done in .
## 3 The Drinfel’d Twist
The material in the preceding Section is in fact a specific example of a more general construction: Suppose $``$ is a HA such that there exists an invertible element $`F`$ which satisfies $`(ϵid)(F)=(idϵ)(F)=1`$ as well as the coproduct identity
$$F_{12}\left(\mathrm{\Delta }id\right)(F)=F_{23}\left(id\mathrm{\Delta }\right)(F).$$
(3.1)
If this is the case, then a new HA $`^F`$, called the Drinfel’d twist of $``$ , can be defined in the following way: $`^F=`$ at the algebra level, and the counit and unit of $`^F`$ are the same as those of $``$. The coproduct and antipode, however, are given in terms of those on $``$ by
$`\mathrm{\Delta }^F(f):=F\mathrm{\Delta }(f)F^1,`$ $`S^F(f):=\sigma ^1S(f)\sigma ,`$ (3.2)
where $`\sigma `$ is the quantity constructed from $`F:=F_\alpha F^\alpha `$ (sum implied) via
$$\sigma :=m\left((idS)\right)(F)F_\alpha S\left(F^\alpha \right).$$
(3.3)
(The inverse can be shown to be $`\sigma ^1=m\left((Sid)\left(F^1\right)\right)`$.) For future reference, we also use the notation $`\mathrm{\Delta }^F(f):=f_{(1)F}f_{(2)F}`$.
Now suppose we start with dually paired HAs $``$ and $`^{}`$, and an element $`F`$ in $`^{}^{}`$ satisfying (3.1) exists; then a new product, $``$, may be defined on $``$ via
$$fg:=f_{(1)}g_{(1)}F^1,f_{(2)}g_{(2)}.$$
(3.4)
We can then check associativity by first computing the triple product $`(fg)h`$:
$`(fg)h`$ $`=`$ $`\left(f_{(1)}g_{(1)}\right)hF^1,f_{(2)}g_{(2)}`$ (3.5)
$`=`$ $`f_{(1)}g_{(1)}h_{(1)}F^1,f_{(2)}g_{(2)}h_{(2)}F^1,f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(1)}h_{(1)}(\mathrm{\Delta }id)\left(F^1\right),f_{(2)}g_{(2)}h_{(2)}F^1,f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(1)}h_{(1)}(\mathrm{\Delta }id)\left(F^1\right)F_{12}^1,f_{(2)}g_{(2)}h_{(2)}.`$
Computing $`f(gh)`$ in a similar fashion replaces the left argument of the inner product above with $`(id\mathrm{\Delta })\left(F^1\right)F_{23}^1`$, and if we take the inverse of (3.1), we see the two are equal, and this proves that $``$ is associative. The counit condition of $`F`$ ensures that $`1`$ is also the $``$-multiplicative identity as well. We therefore denote by $`\widehat{}`$ the unital associative algebra with vector space $``$ and multiplication $``$.
One consequence of this definition of $``$ is that the Drinfel’d twist of $`^{}`$ is a HA of left actions on $`\widehat{}`$: Taking $`x^{}`$ and $`f,g\widehat{}`$,
$`x(fg)`$ $`=`$ $`x\left(f_{(1)}g_{(1)}\right)F^1,f_{(2)}g_{(2)}`$ (3.6)
$`=`$ $`f_{(1)}g_{(1)}x,f_{(2)}g_{(2)}F^1,f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(1)}\mathrm{\Delta }(x)F^1,f_{(2)}g_{(2)}f_{(3)}g_{(3)}`$
$`=`$ $`f_{(1)}g_{(1)}\mathrm{\Delta }(x)F^1,f_{(2)}g_{(2)}`$
$`=`$ $`f_{(1)}g_{(1)}F^1\mathrm{\Delta }^F(x),f_{(2)}g_{(2)}`$
$`=`$ $`f_{(1)}g_{(1)}x_{(1)F}x_{(2)F},f_{(2)}g_{(2)}`$
$`=`$ $`\left(x_{(1)F}f\right)\left(x_{(2)F}g\right).`$
So $`^F`$ has a well-defined action on $`\widehat{}`$, and therefore may be used as a the space of derivations on $`\widehat{}`$.
What if $`^{}`$ is quasitriangular? Then we automatically have an $`F`$ which satisfies (3.1),namely, $`F=R_{21}^1`$. This follows from the coproduct properties of the R-matrix, and we recover all the results in Section 2: We see immediately that the $``$-product given by plugging $`R_{21}`$ in for $`F^1`$ in (3.4) is the same as (2.1). The coproduct is also the same, $`\mathrm{\Delta }^F=\tau \mathrm{\Delta }=\mathrm{\Delta }^{}`$. To compare the antipodes, note that $`\sigma ^1`$ becomes $`m(Sid)\left(R_{21}\right)`$, which is the element, usually denoted $`u`$, which generates the square of the antipode in $`^{}`$: $`uxu^1=S^2(x)`$ . This immediately leads to $`S^F=S^1=S^{}`$, exactly as expected, and we see that this choice of $`F`$ gives $`\left(^{}\right)^F=^{\mathrm{op}}`$, which, as we proved, is the correct choice for the space of derivations on $`\widehat{}`$.
There is potentially a wider class of $`F`$s than there are of R-matrices, because the one coproduct condition on $`F`$ (3.1) is less restrictive than the two coproduct conditions on $`R`$:
$`(\mathrm{\Delta }id)(R)=R_{13}R_{12},`$ $`(id\mathrm{\Delta })(R)=R_{13}R_{23}.`$ (3.7)
This means that, in principle, there may be other associative $``$-products besides the one defined using the R-matrix. However, for the noncommutative string case, the cocommutativity of $`^{}`$ is still a strong enough condition to force $`F`$ to be identical to the $`R`$ given in (2.7), so in this instance, reformulating the $``$-product on $``$ in terms of the Drinfel’d twist does not change the fact that it has the unique form (2.9). Therefore, using the Drinfel’d twist gives exactly the same noncommutativity to the D$`p`$-brane-open string system as using an R-matrix does.
## 4 Noncommutative String Theory
We conclude this work with some brief speculations about how the Drinfel’d twist may appear in the context of a noncommutative string theory.
In , we conjectured that the dependence on $`\theta ^{ij}`$ in the effective field theory would appear only through the $``$-product, and that since this in turn was given in terms of the R-matrix $`R`$, then there would be explicit dependence on $`R`$ in the action when written in terms of the commutative theory. However, this is in fact probably not the case in general, for the following reason: The $``$-product becomes the commutative product when $`\theta ^{ij}=0`$, which, using the explicit form (2.7), corresponds to $`R=11`$. When we are dealing with a cocommutative HA where $`\mathrm{\Delta }^{}=\mathrm{\Delta }`$, this is an admissible R-matrix, but for the most general case where the HA may not be cocommutative, $`11`$ doesn’t work as an R-matrix.
However, the Drinfel’d twist construction is still applicable, because $`F=11`$ satisfies (3.1), and just gives the trivial case $`\left(^{}\right)^F=^{}`$. Thus, if $`\theta `$ is some element of a parameter space, and there exists a continuous map $`\theta F(\theta )`$ which satisfies (3.1) and $`F(0)=11`$, we have a family of spaces $`\widehat{}(\theta )`$ and Drinfel’d twists $`^F(\theta )`$ continuously connected to the undeformed cases $``$ and $`^{}`$, respectively, with elements of the latter being derivations on elements of the former. This deformation does not depend on the cocommutativity, or lack thereof, of $`^{}`$. Furthermore, since the coproduct of $`^{}`$ is dual to the multiplication on $``$, this also implies that we do not even have to start with a commutative $``$ for this procedure to be valid.
When we look at the specific case of a D$`p`$-brane/open string system, where we expect to be able to go continuously from the commutative case with vanishing $`B_{ij}`$ to the noncommuting theory, it therefore seems reasonable to us that the $`\theta `$-dependence in the effective action when expressed as an integral over the commutative space will be entirely through an $`F`$ and not an $`R`$, and that the underlying structure is that of a Drinfel’d twisted HA rather than a quasitriangular one.
There is another reason to favour the Drinfel’d twist: Recall that the recent work on noncommutative string theory has been done with a constant $`B_{ij}`$. Since the full theory should be invariant under the gauge transformation $`B_{MN}B_{MN}+_{[M}\lambda _{N]}`$ for any $`\lambda _M`$, where $`M,N=0,\mathrm{},9`$, taking $`B_{MN}=0`$ for $`M,N=p+1,\mathrm{},9`$ and constant for $`M,N=0,\mathrm{},p`$ is a gauge fixing condition. We should therefore expect to find a Ward identity resulting from this fixing. To us, it seems very likely that this Ward identity is related to (3.1). This explanation is attractive because, if correct, it means we do not have to impose (3.1) by hand; it comes out naturally from taking $`B_{ij}`$ to be constant. And just as a Ward identity must hold to have a self-consistent theory, i.e. gauge invariance, so must (3.1) hold for consistency, i.e$``$ is associative. For an R-matrix to be involved instead of $`F`$, we would have to come up with some way of arriving at the two conditions (3.7), and the single requirement that the theory have $`B`$-gauge invariance would presumably not give these.
Thus, the signs point more toward a Drinfel’d twisted rather than a quasitriangular HA; more specifically, if we think of the ‘undeformed’ theory as one formulated with the HA $``$ (commutative or not), and the ‘deformed’ one as that on the algebra $`\widehat{}`$, then the former should have explicit $`F`$-dependence. This fact may therefore give some clues as to the explicit form of the undeformed low-energy effective action, even though it is presumably very complicated (unlike the deformed version, which may be very nice, e.g. super-Yang-Mills ).
## Acknowledgements
I’d like to thank Florin Panaite for his helpful comments and for pointing out the references , and to Robert Oeckl for . |
warning/0003/gr-qc0003052.html | ar5iv | text | # Singularities in General Relativity coupled to nonlinear electrodynamics
## I Introduction
It is a well-known fact that some of the most important solutions of Einstein’s field equations (e.g. Friedmann-Robertson-Walker and Schwarzschild) are singular. However, our understanding of the nature of these singularities is still incomplete. For instance, the cosmic censorship conjecture was put forward by R. Penrose in 1969 , but there is still no general proof of it. As a consequence of this lack of understanding, solutions that are everywhere regular and share some of the properties of singular solutions deserve attention. This is precisely the case of the “regular black hole” spacetimes recently exhibited in . These solutions were obtained for a very special type of source: an electric field that obeys a nonlinear electrodynamics. The authors of analyzed some of the features of the solution, but left aside others that are relevant. We shall re-examine this solution in detail. More importantly, we shall show in this particular example the far-reaching consequences of the fact that in nonlinear electromagnetism photons do not propagate along null geodesics of the background geometry. They propagate instead along null geodesics of an effective geometry, which depends on the nonlinearities of the theory. This result, derived by Plebańsky for Born-Infeld electrodynamics , was generalized for any nonlinear theory by Gutiérrez et al , and later independently rediscovered by Novello et al . Let us mention that the propagation of photons beyond Maxwell electrodynamics has been studied in several different situations. It has been investigated in curved spacetime, as a consequence of non-minimal coupling of electrodynamics with gravity , and in nontrivial QED vacua as an effective modification induced by quantum fluctuations . Nearly always, these analysis have had some unexpected results. As an example, let us mention the possibility of faster and slower-than-light photons .
Our main concern in this article will be then to show that one must consider the modifications on the trajectories of the photons induced by the nonlinearities of the electromagnetic theory in order to give a complete characterization of spacetimes with a nonlinear electromagnetic source. The structure of the paper is the following. A summary of the the solution given in and the properties studied there will be given in Section II, along with some interesting properties that went unnoticed before. In Section III we briefly review the origin of the effective geometry for photons in nonlinear electrodynamics. We shall use in Section IV the method of the effective geometry to study the features of the structure that photons see when travelling in the geometry given in . We close with some conclusions.
## II Details of the solution
Ayón Beato and García have found an exact solution of Einstein’s equations in the presence of a nonlinear electromagnetic source. The relevant equations are derived from the action
$$𝒮=d^4x\left[\frac{1}{16\pi }R\frac{1}{4\pi }(F)\right],$$
(1)
where $`R`$ is the curvature scalar and $``$ is a nonlinear function of $`F\frac{1}{4}F_{\mu \nu }F^{\mu \nu }`$. Following and this system could also be described using another function obtained by means of a Legendre transformation:
$$2F_F.$$
(2)
($`_F`$ denotes the derivative of $``$ w.r.t. $`F`$). With the definition
$$P_{\mu \nu }_FF_{\mu \nu },$$
(3)
it can be shown that $``$ is a function of $`P\frac{1}{4}P_{\mu \nu }P^{\mu \nu }=(_F)^2F`$, i.e., $`d=(_F)^1d((_F)^2F)=_PdP`$. With the help of $``$ one could express the nonlinear electromagnetic Lagrangian in the action (1) as $`=2P_P`$, which depends on the anti–symmetric tensor $`P_{\mu \nu }`$. The solution of Einstein’s equations coupled to nonlinear electrodynamics obtained in was derived from the following source:
$$(P)=P\frac{\left(13\sqrt{2q^2P}\right)}{\left(1+\sqrt{2q^2P}\right)^3}\frac{3}{2q^2s}\left(\frac{\sqrt{2q^2P}}{1+\sqrt{2q^2P}}\right)^{5/2},$$
(4)
where $`s=|q|/2m`$ and the invariant $`P`$ is a negative quantity. The corresponding Lagrangian is given by
$``$ $`=`$ $`P{\displaystyle \frac{\left(18\sqrt{2q^2P}6q^2P\right)}{\left(1+\sqrt{2q^2P}\right)^4}}`$ (6)
$`{\displaystyle \frac{3}{4q^2s}}{\displaystyle \frac{(2q^2P)^{5/4}\left(32\sqrt{2q^2P}\right)}{\left(1+\sqrt{2q^2P}\right)^{7/2}}}.`$
From Eq.(1) we get the following equations of motion:
$$G_\mu ^\nu =2(_PP_{\mu \lambda }P^{\nu \lambda }\delta _\mu ^\nu (2P_P)),$$
(7)
$$_\mu P^{\alpha \mu }=0.$$
(8)
This system was solved in , and the explicit form of the solution is the following:
$`ds^2`$ $`=`$ $`\left[1{\displaystyle \frac{2mr^2}{(r^2+q^2)^{3/2}}}+{\displaystyle \frac{q^2r^2}{(r^2+q^2)^2}}\right]dt^2`$ (10)
$`\left[1{\displaystyle \frac{2mr^2}{(r^2+q^2)^{3/2}}}+{\displaystyle \frac{q^2r^2}{(r^2+q^2)^2}}\right]^1dr^2r^2d\mathrm{\Omega }^2,`$
$`E_r=qr^4\left[{\displaystyle \frac{r^25q^2}{(r^2+q^2)^4}}+{\displaystyle \frac{15}{2}}{\displaystyle \frac{m}{(r^2+q^2)^{7/2}}}\right].`$ (11)
By means of the substitution $`x=r/|q|`$ we can rewrite $`g_{tt}`$ and $`E_r`$ as follows
$$g_{tt}=A(x,s)1\frac{1}{s}\frac{x^2}{(1+x^2)^{3/2}}+\frac{x^2}{(1+x^2)^2},$$
(12)
$$E_r=\frac{x^4}{q}\left[\frac{x^25}{(x^2+1)^4}+\frac{15}{4s}\frac{1}{(x^2+1)^{7/2}}\right].$$
(13)
The result of the analysis made in is that this metric describes a regular black hole. The position of the horizons was identified there with the values of the coordinate $`x`$ for which $`g_{tt}`$ is zero. These are given by
$$s=\frac{x^2\sqrt{x^2+1}}{x^4+3x^2+1}.$$
(14)
Accordingly, the solution has two horizons (for $`0<s<0.317`$), one horizon (for $`s=0.317`$), or no horizons (for $`s>0.317`$). It was also stated that this solution is regular, on the basis of the finiteness of the three invariants $`R`$, $`R_{\mu \nu }R^{\mu \nu }`$, and $`R_{\mu \nu \alpha \beta }R^{\mu \nu \alpha \beta }`$ <sup>1</sup><sup>1</sup>1 We have checked that all the components of $`R_{ABCD}`$ and $`C_{ABCD}`$ w.r.t a static observer are finite at $`r=0`$..
Let us point out now some features of the solution described by Eqns.(12) and (13) that were not noticed in . First, the behaviour of the radial component of the electric field depends on the value of $`s`$. Specifically, $`E_r`$ may have a zero; its position is given by
$$s=\frac{15}{4}\frac{\sqrt{x^2+1}}{x^25}.$$
(15)
Consequently, $`E_r`$ does not have zeros for $`0<s<3/4`$. For $`s3/4`$, $`E_r`$ has one zero located in the interval $`(0,\sqrt{5})`$ of the coordinate $`x`$. These features of the electric field are depicted in Figure 1 <sup>2</sup><sup>2</sup>2The plots in this paper have been done with gnuplot . .
Another salient feature of $`E_r`$ is that its energy density, calculated as the $`G_t^t`$ component of the Einstein tensor <sup>3</sup><sup>3</sup>3 This and other calculations in this paper were done with the package Riemann . may be negative for some interval of $`x`$. In fact, the expression
$$G_t^t=\rho =\frac{1}{sq^2}\frac{s\sqrt{1+x^2}(x^23)+3(x^2+1)}{(1+x^2)^{7/2}}$$
(16)
is zero for
$$s=3\frac{\sqrt{1+x^2}}{x^23}.$$
(17)
For $`s<1`$, the energy is always positive, but for $`s1`$ it has a zero given by Eq.(17). Figure 2 illustrates the situation.
## III Effective geometry for photons
In this section we give a summary of the method of the effective geometry . We will deal here only with the case in which the Lagrangian of the nonlinear electromagnetic theory is a function of $`F`$ only. The general case in which $``$ depends also on $`G=\frac{1}{2}F^{\mu \nu }\eta _{\mu \nu }^{\alpha \beta }F_{\alpha \beta }`$ is analyzed in . Based on the framework introduced by Hadamard , Novello et al showed that the discontinuities of the electromagnetic field propagate according to the equation
$$(_F\eta ^{\mu \nu }4_{FF}F^{\mu \alpha }F_\alpha {}_{}{}^{\nu })k_\mu k_\nu =0,$$
(18)
where $`\eta _{\mu \nu }`$ is the (flat) background metric, and $`k^\mu `$ is the propagation vector. This expression suggests that the self-interaction of the field $`F^{\mu \nu }`$ can be interpreted as a modification on the spacetime metric $`\eta _{\mu \nu }`$, leading to the effective geometry
$$g_{(\mathrm{eff})}^{\mu \nu }=_F\eta ^{\mu \nu }4_{FF}F_{}^{\mu }{}_{\alpha }{}^{}F^{\alpha \nu }.$$
(19)
Note that only in the particular case of linear Maxwell electrodynamics the discontinuities of the electromagnetic field propagate along the null cones of the Minkowskian background.
The general expression of the effective geometry can be equivalently written in terms of the energy-momentum tensor, given by
$$T_{\mu \nu }\frac{2}{\sqrt{\gamma }}\frac{\delta \mathrm{\Gamma }}{\delta \gamma ^{\mu \nu }},$$
(20)
where $`\mathrm{\Gamma }`$ is the effective action
$$\mathrm{\Gamma }d^4x\sqrt{\gamma }L,$$
(21)
and $`\gamma _{\mu \nu }`$ is the Minkowski metric written in an arbitrary coordinate system; $`\gamma `$ is the corresponding determinant. In the case of one-parameter Lagrangians, $`=(F),`$ we obtain
$$T_{\mu \nu }=4_FF_{\mu }^{}{}_{}{}^{\alpha }F_{\alpha \nu }\eta _{\mu \nu },$$
(22)
where we have chosen an Cartesian coordinate system in which $`\gamma _{\mu \nu }`$ reduces to $`\eta _{\mu \nu }.`$ In terms of this tensor the effective geometry (19) can be re-written as
$$g_{(\mathrm{eff})}^{\mu \nu }=\left(_F+\frac{_{FF}}{_F}\right)\eta ^{\mu \nu }+\frac{_{FF}}{_F}T^{\mu \nu }.$$
(23)
It is shown in that the field discontinuities propagate along the null geodesics of the effective geometry given by Eq.(23). This equation explicitly shows that the stress-energy distribution of the field is the true responsible for the deviation of the geometry felt by photons, from its Minkowskian form <sup>4</sup><sup>4</sup>4For $`T_{\mu \nu }=0`$, the conformal modification in (23) clearly leaves the photon paths unchanged..
We will show now that the modification of the underlying spacetime geometry seen by photons due to nonlinear electrodynamics can be also described as if photons governed by Maxwell electrodynamics were propagating inside a dielectric medium. In this last case, the electromagnetic field is represented by two antisymmetric tensors, the electromagnetic field $`F_{\mu \nu }`$ and the polarization field $`P_{\mu \nu }`$. For electrostatic fields inside isotropic dielectrics it follows that $`P_{\mu \nu }`$ and $`F_{\mu \nu }`$ are related by
$$P_{\mu \nu }=ϵ(E)F_{\mu \nu }.$$
(24)
where $`ϵ`$ is the electric susceptibility. Comparing with Eq.(3) we see that we can make the identification
$$_Fϵ,$$
(25)
which implies
$$_{FF}\frac{ϵ^{}}{4E},$$
(26)
in which $`ϵ^{}dϵ/dE`$ and $`E^2E_\alpha E^\alpha >\mathrm{\hspace{0.17em}0}.`$ Therefore, every Lagrangian $`=(F)`$ which describes a nonlinear electromagnetic theory may be used as a convenient description of Maxwell theory inside isotropic nonlinear dielectric media. Conversely, results obtained in the latter context can be restated in terms of Lagrangians of nonlinear theories. Using this equivalence, the effective geometry can be rewritten as
$$g_{(\mathrm{eff})}^{\mu \nu }=ϵ\eta ^{\mu \nu }\frac{ϵ^{}}{E}\left(E^\mu E^\nu E^2\delta _t^\mu \delta _t^\nu \right).$$
(27)
In other words,
$`g_{(\mathrm{eff})}^{tt}`$ $`=`$ $`ϵ+ϵ^{}E,`$ (28)
$`g_{(\mathrm{eff})}^{ij}`$ $`=`$ $`ϵ\delta ^{ij}{\displaystyle \frac{ϵ^{}}{E}}E^iE^j.`$ (29)
This shows that the discontinuities of the electromagnetic field inside a nonlinear dielectric medium propagate along null cones of an effective geometry (given by Eqn.(27)) which depends on the characteristics of the medium.
Although in the background was flat, the method can also be used in a curved background. The reason is that the equations given in are valid locally in any curved spacetime. Then from the Equivalence Principle follows that the only change in Eq.(19) is that of $`\eta _{\mu \nu }`$ by $`g_{\mu \nu }`$.
## IV Analysis of the “regular black hole”
Using Eqns.(28) and (29) it follows that the effective metric associated to a spherically symmetric solution of Einstein’s equations is given by
$$ds^2=\frac{1}{\mathrm{\Phi }(r)}\left[A(r)dt^2A(r)^1dr^2\right]\frac{r^2}{_F}d\mathrm{\Omega }^2,$$
(30)
where
$$\mathrm{\Phi }=ϵ+\frac{dϵ}{dE_r}E_r=\frac{2q}{r^3}\frac{1}{\frac{dE_r}{dr}}$$
(31)
and
$$ϵ=\frac{1}{E_r}\sqrt{\frac{P_{\mu \nu }P^{\mu \nu }}{2}}.$$
(32)
For the case dealt with in the previous section, the function $`\mathrm{\Phi }`$ takes the form
$`\mathrm{\Phi }(x,s)=`$ (33)
$`{\displaystyle \frac{8(x^2+1)^5s}{x^6(8x^4s104sx^2+80s+45x^2\sqrt{x^2+1}60\sqrt{x^2+1})}}.`$ (34)
From Eq.(30) we see that the $`tt`$ coefficient of the effective metric is given by the quotient $`g_{tt}^{(\mathrm{eff})}=A/\mathrm{\Phi }`$. The function $`\mathrm{\Phi }^1`$ has real zeros for
$$s=\frac{15}{8}\frac{\sqrt{x^2+1}(3x^24)}{x^413x^2+10}.$$
(35)
Taking into account that $`s`$ must be positive, we conclude from Eqn.(35) that the function $`\mathrm{\Phi }^1`$ has one zero for $`s<3/4`$ and two zeros for $`s3/4`$. In both cases the zeros are in the interval $`(0,3.49)`$ of the coordinate $`x`$.
It was shown in that the metric coefficient $`g_{tt}`$ given by Eq.(12) has two zeros for $`s<0.317`$, one zero for $`s=0.317`$, and no zeros for $`s>0.317`$. The zeros in $`g_{tt}`$ were identified in with horizons . We see that due to the effective metric, the geometry seen by the photons is more complex than the geometry seen by ordinary matter. Taking into account the zeros of $`A`$ and those of $`\mathrm{\Phi }^1`$ we conclude that $`g_{tt}^{(\mathrm{eff})}`$ has 3 zeros for $`s<0.371`$, two zeros for $`s=0.371`$, one zero for $`0.317<s<3/4`$, and again two zeros for $`s3/4`$.
To determine the nature of the new zeros in the metric, it is useful to study the effective potential that is felt by the photons. The symmetries of the metric imply that there are two Killing vectors and consequently, two conserved quantities:
$$E_0=g_{tt}\dot{t},\mathrm{and}h_0=\frac{r^2}{_F}\dot{\varphi }$$
(36)
(the overdot means derivative w.r.t the affine parameter). Standard calculations (see for instance ) using $`g_{\mu \nu }^{(\mathrm{eff})}`$ show that the effective potential for photons is given by
$$V_{\mathrm{eff}}=(1\mathrm{\Phi }^2)\frac{E_0^2}{2}+\frac{h_0^2}{x^2}_FA\mathrm{\Phi }$$
(37)
The explicit form of the effective potential is too involved to be displayed here. However, we note that $`V_{\mathrm{eff}}`$ has poles. One of them is at $`x=0`$, and the others are given by the expression of the poles of $`\mathrm{\Phi }`$ (see Eq.(35)), and those of $`_F`$ which are given by Eq.(15). $`_F`$ has no poles for $`0<s<3/4`$, and one pole for $`s3/4`$. Leaving aside the pole at $`x=0`$, it follows that for $`s<3/4`$, the effective potential has only one pole, and for $`s3/4`$, it has three poles. Those that originate in the singularities of the function $`\mathrm{\Phi }`$ are in agreement with the extrema of the electric field, as shown by Eq.(34). We give in Figs. 3, 4 and 5 plots of $`V_{\mathrm{eff}}`$ for different values of the relevant parameters.
Several comments are in order. The singularities in the potential suggest that the effective geometry itself may be singular. This is confirmed by the expression of the scalar curvature $`R^{(\mathrm{eff})}`$, which diverges in the values of $`x`$ given by Eqs.(15) and (35). Let us analize the relative position of these singularities felt by the photons and those of the metric coefficient $`g_{tt}(x,s)`$, given by Eq.(14). The information is conveniently summarized by the following plot:
We see that for a fixed $`s0.317`$ the singularities are situated inside the first horizon. However, for $`s>0.317`$ the singularities are not anymore hidden behind a horizon: we are then in the presence of naked singularities. We must remark that these singularities are only felt by photons. The rest of the matter follows geodesics of the regular spacetime given in .
It can also be seen from the plot that for $`s<0.371`$ the coordinate distance between the two horizons decreases for increasing $`s`$, up to $`s=0.371`$, where the two horizons coalesce.
Before analyzing the path of a photon coming from infinity, let us remark that there is a low potential barrier extending to the right of the outermost singularity for any value of the parameters. This barrier can be seen in Fig.3, and it is also present to the right of Fig.4. A low-energy photon incident from the right will find then this barrier, and will be deflected back to large values of $`x`$. This deflection will be more pronounced with increasing energy. When the energy of the photon is aproximately that of the height of the barrier, the photon can orbit around the center of the field in an unstable orbit. Finally an incident photon with energy greater that the height of the barrier will inevitably encounter the first singularity.
It is easily seen from Eq.(37) that the potential goes to zero for large values of $`x`$. We have also analyzed the effective potential for the case of a negative $`q`$, but the only quantitatively different result is a small increment of the innnermost local maximum seen in Figs.3 and 4.
We move now to another peculiar feature of the effective geometry. It is known that the effective potential for the Schwzarschild and Reissner-Nordstrom geometries is null in the case of photons with $`h_0=0`$. However, from Eq.(37) we see that in this case $`V_{\mathrm{eff}}`$ for the effective geometry reduces to
$$V_{\mathrm{eff}}=(1\mathrm{\Phi }^2)E_0^2$$
(38)
The dependence of this potential on $`\mathrm{\Phi }`$ is the same as in Eq.(37), so the behaviour of $`V_{\mathrm{eff}}`$ with $`x`$ in this case is qualitatively depicted in Figs.3 and 4.
Let us finally point out some unusual geometrical properties of the metric seen by the photons. The effective metric has the same symmetries of the original metric given by Eq.(10). It can be easily shown that the time Killing vector $`/t`$ is null on the hypersurfaces determined by the zeros of $`g_{tt}^{(\mathrm{eff})}`$.
Another interesting property of these surfaces is associated to the redshift of the photons. The redshift $`z`$ of a source as measured by an observer with velocity $`u^\mu `$ can be defined in terms of the frequency by
$$1+z=\frac{(u^\mu k_\mu )_{\mathrm{emitter}}}{(u^\mu k_\mu )_{\mathrm{observer}}}.$$
(39)
Considering a static observer for which $`u^\mu =\delta _0^\mu /\sqrt{g_{tt}}`$ this expression can be written as
$$1+z=\left[\frac{\sqrt{g_{tt}}}{g_{tt}^{(\mathrm{eff})}}\right]_{\mathrm{em}}\left[\frac{g_{tt}^{(\mathrm{eff})}}{\sqrt{g_{tt}}}\right]_{\mathrm{obs}}$$
(40)
Using the expression of the effective metric, and if the observer is at infinity,
$$1+z=\frac{\mathrm{\Phi }}{\sqrt{A}}$$
(41)
We conclude then that the redshift diverges in two cases: when $`A`$ is zero, and when $`\mathrm{\Phi }`$ diverges (see Fig.(6)).
## V Conclusion
The remarkable fact that in nonlinear electrodynamics the trajectories of photons are modified by the nonlinearities of the field equations has not been addressed frequently in the literature. The photons do not propagate following the null cones of the background metric but those of the effective metric. We have shown here the dramatic consequences that this has in a so-called regular black hole. In this case, there are singularities that are seen only by the photons. These singularities can either be hidden behind a horizon or naked, according to the value of the ratio $`q/2m`$. Let us remark that the existence of singularities in these type of solutions is a direct consequence of the existence of extrema of the electric field, as Eq.(31) shows. This is a general property which will always be present in any static and spherically symmetric solution of the system of equations (7) and (8) when the electromagnetic theory is nonlinear.
We have also shown that the effective potential to the right of the outermost singularity resembles that of Schwarzschild and Reissner-Nordstrom. However, contrary to what happens in Maxwell theory, photons with zero angular momentum travel under the influence of an effective potential that is different from zero.
We also exhibited some unusual properties of the solution found in . The electric field may have one or two extrema depending on the value of $`s`$. In the second case, it has a zero. Also, for certain values of $`s`$ the energy of the electric field is negative in some coordinate range. There are at least two more properties, geometrical in origin, that are worth of notice. First, the time Killing vector of the effective geometry is null in the surfaces where the function $`\mathrm{\Phi }`$ diverges. Second, the redshift measured by an observer far from the source diverges on the same surfaces. It is important to remark that these geometrical properties will be present in every solution with the same symmetries if the electric field has extrema.
To close, we would like to emphasize that ordinary matter follows geodesics of the background metric. However, the modifications of the metric induced by the nonlinearities of the electromagnetic field must always be taken into account when studying the propagation of photons. The abovementioned properties are nothing but a consequence of the nonlinearities of the electromagnetic theory.
## ACKNOWLEDGMENTS
SEPB would like to acknowledge financial support from CONICET-Argentina. MN and JMS acknowledge financial support from CNPq. The authors would like to thank K. Bronnikov for a useful remark. |
warning/0003/quant-ph0003021.html | ar5iv | text | # Casimir force at both non-zero temperature and finite conductivity
## ACKNOWLEDGMENTS
The authors are grateful to S. Reynaud for attracting their attention to Ref. and discussion. G.L.K. and V.M.M. are indebted to Center of Theoretical Sciences and Institute of Theoretical Physics of Leipzig University, where this work was performed, for kind hospitality. G.L.K. was supported by Graduate College on Quantum Field Theory at Leipzig University. V.M.M. was supported by Saxonian Ministry for Science and Fine Arts.
\[
\] |
warning/0003/cond-mat0003443.html | ar5iv | text | # An Explicit Form of the Equation of Motion of the Interface in Bicontinuous Phases
## 1 Introduction
The ordering process associated with the first order phase transition has been investigated for a long time since the early works by van der Waals appeared at the end of the last century. It has occupied a major position together with the dynamical critical phenomena as a subject not only of the non-equilibrium statistical physics but also of the field theory. A challenging feature of it is the strong nonlinearity which causes a variety of evolving spatial patterns. This has initiated the interesting concept of ‘pattern formation’ in far-from-equilibrium thermodynamic systems.-
Complicated patterns, at a glance, are seen even in the most simplified theoretical models of a binary mixture, e.g., in the time-dependent Ginzburg-Landau (TDGL) model for the conserved order parameter (COP), or the Cahn-Hilliard (CH) model. The nonlinearity in these systems finally gives rise to a spatial singularity of the order parameter gap, i.e., the interface with a sharp profile that separates distinctly the coexisting couple of phases. This is a primary example of the so-called ‘topological defect’ appearing in continuous fields. The main subject for us in this stage is to reduce a degenerate evolution equation for the interface from the basic TDGL or CH equations for the bulk field.
In spite of remarkable success of the intuitive, droplet theories on a dilute mixture, i.e., some pioneering works in the middle of this century, - the interface evolution equation for more general systems has not been obtained yet in an explicit form. The most general form we have found so far is an integral equation obtained by Kawasaki and Ohta. This is a kind of the curvature flow equation which represents the interface velocity with the mean curvature implicitly in contrast to the Allen-Cahn equation for a system of non-conserved order parameter (NCOP). It has been shown that this is the equivalent simple layer equation for the Dirichlet problem of the Laplace equation that is derived as a quasi-static approximation for a Stefan problem, i.e., a diffusion problem of an autonomous boundary condition evolving through the diffusion process itself. For some simplified systems, e.g., a planar interface or a thin system of spherical droplets, it is straightforward to derive the curvature flow equations, if they are considered as electro-static potential problems. The explicit potential solutions for such systems can be found in every elementary textbook on electro-statics. On the contrary, we will soon come up against a serious difficulty of a complicated geometry of the boundary condition when we try to find out explicit potential solutions for more general systems. For example, in a nearly symmetric mixture there appears a random bicontinuous structure, i.e., the so-called sponge phase. This has an infinitely multiply-connected topology and is characteristic of a three dimensional system. That contrasts to the case of a lamellar phase, which is essentially equivalent to a planar system.
In the previous work the author assumed an almost minimal surface for such a system mainly by intuition gained from some simulated pictures-, which just remind us a periodic minimal surface such as the Schwarz lattice structure, at least in a local view. The theoretical foundation for this assumption was that the mean curvature is included linearly in the boundary value itself as the Gibbs-Thomson condition already, and then it can be neglected within the limits of a linear theory with respect to the small mean curvature. This situation is similar to that of the decay process of a perturbative deformation on a planar interface. We used the common feature of zero mean curvature of the unperturbed basis in both systems. An important difference is that the sponge phase has an evident characteristic length, which was used effectively as the mean electro-static screening length in the previous work.
In the present paper it is shown that the assumption of a minimal surface on the unperturbed basis is not necessary in deriving an explicit solution from the integral equation and the same one-parameter curvature flow equation is rederived for more general bicontinuous systems. We need the bicontinuous nature of the interface only. Here ‘bicontinuous’ means that the interface divides the whole $`𝑹^3`$ space into a couple of (not necessarily symmetric) subspaces connected in each like a planar interface does, and not into subspaces more than three. This is the only condition required to make up a formal expression for the equivalent simple layer in terms of the potential problem.
In §2 the integral equation for the interface velocity is rederived on the basis of Onsager’s variational principle. It is shown that the principle of minimum dissipation is an example of the so-called gradient dynamics and is very useful for such an actual purpose to reduce the degenerate interface equation of motion correctly from the basic transport equation for the bulk field. In §3 a formal explicit expression for the simple layer on an arbitrary connected surface is derived referring to a solvable problem of a plane boundary. The equation of the level function that corresponds to the present curvature flow equation is discussed in §4.
## 2 The Onsager Principle and Basic Equations
In the original work by Kawasaki and Ohta a kind of the path-integral method and a variational principle associated with it were used. It seems instructive for us to follow it in terms of the familiar Onsager principle of minimum dissipation with reference to the recent gradient dynamics. To begin with, let us survey the Onsager variational principle.
Let $`𝒙=\{x_i\}`$ be a set of generalized thermodynamic variables, their phenomenological transport equations being given by
$$\dot{x}_i=\underset{j}{}L_{ij}X_j,$$
(1)
where $`\{L_{ij}\}`$ are the Onsager coefficients and $`𝑿=\{X_i\}`$ is a set of generalized thermodynamic forces, which are defined by using the entropy function $`S`$ as
$$𝑿=_𝒙S\mathrm{or}X_i=\frac{S}{x_i}.$$
(2)
Onsager proposed a minimum principle on the basis of the symmetric property of $`\{L_{ij}\}`$, i.e., $`L_{ij}=L_{ji}`$, and its positive definiteness as follows: Define two kinds of dissipation function and a Lagrangian by
$$\mathrm{\Phi }[\dot{𝒙},\dot{𝒙}]=\frac{1}{2}\underset{i,j}{}L_{}^{1}{}_{ij}{}^{}\dot{x}_i\dot{x}_j,$$
(3)
$$\mathrm{\Psi }[𝑿,𝑿]=\frac{1}{2}\underset{i,j}{}L_{ij}X_iX_j,$$
(4)
and
$`[\dot{𝒙},𝑿]`$ $`=`$ $`\mathrm{\Phi }[\dot{𝒙},\dot{𝒙}]\dot{S}[\dot{𝒙},𝑿]+\mathrm{\Psi }[𝑿,𝑿]`$ (5)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}L_{}^{1}{}_{ij}{}^{}(\dot{x}_i{\displaystyle \underset{k}{}}L_{ik}X_k)(\dot{x}_j{\displaystyle \underset{l}{}}L_{jl}X_l),`$
where
$$\dot{S}[\dot{𝒙},𝑿]=\underset{i}{}\dot{x}_iX_i.$$
(6)
Then the phenomenological equation Eq.(1) is given by
$$\delta =0(\mathrm{minimum})\mathrm{with}\mathrm{respect}\mathrm{to}\dot{𝒙}.$$
(7)
This variational principle was confirmed later as a problem of the most probable path in the path integral method of the linear fluctuation theory. It has been understood widely that the principle itself is nothing but a formal theory as is suggested in the above survey and is no use for an actual purpose to find out the transport equation itself for a given system.
Recently, mathematicians have introduced a notion of ‘gradient dynamics’. In this sense, the above Onsager’s phenomenology is a Lagrange multiplier version<sup>1</sup><sup>1</sup>1The original plan is that a dynamics is to be defined by maximizing $`\dot{S}=\dot{x}_𝒙S`$ with a constraint of a properly chosen inner product $`(\dot{𝒙},\dot{𝒙})=`$ constant. of a gradient dynamics defined by an inner product,
$$(𝒖,𝒗)=\underset{i,j}{}L_{}^{1}{}_{ij}{}^{}u_iv_j=2\mathrm{\Phi }[𝒖,𝒗].$$
(8)
Mathematicians may define an arbitrary dynamics by using an arbitrary inner product. That gives us a good hint: That is, it can be reasonably expected that we may find out the correct interface dynamics easily, when we rewrite the inner product, i.e., the physical dissipation function for the bulk into an interfacial form.
Let $`\{s(𝒓)\}`$ be the scalar field whose evolution is described by the TDGL equation of COP type, or equivalently the CH equation,
$$\frac{}{t}s(𝒓,t)=L^2\frac{\delta }{\delta s(𝒓)}S(\{s(𝒓)\}).$$
(9)
In this case $`S(\{s(𝒓)\})`$ is not the entropy but related to the free energy functional $`F(\{s(𝒓)\})`$ by $`S=F/kT`$, where $`T`$ is the temperature and $`k`$ the Boltzmann constant and
$$F(\{s(𝒓)\})=𝑑𝒓\frac{1}{2\chi _0}\left\{\frac{1}{2}s(𝒓)^2+\frac{1}{4s_{0}^{}{}_{}{}^{2}}s(𝒓)^4+\xi ^2[s(𝒓)]^2\right\}.$$
(10)
The equation (9) has a uniform equilibrium solution,
$$s(𝒓)=\pm s_0,$$
(11)
and a planar interface one, i.e., the so-called kink solution,
$$s_\mathrm{K}(z)=s_0\mathrm{tanh}\frac{z}{2\xi },$$
(12)
where $`z`$ is a normal coordinate perpendicular to the interface. Here $`\xi `$ corresponds to the thickness of the interface. The parameter $`\chi _0`$ is related to the susceptibility by
$$\chi _{0}^{}{}_{}{}^{1}=\frac{\delta ^2F}{\delta s(𝒓)^2}|_{\pm s_0}.$$
(13)
From Eq.(12) we get an expression of the surface tension,
$$\sigma =\frac{\xi ^2}{\chi _0}_{\mathrm{}}^{\mathrm{}}(s_\mathrm{K})^2𝑑z=\frac{\xi (2s_0)^2}{6\chi _0},$$
(14)
as the excess free energy stored in the interface layer where $`s\pm s_0`$.
Now let us derive the interfacial version from the bulk dissipation function defined by
$$\mathrm{\Phi }=\frac{1}{2L}G_0(𝒓𝒓^{})\frac{s(𝒓)}{t}\frac{s(𝒓^{})}{t}𝑑𝒓𝑑𝒓^{},$$
(15)
where $`G_0(𝒓𝒓^{})`$ is the inverse of the Onsager coefficient $`^2`$ in Eq.(9) defined by
$$^2G_0(𝒓𝒓^{})=\delta (𝒓𝒓^{}),$$
(16)
which is the Green function, i.e., the Newton (or the Coulomb) potential. If we assume that the every element $`dA`$ of the interface S at $`𝒂`$ propagates without deforming its profile, it is straightforward to show that
$$_{\mathrm{}}^{\mathrm{}}\frac{s(𝒓)}{t}𝑑z=2s_0v_n(𝒂).$$
(17)
from a geometrical consideration, where $`v_n(𝒂)`$ is the interface normal velocity. Thus, the interface version of the dissipation function Eq.(15) is given by
$$\mathrm{\Phi }=\frac{(2s_0)^2}{2L}_\mathrm{S}_\mathrm{S}G_0(𝒂𝒂^{})v_n(𝒂)v_n(𝒂^{})𝑑A𝑑A^{},$$
(18)
where the volume elements $`d𝒓`$ and $`d𝒓^{}`$ in Eq.(15) are replaced by $`dzdA`$ and $`dz^{}dA^{}`$. On the other hand, by definition of the surface free energy $`\sigma `$, the interface version of $`\dot{S}`$ is written as
$$\dot{S}=\frac{1}{kT}_\mathrm{S}\sigma H(𝒂)v_n(𝒂)𝑑A,$$
(19)
according to the formula on the change of the surface area caused by the normal displacement $`\delta z`$ of the surface S,
$$\delta _\mathrm{S}𝑑A=_\mathrm{S}H(𝒂)\delta z𝑑A,$$
(20)
where $`H(𝒂)`$ is the mean curvature of the surface at $`𝒂`$. Then, by using the variational principle $`\delta (\mathrm{\Phi }\dot{S})=0`$ with respect to $`v_n(𝒂)`$ we obtain an integral equation,
$$_\mathrm{S}G_0(𝒂𝒂^{})v_n(𝒂^{})𝑑A^{}=\frac{D\xi }{6}[H(𝒂)\overline{H}],$$
(21)
which is just the equation obtained by Kawasaki and Ohta. Here a new notation,
$$D=\frac{L}{kT\chi _0}$$
(22)
is introduced for the simplicity together with the relation Eq.(14). Note that a constant, $`\overline{H}`$, is incorporated as the Lagrange unknown multiplier in the present variational principle, corresponding to the constraint of order parameter conservation,
$$_\mathrm{S}v_n(𝒂)𝑑A=0.$$
(23)
If the inverse $`\mathrm{\Gamma }_0(𝒂,𝒂^{})`$ of the integral kernel $`G_0(𝒂𝒂^{})`$ satisfying
$$_\mathrm{S}\mathrm{\Gamma }_0(𝒂,𝒂^{\prime \prime })G_0(𝒂^{\prime \prime }𝒂^{})𝑑A^{\prime \prime }=\delta (𝒂𝒂^{}),$$
(24)
were given, the constant $`\overline{H}`$ would be determined by a weighted average,
$$\overline{H}=_\mathrm{S}_\mathrm{S}\mathrm{\Gamma }_0(𝒂,𝒂^{})H(𝒂^{})𝑑A𝑑A^{}/_\mathrm{S}_\mathrm{S}\mathrm{\Gamma }_0(𝒂,𝒂^{})𝑑A𝑑A^{}.$$
(25)
Apparently the integral equation Eq.(21) having a kernel of the Newton potential is the same as that of the equivalent simple layer for the Dirichlet problem of the Laplace equation. Let us relate it to the diffusion equation in the followings.
Because the order parameter deviation, $`\delta s`$, from the saturated level $`\pm s_0`$ at each phase is expected to be very small in the region out of the boundary layer ($`|z|>\xi `$), the chemical potential can be approximated by
$$\mu \frac{^2F}{s^2}|_{\pm s_0}\delta s=\chi _{0}^{}{}_{}{}^{1}\delta s.$$
(26)
Therefore, the limiting process in this region must be an ordinary diffusion process defined by
$$\frac{}{t}\psi =D^2\psi ,(D=L/kT\chi _0)$$
(27)
where a normalized variable, $`\psi =\delta s/s_0`$, is introduced. By using a kind of the boundary layer method we find a boundary condition,
$$2s_0\mu (0)=\sigma H(𝒂),$$
(28)
that is, the well-known Gibbs-Thomson condition for the curved interface having the mean curvature $`H=𝒏`$, where $`𝒏`$ is the normal unit vector of the interface. Then, the boundary value for $`\psi `$ must be given by
$$\psi _\mathrm{S}(𝒂)=\frac{1}{3}\xi H(𝒂).$$
(29)
It should be noted that the boundary S deforms itself via this diffusion with a normal velocity,
$$v_n(𝒂)=\frac{D}{2}\mathrm{}\left[𝒏\psi \right]_𝒂,$$
(30)
where $`\mathrm{}[\mathrm{}]_𝒂`$ denotes a gap of the enclosed quantity through the boundary layer at $`𝒂`$. These equations make a Stefan problem with a time-dependent, autonomous boundary condition.
In the late stage of the phase separation process where the characteristic length, $`\lambda `$, of the spatial pattern is sufficiently greater than the interface thickness $`\xi `$, the propagating velocity $`v_n(D\psi D\xi H/\lambda D\xi /\lambda ^2)`$ and the diffusion velocity $`v_\mathrm{D}(D\psi /\psi D/\lambda )`$ satisfy the condition
$$\frac{v_n}{v_\mathrm{D}}\frac{\xi }{\lambda }1.$$
(31)
Therefore, a quasi-static approximation,
$$^2\psi =0,$$
(32)
is applicable in this stage. Thus, the above Stefan problem becomes a Dirichlet problem of the Laplace equation. The propagating velocity Eq.(30) together with the boundary condition Eq.(29) is now given by the gap of the ‘electric field’ across the boundary, i.e., the equivalent ‘simple charge layer’, which should satisfy the integral equation Eq.(21). The additional parameter $`\overline{H}`$ should be regarded as a compatibility condition in this framework as follows: When $`\{\psi (𝒓)\}`$ is the solution for the boundary condition $`\{\psi _\mathrm{S}(𝒂)\}`$, $`\{\psi (𝒓)+c\}`$ is the solution for another boundary condition $`\{\psi _\mathrm{S}(𝒂)+c\}`$ with an arbitrary additional constant $`c`$. Obviously the family of the boundary conditions $`\{\psi _\mathrm{S}(𝒂)+c\}`$ have the same value of the field gap $`\mathrm{}[𝒏\psi ]_𝒂`$ or the simple layer solution. If the equivalent simple layer solution for this family exists, it determines a unique boundary value given by the left hand side of the integral equation. Then the right hand side should have an adjustable constant for the compatibility when an arbitrary boundary condition $`\{\psi _\mathrm{S}(𝒂)\}`$ is given. Note that this is true only if no flux lines escape out of the system into infinite points as is mentioned below. That corresponds to the conservation condition Eq.(23), because of the Gauss theorem, i.e., the total escaping flux being given by the total charge,
$$_\mathrm{S}v_n(𝒂)𝑑A,$$
(33)
included in the system. Thus the constant $`\overline{H}`$ must be determined by the conservation condition again. An exceptional example not covered by this rule is a problem of an isolated system composed of finite closed surface(s), where the flux may escape out. Of course we know that this is an ordinary case in the usual electro-static problem. For the simplicity suppose a finite sphere of radius $`R`$ and a constant boundary value $`\psi _\mathrm{S}`$ on it. The simple layer in this case is given by $`\mathrm{}[𝒏\psi ]=\psi _\mathrm{S}/R`$, because $`\psi (r)=R\psi _\mathrm{S}/r`$ for $`rR`$ and $`\psi (r)=\psi _\mathrm{S}`$ for $`r<R`$. Evidently, another boundary condition, $`\psi _\mathrm{S}+c`$, has a different value of the simple layer, $`(\psi _\mathrm{S}+c)/R`$. Thus, for the problem of isolated closed surface(s), the compatibility constant is not necessary. In contrast to it, for the extraordinary case of infinitely extended systems such as a system of spheres scattered homogeneously in the whole $`𝑹^3`$ space or a bicontinuous system spreading over the whole $`𝑹^3`$ space, etc., we need the compatibility condition. In these systems the constant boundary values induce a unique constant potential in both sides of the interface(s).
## 3 Derivation of the Explicit Solution
### 3.1 A planar boundary
First let us consider a Dirichlet problem of the Laplace equation,
$$^2\psi =0,$$
(34)
when the values $`\{\psi _\mathrm{S}(x,y)\}`$ on the $`xy`$-plane, S, are given. Let $`\sigma _{\mathrm{ind}}(x^{},y^{};z)`$ be the surface charge density on the grounded conductor S induced by a probe charge $`1`$ located at a referred point $`(0,0,z)`$. Then the potential $`\psi `$ at the referred point $`(0,0,z)`$ is expressed in an explicit form as
$$\psi (0,0,z)=\sigma _{\mathrm{ind}}(x^{},y^{};z)\psi _\mathrm{S}(x^{},y^{})𝑑x^{}𝑑y^{},$$
(35)
by using the elementary method of the Green function. The induced charge density $`\sigma _{\mathrm{ind}}`$ on the plane boundary can be found in every textbook on electro-statics as the simplest example of the mirror image method. By using the actual expression for it together with the Taylor expansion,
$$\psi _\mathrm{S}(x^{},y^{})=\mathrm{exp}(𝒂^{}_\mathrm{S})\psi _\mathrm{S}(0,0),$$
(36)
where $`𝒂^{}=(x^{},y^{})`$ and $`_\mathrm{S}=(/x,/y)`$, Eq.(35) is rewritten as
$`\psi (0,0,z)`$ $`=`$ $`{\displaystyle \frac{z}{2\pi (x^2+y^2+z^2)^{3/2}}\mathrm{exp}(𝒂^{}_\mathrm{S})\psi _\mathrm{S}(0,0)𝑑x^{}𝑑y^{}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑a^{}{\displaystyle \frac{za^{}}{(a^2+z^2)^{3/2}}}J_0(a^{}|i_\mathrm{S}|)\psi _\mathrm{S}(0,0)`$
$`=`$ $`\mathrm{exp}(z|i_\mathrm{S}|)\psi _\mathrm{S}(0,0),`$
where $`J_n(r)`$ is the usual Bessel function of $`n`$-th order. Then, by differentiating it with respect to $`z`$, we find
$$𝒏\psi |_{z=+0}=|i_\mathrm{S}|\psi _\mathrm{S}(𝒂).$$
(38)
Note that this result itself means merely a transform of the boundary condition from a Dirichlet type into a Neumann type. Combining with the same, symmetric result for the opposite side, $`z0`$, we obtain the formula for the equivalent simple layer,
$$\sigma _{\mathrm{eq}}(𝒂)=\mathrm{}\left[𝒏\psi \right]_𝒂=2|i_\mathrm{S}|[\psi _\mathrm{S}(𝒂)\overline{\psi _\mathrm{S}}],$$
(39)
by means of which we can write the potential in another explicit form,
$$\psi (𝒓)=_\mathrm{S}G_0(𝒓𝒂^{})\sigma _{\mathrm{eq}}(𝒂^{})𝑑A^{}.$$
(40)
The compatibility constant mentioned in the last part in §2 is incorporated here, although it would disappear after the differentiation in the present case. It should be noted that, from the view-point of the Green function method this term is required, in principle, for the condition $`\psi 0`$ when $`|z|\mathrm{}`$ on applying the Green theorem.
The fictitious operator $`|i_\mathrm{S}|`$ was introduced by Ohta and Nozaki in a perturbational theory for a planar interface. They defined it by a Gaussian integral. However, in the context of the present derivation it should be interpreted as the following limit,
$$|i_\mathrm{S}|=\underset{\lambda \mathrm{}}{lim}\frac{1}{\lambda }V(\lambda ^2_{\mathrm{S}}^{}{}_{}{}^{2}),$$
(41)
where $`_{\mathrm{S}}^{}{}_{}{}^{2}(=^2/x^2+^2/y^2)`$ is the surface Laplacian and the function $`V(Q^2)`$ is defined by
$`V(Q^2)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n1)[(2n)!!]^2}}(Q^2)^n`$ (42)
$`=`$ $`J_0(Q)QJ_1(Q)+Q{\displaystyle _0^Q}J_0(Q^{})𝑑Q^{}.`$
The asymptotic behavior for large $`|Q|`$ is evaluated as
$$V(Q^2)|Q|,$$
(43)
which results in the well-known $`|q|^3`$ dispersion relation- for the relaxation mode on a planar interface obeying the decay law $`\mathrm{exp}(D\xi q^3t/3)`$. The formula Eq.(42) was derived in the previous paper by introducing a cut-off length or the upper limit $`\lambda `$ for the integration in Eq.(3.1). This time the cut-off is temporarily required to integrate the series expansion for the Bessel function $`J_0(a^{}|i_\mathrm{S}|)`$, or for $`\mathrm{exp}(𝒂^{}_\mathrm{S})`$, term by term. It will be regarded as a physical screening length and has an essential role in the final formula for bicontinuous systems in the followings.
The above formulation is applicable to the two dimensional system also, where the function $`V(Q^2)`$ should be replaced by
$`V(Q^2)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n1)(2n)!}}(Q^2)^n`$ (44)
$`=`$ $`{\displaystyle \frac{2}{\pi }}\left[\mathrm{cos}Q+Q{\displaystyle _0^Q}{\displaystyle \frac{\mathrm{sin}Q^{}}{Q^{}}}𝑑Q^{}\right],(d=2).`$
The asymptotic behavior is the same as Eq.(43). Here it should be noted that the operator $`|i/x|`$ in this case is neither $`/x`$ nor $`|f/x|`$ if operated on $`f(x)`$.
### 3.2 General curved surface boundaries
Now let us discuss the general case of a parametric, curved surface S. Let $`(u_1,u_2,u_3)`$ be a curvilinear orthogonal coordinate system, where the surface S is defined by $`u_3=0`$. The Laplacian in this system is given by
$$^2=\frac{1}{\gamma }\underset{i=1}{\overset{3}{}}\frac{}{u_i}\frac{\gamma }{g_{i}^{}{}_{}{}^{2}}\frac{}{u_i},$$
(45)
where
$$g_i=\left|\frac{𝒓}{u_i}\right|=1/|u_i|,$$
(46)
is the linear metric ($`i`$-th element of the diagonalized metric tensor) and $`\gamma `$ is the Jacobian defined by $`\gamma =g_1g_2g_3`$. We need the infinitesimal vicinity of S only in the followings. Then it is always possible without loosing the generality to assume $`g_3(u_1,u_2,u_3)=1`$ everywhere by introducing a family of parallels of S for the simplicity. In addition, let $`(u_1,u_2)`$ be the orthogonal coordinate associated with the directions of principal curvature on S. Suppose that a probe charge $`1`$ is located at a referred point $`(0,0,z)`$ in this curvilinear coordinate and the surface S is grounded. The Green function $`G(u_1,u_2,u_3;z)`$ of the Laplace equation for this boundary condition obeys the Poisson equation,
$$\underset{i=1}{\overset{3}{}}\frac{}{u_i}ϵ_i\frac{}{u_i}G=\delta (u_1)\delta (u_2)\delta (u_3z),$$
(47)
with the constraint,
$$G(u_1,u_2,0;z)=0.$$
(48)
Here new parameters $`ϵ_1=g_2/g_1,ϵ_2=g_1/g_2`$ and $`ϵ_3=g_1g_2`$ are introduced for convenience. Note that the denominator $`\gamma `$ in Eq.(45) canceled out with the Jacobian, which had appeared in the right hand side of Eq.(47). This can be regarded as an electro-static problem with a planar boundary condition in a Euclidean space $`(u_1,u_2,u_3)`$ having an anisotropic, heterogeneous dielectric tensor, if the whole parametric space $`(u_1,u_2,u_3)`$ corresponds to the original physical space, and vice versa. For the time being let us assume it on condition that we need only the limit $`z0`$, and let it be discussed later.
Therefore, the elementary mirror image method is applicable to this case, taking account of the refraction of the flux lines,
$$𝑫=(ϵ_1\frac{G}{u_1},ϵ_2\frac{G}{u_2},ϵ_3\frac{G}{u_3}).$$
(49)
For the present purpose to find $`𝒏\psi `$ on S we need only the flux lines that run along the infinitesimal vicinity of S in order to calculate the induced charge $`\sigma _{\mathrm{ind}}`$, because the fields produced by the probe charge at $`z`$ and by its image at $`z`$ can be calculated separately.
The method of a reference frame used in the previous paper to calculate the refracted flux assuming a minimal surface can be extended to more general surfaces as follows: If we take a normalization, $`ϵ_i(0,0,z)=1`$, for the simplicity, the flux that arrives in the area element $`du_1du_2`$ on the plane S is related to the solid angle $`d\mathrm{\Omega }`$ when the flux started from the source point $`(0,0,z)`$, that is, $`D_3du_1du_2=d\mathrm{\Omega }/4\pi `$. Let $`d\stackrel{~}{u}_1d\stackrel{~}{u}_2`$ be the area element that $`d\mathrm{\Omega }`$ would cut from the plane S when it were extended straight without refraction. Then the solid angle $`d\mathrm{\Omega }`$ is given by
$$d\mathrm{\Omega }=\frac{z}{(\stackrel{~}{u}_{1}^{}{}_{}{}^{2}+\stackrel{~}{u}_{2}^{}{}_{}{}^{2}+z^2)^{3/2}}d\stackrel{~}{u}_1d\stackrel{~}{u}_2.$$
(50)
Thus, the required surface charge induced in the area element $`dA`$ on the original curved surface S is finally given by
$$\sigma _{\mathrm{ind}}(u_1,u_2;z)dA=\frac{G}{u_3}|_{u_3=0}dA=\frac{z}{2\pi (\stackrel{~}{u}_{1}^{}{}_{}{}^{2}+\stackrel{~}{u}_{2}^{}{}_{}{}^{2}+z^2)^{3/2}}d\stackrel{~}{u}_1d\stackrel{~}{u}_2,$$
(51)
where the relation $`dA=ϵ_3du_1du_2`$ is used and a factor 2 caused by the image charge is introduced. Note that the actual expression for the refraction law of the flux, or that for the mapping $`(u_1,u_2)(\stackrel{~}{u}_1,\stackrel{~}{u}_2)`$ are not necessary in the above calculations.
This is the method of the pseudo-conformal transformation introduced in the previous paper. Thus the assumption of a minimal surface used there has none of special meaning for us no longer.
### 3.3 The bicontinuous phase
The remaining procedure is almost the same as that for the planar boundary except for the following three facts: First, it should be noted that the operator we will face in the expected formula related to Eq.(3.1) after the Taylor expansion must be
$$\frac{^2}{\stackrel{~}{u}_{1}^{}{}_{}{}^{2}}+\frac{^2}{\stackrel{~}{u}_{2}^{}{}_{}{}^{2}}.$$
(52)
However, in the limit $`z0`$, this can be replaced by the desired surface Laplacian defined by
$$_{\mathrm{S}}^{}{}_{}{}^{2}=\frac{1}{g_1g_2}\left(\frac{}{u_1}\frac{g_2}{g_1}\frac{}{u_1}+\frac{}{u_2}\frac{g_1}{g_2}\frac{}{u_2}\right),$$
(53)
as is shown in Appendix A. Note that this fact does not mean a conformal mapping as is mentioned there.
Second, there still remains the question whether the whole physical space is covered by a parametric space $`(u_1,u_2,u_3)`$ or not. Apparently the answer is no in general except for simple curved surfaces having the same topology as that of a plane. In a complicated sponge structure the region described by a given set of parameters $`(u_1,u_2,u_3)`$ must be limited in the local cave around the referred point. In order to avoid this difficulty, let us introduce an upper-limit $`\lambda `$ for the integration in Eq.(3.1). This cut-off length was defined intuitively as an electro-static screening length, i.e., the effective diameter of local caves of the sponge-shaped conductor in the previous paper. In fact we have such a unique (but time-dependent), well-defined characteristic length in the bicontinuous phase that grows up starting from a quenched, homogeneous mixtures. It shou1d be noted that we have another difficulty that at least one of the parameters $`(u_1,u_2)`$ may become multi-valued, for example, on a part of the surface having a rotational symmetry like a catenoid. It seems that we have no problem in applying the image method to this case if we do not introduce any cut-off, as is exemplified using a two-dimensional solvable problem of a circle boundary in Appendix B. The cut-off procedure, however, causes literally cut-off of the multi-valued part in this case. Fortunately, in addition to the fact that we need only the very vicinity of S, the weight of the contribution of the boundary values at the points away from the referred point decreases as $`\sigma _{\mathrm{ind}}1/\stackrel{~}{r}^3`$, where $`\stackrel{~}{r}=(\stackrel{~}{u}_{1}^{}{}_{}{}^{2}+\stackrel{~}{u}_{2}^{}{}_{}{}^{2})^{1/2}`$ is essentially the distance along S as is suggested by the equality between the operator Eq.(52) and $`_{\mathrm{S}}^{}{}_{}{}^{2}`$. Thus, the cut-off procedure may be introduced without ruining the approximation.
The last point is the bicontinuous nature of the interface S. This is related to the equivalence of the boundary condition between the inner and the outer problems of S. For example, let us consider a non-bicontinuous system composed of several closed surfaces, $`\mathrm{S}_1,\mathrm{S}_2,\mathrm{}`$ . Evidently, each inner problem associated with each $`\mathrm{S}_i`$ has a single boundary $`\mathrm{S}_i`$ itself. On the other hand the sole outer problem has a united boundary $`\mathrm{S}_1+\mathrm{S}_2+\mathrm{}`$ . Thus the boundary conditions are not equivalent to each other. In contrast to it, we have just two regions separated by a single sheet of the interface S in the bicontinuous phase that is extended infinitely or is connected periodically in all three directions, and the difference between the ‘inner’ and ‘outer’ problems has no sense there. Both problems have exactly the same boundary in this case. Then the induced charge densities for both problems must coincide with each other when $`z\pm 0`$ as is seen in the above derivation, even if the both bulk regions are asymmetric at an average. Therefore, the values of the potential gradient $`𝒏\psi (𝒂)`$ constructed with the common boundary values $`\{\psi _\mathrm{S}(𝒂)\}`$ on S become equivalent one another.
Thanks to these facts we can construct the equivalent simple layer by using Eq.(39). Thus the explicit expression for the curvature flow equation,
$$v_n(𝒂)=\frac{D\xi }{3\lambda (t)}V(\lambda (t)^2_{\mathrm{S}}^{}{}_{}{}^{2})[H(𝒂)\overline{H}],$$
(54)
is obtained again for general bicontinuous phases. The parameter $`\lambda (t)`$ is the time-dependent characteristic length, which is to be determined self-consistently by this equation. On the assumption of scaling, i.e., the time-dependent similarity law, the well-known time-dependence $`\lambda (t)t^{1/3}`$ may be found by using a kind of dimension analysis on this equation. Note that the compatibility constant $`\overline{H}`$ in the right hand side is given by the simple surface average in this approximation, i.e.,
$$\overline{H}=_\mathrm{S}H(𝒂)𝑑A/_\mathrm{S}𝑑A.$$
(55)
For small $`Q`$, the function $`V(Q^2)`$ is expanded as
$$V(Q^2)\{\begin{array}{cc}1+\frac{Q^2}{4}+\mathrm{}\hfill & (d=3),\hfill \\ & \\ \frac{2}{\pi }(1+\frac{Q^2}{2}+\mathrm{})\hfill & (d=2).\hfill \end{array}$$
(56)
The first term expresses the mean-field evaporation-condensation process $`((H(𝒂)\overline{H}))`$, due to the imbalance of the mean curvature around its average value $`\overline{H}`$. The second term may be interpreted as a kind of surface diffusion $`(_{\mathrm{S}}^{}{}_{}{}^{2}H(𝒂))`$ due to the local nonuniformity of the curvature. On the contrary the asymptotic value of $`V(Q^2)`$ for $`Q\pm \mathrm{}`$ is given by Eq.(43), i.e., $`V(Q^2)|Q|`$, which corresponds to the $`|q|^3`$ relaxation mode $`(|_\mathrm{S}|H(𝒂))`$. Then the short wave-length local fluctuations ($`|q|\lambda (t)1`$) on the interface must decay rapidly and we can expect that the interface is always smooth within a spatial scope of $`\lambda (t)`$ almost everywhere.
The above behaviors of $`V(Q^2)`$ are shown in Fig.1. The asymptotic estimation ($`|Q|`$) is practically applicable for $`Q>2`$. However, the important modes for the interface evolution are included in the small wave-number region, $`|q\lambda (t)|<1`$, where $`V(Q^2)`$ deviates significantly from its asymptote.
## 4 Discussions
The Onsager principle of minimum dissipation seems to be really useful in such a practical problem to reduce the interface evolution equation from the bulk equation, at least if the effect of thermal fluctuations can be neglected. In the unstable phase separation process the main role of thermal fluctuations is creating nuclei or microscopic droplets at the very early stage. This effect may be incorporated as the initial condition in the picture of the interface dynamics. Another example is the NCOP system. The inverse of the Onsager coefficients in this case is merely $`L^1\delta (𝒓𝒓^{})`$. Then the dissipation function is given by
$$\mathrm{\Phi }=\frac{1}{2L}\left(\frac{s(𝒓)}{t}\right)^2𝑑𝒓=\frac{\chi _0\sigma }{2L\xi ^2}_\mathrm{S}v_n(𝒂)^2𝑑A,$$
(57)
where the definition of the surface tension Eq.(14) is used. On the other hand, $`\dot{S}`$ becomes
$$\dot{S}=\frac{1}{kT}_\mathrm{S}\left\{\sigma H(𝒂)2s_0h\right\}v_n(𝒂)𝑑A,$$
(58)
where $`h`$ is the uniform external field that should be included in the free energy in the following form,
$$hs(𝒓)𝑑𝒓,$$
for a NCOP system. Then using the variational method we find
$$v_n(𝒂)=\frac{L\xi ^2}{kT\chi _0}\left[H(𝒂)\frac{2s_0h}{\sigma }\right].$$
(59)
That is the well-known Allen-Cahn equation. In addition, putting $`H=(d1)/R`$ for the droplet system, we obtain a radius evolution equation,
$$\dot{R}=\frac{(d1)L\xi ^2}{kT\chi _0}\left(\frac{1}{R_c}\frac{1}{R}\right),$$
(60)
where $`R_c=(d1)\sigma /2s_0h`$ is the radius of the critical droplet. Thus we can get these important formulae straightforwardly and correctly by using the Onsager principle.
The one-parameter curvature flow equation for general COP systems is rederived without assuming a minimal surface. The bicontinuous nature is shown to be indispensable in the present derivation. Of course the result is applicable to a planar interface if we adopt an infinite cut-off length, $`\lambda \mathrm{}`$. Especially it should be noted that it gives an exact explicit solution of the Dirichlet problem for an arbitrary curved surface having the same topology as that of a plane. Such formulae have not been found in the textbook on the electro-statics or on the potential problem.
On the contrary, this result cannot be applied to a system of scattered spheres because it has not the bicontinuous nature, although it has been one of the simplest examples to study. Of course a more accurate argument on this system is possible by applying the standard electro-statics directly to Eq.(21) as was performed by Kawasaki and Ohta. Suppose each sphere of radius $`R_i`$ and charged by $`Q_i`$ is sufficiently separated from each other. Then Eq.(21) for the potential on the surface of the sphere $`i`$ becomes
$$\frac{Q_i}{4\pi R_i}+\underset{ji}{}\frac{Q_j}{4\pi R_{ij}}=\frac{D\xi }{6}\left[\overline{H}\frac{2}{R_i}\right],$$
(61)
where $`R_{ij}`$ is the distance between a couple of spheres $`i`$ and $`j`$ and is assumed as $`R_{ij}R_i,R_j`$. Then the surface charge density given by $`\dot{R}_i=Q_i/4\pi R_{i}^{}{}_{}{}^{2}`$ satisfies
$$\dot{R}_i+\underset{ji}{}\frac{R_{j}^{}{}_{}{}^{2}}{R_iR_{ij}}\dot{R}_j=\frac{D\xi }{3}\frac{1}{R_i}\left[\frac{1}{R_c(t)}\frac{1}{R_i}\right].$$
(62)
The second term in the left-hand side shows plainly the long-range interaction of the interface through the diffusion process. When this term is neglected for a sufficiently thin system, this becomes the Lifshitz-Slyozov equation, where the critical radius $`R_c(t)`$ in the right-hand side is to be determined by
$$R_c(t)=R=\underset{i}{}R_i/\underset{i}{}1,$$
(63)
according to the conservation condition $`_i4\pi R_{i}^{}{}_{}{}^{2}\dot{R}_i=0`$ in this system.
Lastly let us discuss the equation of the level function associated with the present curvature flow equation. This method was used first by Ohta et al in a special problem of the ordering process in a NCOP system with $`h=0`$ and has been developed by mathematicians independently: Let $`\{u(𝒓,t)\}`$ be the fictitious scalar field, the interface being defined by $`u(𝒓,t)=0`$. Then the interface normal velocity is represented as $`v_n=|u|^1(u/t)_{u=0}`$. Suppose we have an explicit curvature flow equation,
$$v_n(𝒂)=(\{H(𝒂^{})\}),$$
(64)
where $``$ may be a functional. Here the mean curvature is given by
$$H=𝒏=\left[|u|^1_\mathrm{S}u\right]_{u=0},$$
(65)
where $`_\mathrm{S}=𝒏(𝒏)`$. If we extend these equations to the bulk ($`u0`$), that makes a closure of $`\{u(𝒓,t)\}`$.
Thus the level function equation for our case is written as
$$\frac{u}{t}=\frac{D\xi }{3\lambda (t)}|u|V(\lambda (t)^2_{\mathrm{S}}^{}{}_{}{}^{2})\left[|u|^1_\mathrm{S}u+\overline{H}\right].$$
(66)
For a small value of $`u`$ we have a formula,
$$H(u)H(0)u\left[|u|^1(H^22K)+_{\mathrm{S}}^{}{}_{}{}^{2}|u|^1\right]_{u=0},$$
(67)
where $`K`$ is the Gauss curvature defined by $`K=1/R_1R_2`$ using the radii of the principal curvature, $`R_1`$ and $`R_2`$. Then $`\overline{H}`$ in the right hand side of Eq.(66) may be replaced by
$$\overline{H}=\overline{H(0)}+u\overline{\left[|u|^1(H^22K)\right]}_{u=0}.$$
(68)
The second term in the right hand side of Eq.(67) canceled out after the surface integration. Note that the quantity $`H^22K(=R_{1}^{}{}_{}{}^{2}+R_{2}^{}{}_{}{}^{2})`$ is positive definite. Especially for a homogeneous bicontinuous phase of a symmetric mixture, this can be estimated as
$$H^22K2K2\lambda (t)^2,$$
(69)
or alternatively we may use it as a definition of the parameter $`\lambda (t)`$.
Thus the level function equation for a COP system has very tough form and has several difficulties to be solved. None of satisfactory analyses on it has been found so far.
## Acknowledgements
The author would like to thank the members of the statistical physics group of Nara Women’s University, who invited him to have a lecture on the phase ordering process in the autumn of 1999 and stimulated him into the present work in preparing a lecture note.
## Appendix A on the pseudo-conformality
Let $`f(u_1,u_2)`$ be an arbitrary continuous function on S and suppose a harmonic function satisfying the Laplace equation
$$\left(_{\mathrm{S}}^{}{}_{}{}^{2}+\frac{1}{ϵ_3}\frac{}{u_3}ϵ_3\frac{}{u_3}\right)\psi (u_1,u_2,u_3)=0,$$
(70)
with the boundary condition
$$\psi (u_1,u_2,0)=f(u_1,u_2).$$
(71)
Let us generalize the reference frame on S defined in §3 to the bulk, i.e., to the region $`u_30`$ in the same manner. Then by using the Green function method a formal solution for the harmonic function $`\psi `$ is given by
$$\psi (u_1,u_2,u_3)=_\mathrm{S}\frac{\stackrel{~}{G}}{u_3^{}}|_{u_3^{}=0}f(𝒂^{})d\stackrel{~}{A}^{},$$
(72)
where
$$\stackrel{~}{G}(\stackrel{~}{u}_1^{},\stackrel{~}{u}_2^{},u_3^{};\stackrel{~}{u}_1,\stackrel{~}{u}_2,u_3)=\stackrel{~}{G}^+\stackrel{~}{G}^{},$$
(73)
and
$$\stackrel{~}{G}^\pm =\frac{1}{4\pi [(\stackrel{~}{u}_1^{}\stackrel{~}{u}_1)^2+(\stackrel{~}{u}_2^{}\stackrel{~}{u}_2)^2+(u_3^{}u_3)^2]^{1/2}},$$
(74)
which satisfy a Poisson equation
$$\left(\frac{^2}{\stackrel{~}{u}_{1}^{}{}_{}{}^{2}}+\frac{^2}{\stackrel{~}{u}_{2}^{}{}_{}{}^{2}}+\frac{^2}{u_{3}^{}{}_{}{}^{2}}\right)\stackrel{~}{G}^\pm =\delta (\stackrel{~}{u}_1^{}\stackrel{~}{u}_1)\delta (\stackrel{~}{u}_2^{}\stackrel{~}{u}_2)\delta (u_3^{}u_3).$$
(75)
Note that the reference frame is defined on taking the origin at the point just under the probe charge in §3. However, once the new Euclidean coordinate system is defined, one may shift the origin to any point. Thus, the reference coordinate $`(\stackrel{~}{u}_1,\stackrel{~}{u}_2,u_3)`$ is also used for the point of the probe charge in the above equations.
Now combining Eqs.(70),(72) and (75) it is straightforward to find
$$\left(\frac{^2}{\stackrel{~}{u}_{1}^{}{}_{}{}^{2}}+\frac{^2}{\stackrel{~}{u}_{2}^{}{}_{}{}^{2}}\right)f=_{\mathrm{S}}^{}{}_{}{}^{2}f,$$
(76)
in the limit, $`u_30`$. Here the following relation for small $`u_3`$,
$$\frac{1}{ϵ_3}\frac{}{u_3}ϵ_3\frac{\stackrel{~}{G}^\pm }{u_3}=H\frac{\stackrel{~}{G}^\pm }{u_3}+\frac{^2\stackrel{~}{G}^\pm }{u_{3}^{}{}_{}{}^{2}}\frac{^2\stackrel{~}{G}^\pm }{u_{3}^{}{}_{}{}^{2}},$$
(77)
is used. Of course this approximation is exact for a minimal surface ($`H=0`$). In general one may use the fact that $`\stackrel{~}{\sigma }=(\stackrel{~}{G}^\pm /u_3^{})_{u_3^{}=0}`$, which has an almost $`\delta `$-function shape, satisfies
$$\frac{\stackrel{~}{\sigma }}{u_3}\frac{1}{u_3}\stackrel{~}{\sigma },\frac{^2\stackrel{~}{\sigma }}{u_{3}^{}{}_{}{}^{2}}\frac{1}{u_{3}^{}{}_{}{}^{2}}\stackrel{~}{\sigma },$$
(78)
and so on, when $`u_3`$ is sufficiently small.
Using the same argument repeatedly on the functions $`_{\mathrm{S}}^{}{}_{}{}^{2}f,_{\mathrm{S}}^{}{}_{}{}^{4}f`$ …, we obtain
$$\left(\frac{^2}{\stackrel{~}{u}_{1}^{}{}_{}{}^{2}}+\frac{^2}{\stackrel{~}{u}_{2}^{}{}_{}{}^{2}}\right)^nf=_{\mathrm{S}}^{}{}_{}{}^{2n}f,$$
(79)
for an arbitrary integer $`n`$. This is the result desired in §3.
It should be noted that this does not mean a conformal transformation on S, because the definition of the reference coordinate itself depends on the probed point. Further the reference coordinate $`(\stackrel{~}{u}_1,\stackrel{~}{u}_2)`$ itself is not an orthogonal system on S in general. Then, strictly speaking, the operator $`^2/\stackrel{~}{u}_{1}^{}{}_{}{}^{2}+^2/\stackrel{~}{u}_{2}^{}{}_{}{}^{2}`$ is not an invariant Laplacian but merely a Taylor expansion operator around the referred point.
## Appendix B a solvable example in $`d=2`$
It is instructive to discuss a solvable problem in the present scheme, though it is stupid practically. Let us consider a circle of radius $`R`$ and assume a probe charge is located at a distance $`z`$ from it. A new curvilinear coordinate $`(u_1,u_2)`$ may be defined by
$$u_1=(R+z)\phi ,u_2=rR,$$
(80)
where $`(r,\phi )`$ is the usual polar coordinate. The two dimensional Laplacian is written in this coordinate as
$$^2=\frac{1}{g_1}\left(\frac{}{u_1}\frac{1}{g_1}\frac{}{u_1}+\frac{}{u_2}g_1\frac{}{u_2}\right),$$
(81)
where
$$g_1=|u_1|^1=\frac{R+u_2}{R+z},g_2=|u_2|^1=1.$$
(82)
Let us derive the refraction law of the flux $`𝑫`$ in the Euclidean space $`(u_1,u_2)`$ due to the anisotropic dielectric tensor $`(ϵ_1,ϵ_2)=(g_{1}^{}{}_{}{}^{1},g_1)`$, which depends on $`u_2`$ only. Let $`\alpha `$ be the angle between the flux line and the $`u_1`$-axis. Then we find that
$$ϵ_1\mathrm{tan}\alpha =\mathrm{constant},$$
(83)
because of the continuity relations across a boundary line parallel to the $`u_1`$-axis,
$$D\mathrm{sin}\alpha =D^{}\mathrm{sin}\alpha ^{},\frac{D\mathrm{cos}\alpha }{ϵ_1}=\frac{D^{}\mathrm{cos}\alpha ^{}}{ϵ_1^{}},$$
(84)
for the perpendicular and the parallel components, respectively. By using Eq.(83) the flux equation is given by
$$\frac{du_2}{du_1}=\frac{\mathrm{tan}\alpha _0}{ϵ_1}=\frac{R+u_2}{R+z}\mathrm{tan}\alpha _0,$$
(85)
or
$$\mathrm{log}(R+u_2)=\mathrm{log}(R+z)\frac{\mathrm{tan}\alpha _0}{R+z}u_1,$$
(86)
where $`\alpha _0`$ is the initial angle, which defines the reference coordinate by
$$\stackrel{~}{u}_1=z/\mathrm{tan}\alpha _0.$$
(87)
On putting $`u_2=0`$ in Eq.(86) we find the mapping, $`u_1\stackrel{~}{u}_1`$, i.e.,
$$\stackrel{~}{u}_1=\frac{z}{(R+z)\mathrm{log}(1+z/R)}u_1=\frac{z}{\mathrm{log}(1+z/R)}\phi .$$
(88)
Of course this is a multi-valued mapping, $`(\pi ,\pi )(\mathrm{},\mathrm{})`$. Then the induced charge density is given by
$$\sigma _{\mathrm{ind}}d\phi =\frac{1}{\pi }\frac{z}{\stackrel{~}{u}_{1}^{}{}_{}{}^{2}+z^2}d\stackrel{~}{u}_1=\frac{x}{\pi }\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{(\phi +2n\pi )^2+x^2}d\phi ,$$
(89)
where a new parameter $`x=\mathrm{log}(R+z)`$ is used. By using some formulae this yields
$`\sigma _{\mathrm{ind}}d\phi `$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{\mathrm{sinh}x}{\mathrm{cosh}x\mathrm{cos}\phi }}d\phi `$ (90)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{(2R+z)z}{R^2+(R+z)^22R(R+z)\mathrm{cos}\phi }}d\phi .`$
That is the well-known result obtained directly with use of a mirror image method.
Thus the multi-valued mapping has no problem in the present reference frame method. However, the cut-off procedure causes literally the cut-off of the tail part of the induced charge density. Here it is wrong to conclude that this effect would become infinitesimal in the limit $`z0`$ because the cut-off charge is of order of $`z/\lambda (t)`$. Note that the final formula for the equivalent simple layer is obtained after differentiating with respect to $`z`$. The weight for the contribution from the remote points becomes $`d\stackrel{~}{𝒖}/\stackrel{~}{u}^d`$, lacking the factor $`z`$. That estimates $`1/\lambda (t)`$ for the cut-off effect. Thus the cut-off loss must be recovered by some trick, such as the mean field term $`\overline{H}/\lambda (t)`$ in Eq.(54). |
warning/0003/physics0003007.html | ar5iv | text | # Performance of Discrete Heat Engines and Heat Pumps in Finite Time.
## I Introduction
Analysis of heat engines has been a major source of thermodynamic insight. The second law of thermodynamics resulted from Carnot’s study of the reversible heat engine . Study of the endo-reversible Newtonian engine began the field of finite time thermodynamics . Analysis of a virtual heat engine by Szilard led to the connection between thermodynamics and information theory . Recently this connection has been extended to the regime of quantum computation .
Quantum models of heat engines show a remarkable similarity to engines obeying macroscopic dynamics. The Carnot efficiency is a well established limit for the efficiency of lasers as well as other quantum engines . Moreover, even the irreversible operation of quantum engines with finite power output has many similarities to macroscopic endo-reversible engines .
It is this line of thought that serves as a motivation for a detailed analysis of a discrete four stroke quantum engine. In a previous study , the same model served to find the limits of the finite time performance of such an engine but with the emphasis on power optimization. In that study the working medium was composed of discrete level systems with the dynamics governed by a master equation. The purpose was to gain insight into the optimal engine’s performance with respect to time allocation when external parameters such as: the applied fields, the bath temperatures and the relaxation rates were fixed.
The present analysis emphasizes the reverse operation of the heat engine as a heat pump. For an adequate description of this mode of operation inner friction has to be a consideration. Without it the model is deficient with respect to optimizing the cooling power. Another addition is the optimization of the external fields. This is a common practice when cold temperatures are approached. With the addition of these two attributes, the four stroke quantum model is analyzed both as a heat engine and as a refrigerator.
Inner friction is found to have a profound influence on performance of the refrigerator. A direct consequence of the friction is a lower bound on the cycle time. This lower bound excludes the non-realistic global optimization solutions found for frictionless cases where the cooling power can be optimized beyond bounds. This observation, has led to the suggestion of replacing the optimization of the cooling power by the optimization of the cooling efficiency per unit time . Including friction is therefore essential for more realistic models of heat engines and refrigerators with the natural optimization goal becomes either the power output or the cooling power. The source of friction is not considered explicitly in the present model. Physically friction is the result of non-adiabatic phenomena which are the result of the rapid change in the energy level structure of the system. For example friction can be caused by the missalignement of the external fields with the internal polarization of the working medium. For a more explicit description of the friction the interactions between the individual particles composing the working fluid have to be considered. The present model is a microscopic analogue of the Ericsson refrigeration cycle where the working fluid consists of magnetic salts. The advantage of the microscopic model is that the use of the phenomenological heat transfer laws can be avoided . The results of the present model are compared to a recent analysis of macroscopic chillers . In that study, a universal modeling was demonstrated. It is found that the discrete quantum version of heat pumps has behavior similar to that of macroscopic chillers.
There is a growing interest in the topic of cooling atoms and molecules to temperatures very close to absolute zero . Most of the analysis of the cooling schemes employed are based on quantum dynamical models. New insight can be gained by employing a thermodynamic perspective. In particular the temperatures achieved are so low that the third law of thermodynamics has to be considered. The discrete level heat pump can serve as a model to study the third law limitations. The finite time perspective of the third law is a statement on the asymptotic rate of cooling as the absolute temperature is approached. These restrictions are imposed on the optimal cooling rate. The behavior of the optimal cooling rate as the absolute temperature is approached is a third law upper bound on the cooling rate. The main finding of this paper is that the optimal cooling rate converges to zero linearly with temperature, and the entropy production reaches a constant when the cold bath temperature approaches absolute zero.
## II Basic Assumptions and Formal Background for the Heat Engine and the Heat Pump
Heat engines and heat pumps are characterized by three attributes: the working medium, the cycle of operation, and the dynamics which govern the cycle. Heat baths by definition are large enough so that their temperatures is constant during the cycle of operation. The heat engine and the heat pump are constructed from the same components and differ only by their cycle of operation.
### A The Working Medium
The working medium consists of an ideal ensemble of many non-interacting discrete level systems. Specifically, the analysis is carried out on two-level systems (TLS) but an ensemble of harmonic oscillators would lead to equivalent results.
The TLS systems are envisioned as spin-1/2 systems. The lack of spin-spin interactions enables the description of the energy exchange between the working medium and the surroundings in terms of a single TLS. The state of the system is then defined by the average occupation probabilities $`P_+`$ and $`P_{}`$ corresponding to the energies $`\frac{1}{2}\omega `$ and $`\frac{1}{2}\omega `$, where $`\omega `$ is the energy gap between the two levels. The average energy per spin is given by
$`E=P_+\left({\displaystyle \frac{1}{2}}\omega \right)+P_{}\left({\displaystyle \frac{1}{2}}\omega \right)`$ (1)
The polarization, $`S`$, is defined by
$`S={\displaystyle \frac{1}{2}}(P_+P_{}),`$ (2)
and thus the energy can be written as $`E=\omega S`$. Energy change of the working medium can occur either by population transfer from one level to the other (changing S) or by changing the energy gap between the two levels (changing $`\omega `$). Hence
$`dE=Sd\omega +\omega dS.`$ (3)
Population transfer is the microscopic realization of heat exchange. The energy change due to external field variation is associated with work. Eq. (3) is therefore the first law of thermodynamics:
$`D𝒲Sd\omega ;D𝒬\omega dS.`$ (4)
Finally, for TLS the internal temperature, $`T^{}`$, is always defined via the relation
$`S={\displaystyle \frac{1}{2}}\mathrm{tanh}\left({\displaystyle \frac{\omega }{2k_BT^{}}}\right).`$ (5)
Note that the polarization $`S`$ is negative as long as the temperature is positive.
### B The Cycle of Operation
#### 1 Heat engines cycle
The cycle of operation is analyzed in terms of the polarization and frequency $`(S,\omega )`$. A schematic display is shown in Fig.(1) for a constant total cycle time, $`\tau `$. The present engine is an irreversible four stroke engine resembling the Stirling cycle, with the addition of internal friction. The direction of motion along the cycle is chosen such that net positive work is produced.
The four branches of the engine will be now briefly described.
On the first branch, $`AB`$, the working medium is coupled to the hot bath of temperature $`T_h`$ for period $`\tau _h`$, while the energy gap is kept fixed at the value $`\omega _b`$. The conditions are such that the internal temperature of the medium is lower than $`T_h`$. In this branch, the polarization is changing from the initial polarization $`S_2`$ to the polarization $`S_1`$. The inequality to be fulfilled is therefore:
$`S_1<{\displaystyle \frac{1}{2}}\mathrm{tanh}\left({\displaystyle \frac{\omega _b}{2k_BT_h}}\right).`$ (6)
Since $`\omega `$ is kept fixed, no work is done and the only energy transfer is the heat $`\omega _b(S_1S_2)`$ absorbed by the working medium.
In the second branch, $`BC`$ the working medium is decoupled from the hot bath for a period $`\tau _a`$, and the energy gap is varied linearly in time, from $`\omega _b`$ to $`\omega _a`$. In this branch work is done to overcome the inner friction which develops heat, causing the polarization to increase from $`S_1`$ to $`S_3`$ (Cf. Fig. 1). The change of the internal temperature is the result of two opposite contributes. First lowering the energy gap leads to a lower inner temperature for constant polarization $`S`$. Second increase in polarization due to friction, leads to an increase of the inner temperature for fixed $`\omega `$. The inner temperature $`T^{}`$ at point C might therefore be lower or higher than the initial temperature at point B.
The third branch $`CD`$, is similar to the first. The working medium is now coupled to a cold bath at temperature $`T_c`$ for time $`\tau _c`$. The polarization changes on this branch from $`S_3`$ to the polarization $`S_4`$. For the cycle to close, $`S_4`$ should be lower than $`S_2`$. At the end of the cycle the internal temperature of the working medium should be higher than the cold bath temperature, $`T^{}>T_c`$, leading to:
$`S_4>{\displaystyle \frac{1}{2}}\mathrm{tanh}\left({\displaystyle \frac{\omega _a}{2k_BT_c}}\right).`$ (7)
Since $`S_4<S_1`$ (Fig. 1), it follows from Eq. ( 6) and Eq.( 7), that:
$`\left({\displaystyle \frac{\omega _a}{T_c}}\right)>\left({\displaystyle \frac{\omega _b}{T_h}}\right)`$ (8)
Inequality (8) is equivalent to the Carnot efficiency bound, from Eq. (8) one gets:
$`1\left({\displaystyle \frac{\omega _a}{\omega _b}}\right)<1\left({\displaystyle \frac{T_c}{T_h}}\right)=\eta _{Carnot}`$ (9)
The present model is a quantum analogue of the Stirling engine which also has Carnot’s efficiency as an upper bound.
The polarization $`S`$ changes uni-directionally along the ’adiabats’ due to the increase of the excited level population as a result of the heat developed in the working fluid when work is done against friction, irrespective of the direction of the field change.
The fourth branch $`DA`$, closes the cycle and is similar to the second. The working medium is decoupled from the cold bath. In a period $`\tau _b`$ the energy gap is changing back to its original value, $`\omega _b`$. The polarization increases from $`S_4`$ to the original value $`S_2`$.
#### 2 Refrigerator cycle
The purpose of a heat pump is to remove heat from the cold reservoir by employing external work. The cycles of operation in the $`(S,\omega )`$ plane is schematically shown in Fig. 2,
The cycle of operation resembles the Ericsson refrigeration cycle . The differences are in the dynamics of the microscopic working fluid which are described in subsection II C. The work and heat transfer for the heat pump is summarized in Table II.
The four branches for the heat pump become:
In the first branch, $`DC`$, the working medium is coupled to the cold bath of temperature $`T_c`$ for time $`\tau _c`$, while the energy gap is kept fixed at the value $`\omega _a`$. The conditions are such that the internal temperature of the medium is lower than $`T_c`$ during $`\tau _c`$. Along this branch, the polarization changes from the initial polarization $`S_1`$ to the polarization $`S_2`$. Since $`\omega `$ is kept fixed, no work is done and the only energy transfer is the heat $`\omega _a(S_2S_1)`$ absorbed by the working medium. On this branch:
$`S_2<{\displaystyle \frac{1}{2}}\mathrm{tanh}\left({\displaystyle \frac{\omega _a}{2k_BT_c}}\right).`$ (10)
In the second branch, $`CB`$ the working medium is decoupled from the cold bath, and the energy gap is varied. In the frictionless case the polarization $`S_2`$ is constant (Left of Fig. 2). The only energy exchange is the work done on the system ( Table II). When friction is added the polarization is changing from $`S_2`$ to $`S_3`$ in a period $`\tau _a`$. The energy gap changes from $`\omega _a`$ to $`\omega _b`$ (Right of Fig. 2), according to a linear law. In addition to work, heat is developing as a result of the inner friction ( Table II).
The third branch $`BA`$, is similar to the first. The working medium is coupled to the hot bath at temperature $`T_h`$, for time $`\tau _h`$, keeping the energy gap $`\omega _b`$ fixed. In this branch the polarization changes from $`S_2`$ to $`S_1`$ in the frictionless case, and from $`S_3`$ to $`S_4`$ when friction is added. The constraint is that the internal temperature of the working medium should be higher than the hot bath temperature during the time $`\tau _h`$, $`T^{}>T_h`$, leading to the inequality (Fig. 2),
$`S_1>S_4>{\displaystyle \frac{1}{2}}\mathrm{tanh}\left({\displaystyle \frac{\omega _b}{2k_BT_h}}\right).`$ (11)
therefore $`S_2>S_1`$. From Eqs. (10) and (11), the condition for the interrelation between the bath temperatures and the field values becomes:
$`\left({\displaystyle \frac{\omega _a}{T_c}}\right)<\left({\displaystyle \frac{\omega _b}{T_h}}\right)`$ (12)
which is just the opposite inequality of the heat engine, (Eq. 8). In the heat pump work is done on the working fluid and since no useful work is done Carnot’s bound is not violated.
The fourth branch $`AD`$, closes the cycle and is similar to the second. The working medium is decoupled from the cold bath, and the energy gap changes back, during a period $`\tau _b`$ to its original value, $`\omega _b`$.
The results are summarized in Table II.
### C Dynamics of the working medium
The dynamics of the system along the heat exchange branches is represented by changes in the level population of the two-level-system. This is a reduced description in which the dynamical response of the bath is cast in kinetic terms . Since the dynamics has been described previously only a brief summary of the main points is presented here, emphasizing the differences in the energy exchanges on the ’adiabats’.
#### 1 The dynamics of the heat exchange branches
The dynamics of the population at the two levels, $`P_+`$ and $`P_{}`$, are described via a master equation
$`\{\begin{array}{c}\frac{dP_+}{dt}=k_{}P_++k_{}P_{}\hfill \\ \frac{dP_{}}{dt}=k_{}P_+k_{}P_{}\hfill \end{array},`$ (15)
where $`k_{}`$ and $`k_{}`$ are the transition rates from the upper to the lower level and vice versa. The explicit form of these coefficients depend on the nature of the bath and the system bath coupling interactions. The thermodynamics partition between system and bath is consistent with a weak coupling assumption . Temperature enters through detailed balance. The equation of motion for the polarization $`S`$ obtained from Eq. (15) becomes:
$`{\displaystyle \frac{dS}{dt}}=\mathrm{\Gamma }(SS^{eq})`$ (16)
where
$`\mathrm{\Gamma }=k_{}+k_{}`$ (17)
and
$`S^{eq}={\displaystyle \frac{1}{2}}{\displaystyle \frac{k_{}k_{}}{k_{}+k_{}}}={\displaystyle \frac{1}{2}}\mathrm{tanh}\left({\displaystyle \frac{\omega }{2k_BT}}\right)`$ (18)
where $`S^{eq}`$ is the corresponding equilibrium polarization. It should be noticed that in a TLS there is a one to one correspondence between temperature and polarization thus internal temperature is well defined even for non-equilibrium situations.
The general solution of Eq (16) is,
$`S(t)=S^{eq}+(S(0)S^{eq})e^{\mathrm{\Gamma }t}.`$ (19)
where S(0) is the polarization at the beginning of the branch.
From Eqs. (16) and (18) the rate of heat change becomes:
$`\dot{𝒬}=\omega \dot{S}`$ (20)
See also .
For convenience, new time variables are defined:
$`x=e^{\mathrm{\Gamma }_c\tau _c},y=e^{\mathrm{\Gamma }_h\tau _h}`$ (21)
These expressions represent a nonlinear mapping of the time allocated to the hot and cold branches by the heat conductivity $`\mathrm{\Gamma }`$. As a result, the time allocation and the heat conductivity parameter become dependent on each other.
Figure 1 and 2 show that the friction induces an asymmetry between the time allocated to the hot and cold branches since more heat has to be dissipated on the cold branch.
#### 2 The dynamics on the ’adiabats’
The external field $`\omega `$ and its rate of change $`\dot{\omega }`$ are control parameters of the engine. For simplicity it is assumed that the field changes linearly with time:
$`\omega (t)=\dot{\omega }t+\omega (0)`$ (22)
Rapid change in the field causes non-adiabatic behavior which to lowest order is proportional to the rate of change $`\dot{\omega }`$. In this context non-adiabatic is understood in its quantum mechanical meaning. Any realistic assumption beyond the ideal non-interacting TLS will lead to such non-adiabatic behavior. It is therefore assumed that the phenomena can be described by a friction coefficient $`\sigma `$ which forces a constant speed polarization change $`\dot{S}`$:
$`\dot{S}=\left({\displaystyle \frac{\sigma }{t^{}}}\right)^2`$ (23)
where $`t^{}`$ is the time allocated to the corresponding ’adiabat’. Therefore, the polarization as a function of time becomes:
$`S(t)=S(0)+\left({\displaystyle \frac{\sigma }{t^{}}}\right)^2t`$ (24)
where $`t0`$ , $`tt^{}`$. A modeling assumption of internally dissipative friction, similar to Eq.(23), was also made by Gordon and Huleihil (). Friction does not operate on the heat-exchange branches, there is no nonadiabtic effect since the fields $`\omega _a`$ and $`\omega _b`$ are constant in time. The irreversibilities on those branches are due to the transition rates ($`\mathrm{\Gamma }`$) of the master equation.
From Fig.(1), Eq. (4), and Eq. (24) the polarization, for the $`BC`$ branch of the heat engine becomes:
$`S_C=S_3=S_1+\left({\displaystyle \frac{\sigma ^2}{\tau _a}}\right).`$ (25)
The work done on this branch is:
$`𝒲_{BC}={\displaystyle _0^{\tau _a}}D𝒲={\displaystyle _0^{\tau _a}}S\dot{\omega }𝑑t=(\omega _a\omega _b)\left(S_1+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sigma ^2}{\tau _a}}\right)\right)`$ (26)
The heat generated on this branch in the working fluid, which is the work against the friction, becomes:
$`𝒬_{BC}={\displaystyle _0^{\tau _a}}D𝒬={\displaystyle _0^{\tau _a}}\omega \dot{S}𝑑t={\displaystyle \frac{\sigma ^2(\omega _a+\omega _b)}{2\tau _a}}`$ (27)
This work is dependent on the friction coefficient and inversely on the time allocated to the ’adiabats’. The computation for the other branches of the heat engine and heat pump are similar.
#### 3 Explicit expressions for the polarizations imposed by the closing of the cycle.
By forcing the cycle to close, the four corners of the cycle observed in Fig. 1 are linked. Applying Eq. (19) leads to the equations:
$`\begin{array}{c}S_1=S_2y+S_h^{eq}(1y)\hfill \\ S_3=S_1+\frac{\sigma ^2}{\tau _a}\hfill \\ S_4=S_3x+S_c^{eq}(1x)\hfill \\ S_2=S_4+\frac{\sigma ^2}{\tau _b}\hfill \end{array}`$ (32)
The solutions for $`S_1`$, $`S_2`$ and $`S_1S_2`$ are
$`\begin{array}{c}S_1=S_c^{eq}+\frac{\mathrm{\Delta }S^{eq}(1y)+\sigma ^2yG(x)}{(1xy)}=S_h^{eq}\frac{\mathrm{\Delta }S^{eq}y(1x)\sigma ^2yG(x)}{(1xy)}\hfill \\ S_2=S_c^{eq}+\frac{\mathrm{\Delta }S^{eq}x(1y)+\sigma ^2G(x)}{(1xy)}=S_h^{eq}\frac{\mathrm{\Delta }S^{eq}(1x)\sigma ^2G(x)}{(1xy)}\hfill \end{array}`$ (35)
and
$`S_1S_2=(\mathrm{\Delta }S^{eq})F(x,y){\displaystyle \frac{\sigma ^2(1y)G(x)}{(1xy)}}`$ (36)
where
$`F(x,y)={\displaystyle \frac{(1x)(1y)}{(1xy)}},\mathrm{\Delta }S^{eq}=(S_h^{eq}S_c^{eq}),G(x)=(x/\tau _a+1/\tau _b)`$ (37)
The constraint that the cycle must close leads to conditions on the polarizations $`S_1`$ and $`S_2`$ and on the minimum cycle time $`\tau _{c,min}`$. Eqs. (35) shows that both $`S_1`$ and $`S_2`$ are bounded by $`S_h^{eq}`$ and $`S_c^{eq}`$. The minimum cycle time is obtained when the polarizations coincide with the hot bath polarization: $`S_1`$=$`S_2`$=$`S_h^{eq}`$. In this case, $`\tau _h`$=0, and from Eqs. (21) and (36) the minimum time allocation on the cold bath $`\tau _{c,min}`$ is computed,
$`x_{max}={\displaystyle \frac{(S_h^{eq}S_c^{eq})\sigma ^2/\tau _b}{(S_h^{eq}S_c^{eq})+\sigma ^2/\tau _a}}`$ (38)
or
$`\tau _{c,min}=1/\mathrm{\Gamma }_c\mathrm{lg}{\displaystyle \frac{(S_h^{eq}S_c^{eq})\sigma ^2/\tau _b}{(S_h^{eq}S_c^{eq})+\sigma ^2/\tau _a}}`$ (39)
From this expression for $`\tau _{c,min}`$ the lower bound for the overall cycle time, is obtained (The left of Fig. 3) :
$`\tau \tau _{min}=\tau _{c,min}+\tau _a+\tau _b`$ (40)
When the minimum cycle time Eq. (39) diverges, the cycle cannot be closed. This condition imposes an upper bound on the friction coefficient $`\sigma `$
$`\sigma \sigma ^{up}=\sqrt{\tau _b(S_h^{eq}S_c^{eq})}.`$ (41)
or
$`\tau _b>\tau _{b,min}={\displaystyle \frac{\sigma ^2}{(S_h^{eq}S_c^{eq})}}.`$ (42)
Closing of the cycle imposes similar constraints on the minimal cycle time under friction for the heat pump. The value of the polarization difference $`S_2S_1`$ using the notation of Fig. 2 becomes:
$`S_2S_1=(S_2^{eq}S_1^{eq})F(x,y){\displaystyle \frac{\sigma ^2(1x)(y/\tau _a+1/\tau _b)}{(1xy)}}`$ (43)
The minimum cycle time is calculated in the limit when $`\tau _c`$=0, leading to $`S_2`$=$`S_1`$=$`S_2^{eq}`$. From Eqs. (21) and (43) the minimum time allocation on the hot branch $`\tau _{h,min}`$ is computed:
$`y_{max}={\displaystyle \frac{(S_2^{eq}S_1^{eq})\sigma ^2/\tau _b}{(S_2^{eq}S_1^{eq})+\sigma ^2/\tau _a}}`$ (44)
$`\tau _{h,min}=1/\mathrm{\Gamma }_h\mathrm{lg}{\displaystyle \frac{(S_2^{eq}S_1^{eq})\sigma ^2/\tau _b}{(S_2^{eq}S_1^{eq})+\sigma ^2/\tau _a}},`$ (45)
where $`S_2^{eq}`$ is point F and $`S_1^{eq}`$ is point E on Fig. 2. Using $`\tau _{h,min}`$ the lower bound for the overall cycle time, is computed
$`\tau \tau _{min}=\tau _{h,min}+\tau _a+\tau _b`$ (46)
Closing the cycle imposes a minimum cycle time for both the heat engine and the heat pump, which is a monotonically increasing function of the friction coefficient $`\sigma `$. The divergence of $`\tau _{min}`$ imposes a maximum value for the friction coefficient $`\sigma `$.
### D Finite Time Analysis
#### 1 Quantities to be Optimized.
The primary variable to be optimized is the power of the heat engine and the heat-flow extracted from the cold reservoir of the heat pump. For a preset cycle time, optimization of the power is equivalent to optimization of the total work, while optimization of heat flow is equivalent to the optimization of the heat absorbed. The entropy production will also be analyzed.
(1) The total work done on the environment per cycle of the Heat Engine.
The total work of the engine, is the sum of the work on each branch: Cf. (Table I and Fig. 1):
$`𝒲_{cyle1}={\displaystyle D𝒲}=\left(W_{AB}+W_{BC}+W_{CD}+W_{DA}\right)`$ (47)
which becomes:
$`𝒲_{cyle1}=(\omega _b\omega _a)(S_1S_2)\sigma ^2\omega _a(1/\tau _a+1/\tau _b)`$ (48)
The negative sign is due to the convention of positive $`𝒲`$ when work is done on the system.
Analyzing Eq. (48), the work is partitioned into three positive and negative areas. The positive area (left rotation)
$`𝒲_p=(\omega _b\omega _a)(S_1S_2)`$ (49)
is defined by the points $`A,B,C^1,D^1`$ in Fig. 1. The two negative areas (right rotation)
$`𝒲_n=\sigma ^2\omega _a(1/\tau _a)+\sigma ^2\omega _a(1/\tau _b)`$ (50)
are defined by the points $`C,C^1,S_1,S_3`$ and $`D^1,D,S_4,S_2`$ in Fig. 1.
The cycle which achieves the minimum cycle time $`\tau `$ = $`\tau _{c,min}`$, produces zero positive work $`𝒲_p=0`$. The corners A and B coincide at E, and $`C^1`$ coincides with $`D^1`$. The negative work of Eq. (50), is defined by the corners $`C,D,S_4,S_3`$ and is ’cut’ by the $`S_h^{eq}`$ line (Cf. the right of Fig. 4). The cycle has negative total work, meaning that work is done on the working fluid against friction. When $`\tau `$ increases beyond $`\tau _{c,min}`$ , $`S_1`$ diverts from $`S_2`$, becoming lower than $`S_h^{eq}`$ (Cf.Eq. (35)). At a certain point, the work done against friction is exactly balanced by the useful work of the engine. The minimum time in which this balance is achieved is designated $`\tau _0`$. Its value which can be deduced from Eq. (48) is worked out in appendix B.
The minimum cycle time $`\tau _{min}`$ is compared to $`\tau _0`$, the minimum time needed to obtain positive power shown in the right of Fig. 3 as a function of the friction $`\sigma `$. Both functions increase with friction, but $`\tau _0`$ diverges at a much lower friction parameter. Above this friction parameter no useful work can be obtained from the engine. The divergence of $`\tau _{min}`$ corresponds to a larger friction value where the cycle cannot be closed.
When the total time allocation is sufficient, i.e. $`\tau >\tau _0`$, work is done on the environment, and $`S_1`$ starts to increase. For long cycle times $`S_1`$ will approach $`S_h^{eq}`$, while $`S_2`$ will approach $`S_c^{eq}`$. The constant negative area will become negligible in comparison to the positive area ( Fig. 5).
To study the influence of friction on the work output the polarization difference from Eq. (36) S<sub>1</sub>-S<sub>2</sub> is inserted into the work expression Eq. (48), leading to:
$`𝒲_{cyle1}=(\omega _b\omega _a)(S_h^{eq}S_c^{eq})F(x,y)𝒲_{\sigma 1}`$ (51)
where
$`𝒲_{\sigma 1}=\sigma ^2\left({\displaystyle \frac{\omega _b(1y)(x/\tau _a+1/\tau _b)}{1xy}}+{\displaystyle \frac{\omega _a(1x)(1/\tau _a+y/\tau _b)}{1xy}}\right)`$ (52)
$`𝒲_{\sigma 1}`$ is the additional ’cost’ due to friction and is always positive.
The emergence of positive power $`𝒫`$ is shown in Fig. 4. For a fixed cycle time the optimization of work is equivalent to the optimization of power.
The first two cycles have a cycle time shorter than $`\tau _0`$, and therefore do not produce useful work. For cycle 3, $`\tau >\tau _0`$ and positive work is obtained when the time allocation on the cold bath is sufficient $`\tau _c0.08`$.
For longer total cycle times, the ratio between the negative area to the positive area decreases as can be seen in Fig. 5.
The position of the cycles in the $`S`$, $`\omega `$ coordinates relative to $`S_h^{eq}`$ and $`S_c^{eq}`$ changes as a function of the cycle time. Insight to the origin of the behavior of the ’moving’ cycles is presented in Fig. 11 of Appendix A.
The calculation of the total work done on the working fluid per cycle, $`𝒲_{cycle3}^{on}`$ for the heat pump is described in appendix D. See also ( Cf. Table (II) and Cf. Fig. 2).
(2) The heat-flow($`𝒬_F`$)
The heat-flow, $`𝒬_F`$, extracted from the cold reservoir is:
$`𝒬_F=\omega _a(S_2S_1)/\tau `$ (53)
Due to the dependence of $`𝒬_F`$ only on $`S_2`$-$`S_1`$, the cycle is similar to the cycle of the heat engine.
(3) The entropy production ($`\mathrm{\Delta }𝒮^u`$).
The entropy production of the universe, $`\mathrm{\Delta }𝒮^u`$, is concentrated on the boundaries with the baths since, for a closed cycle, the entropy of the working fluid is constant. The computational details for both the heat engine and the heat pump are shown in appendix C. The entropy production and the power have a reciprocal relation (See Fig. 12). For example, the entropy production increases with $`\sigma `$, while the power decreases.
(4) Efficiency.
The efficiency of the heat engine is the ratio of useful work to the heat extracted from the hot bath.
$$\eta _{H.E.}=\frac{𝒲_{cycle}}{𝒬_{absorbed}}=\eta _{H.E.}^{fricles}\left(\frac{\sigma ^2\omega _a(1/\tau _a+1/\tau _b)}{\omega _b(S_1S_2)}\right)$$
(54)
where $`\eta _{H.E.}^{fricles}=(1\omega _a/\omega _b)`$
When the cycle time approaches its minimum $`\tau \tau _{min},`$ the efficiency diverges: $`\eta _{H.E.}\mathrm{}`$. The efficiency becomes positive only when $`\tau \tau _0`$. Using Eq. (54) a bound for the efficiency is obtained:
$$0<\eta _{H.E.}\eta _{H.E.}^{fricles}\frac{T_c}{T_h}\left(\frac{\sigma ^2(1/\tau _a+1/\tau _b)}{(S_1S_2)}\right)$$
(55)
The cooling efficiency of the refrigerator will be:
$$\eta _{Rf}=\frac{𝒬_{DC}}{𝒲_{cycle}^{on}}=\frac{\omega _a(S_2S_1)}{\left((\omega _b\omega _a)(S_2S_1)+\sigma ^2\omega _b(1/\tau _a+1/\tau _b)\right)}$$
(56)
or:
$$\frac{1}{\eta _{Rf}}+1=\frac{1}{COP}+1=\frac{\omega _b}{\omega _a}\left(1+\frac{\sigma ^2(1/\tau _a+1/\tau _b)}{(S_2S_1)}\right)>\frac{T_h}{T_c}\left(1+\frac{\sigma ^2(1/\tau _a+1/\tau _b)}{(S_2S_1)}\right),$$
(57)
leading to the expression for the efficiency:
$$\eta _{Rf}=\frac{\omega _a}{\omega _b}\frac{1}{\eta _{H.E.}^{fricles}+\frac{\sigma ^2(1/\tau _a+1/\tau _b)}{S_2S_1}}<\frac{T_c}{T_h}\frac{1}{\eta _{H.E.}^{fricles}+\frac{\sigma ^2(1/\tau _a+1/\tau _b)}{S_2S_1}}$$
(58)
For both the heat engine and the heat pump, the efficiency is explicitly dependent on time allocation, cycle time, and bath temperatures.
#### 2 Optimization
The performance of both the heat engine and the heat pump can be optimized with respect to:
* (a) The overall time period $`\tau `$ of the cycle, and its allocation between the hot and cold branches.
* (b) The overall optimal time allocation between all branches. (This optimization is performed only for the heat pump.)
* (c) The external fields, ($`\omega _a`$, $`\omega _b`$).
(a) Optimization with respect to time allocation.
The optimization of time allocation is carried out with the constant fields $`\omega _a`$ and $`\omega _b`$. The Lagrangian for the work output becomes:
$`(x,y,\lambda )=𝒲_{cycle}+\lambda \left(\tau +{\displaystyle \frac{1}{\mathrm{\Gamma }_c}}\mathrm{ln}(x)+{\displaystyle \frac{1}{\mathrm{\Gamma }_h}}\mathrm{ln}(y)\tau _a\tau _b\right),`$ (59)
where $`\lambda `$ is a Lagrange multiplier. Equating the partial derivatives of $`(x,y,\lambda )`$ with respect to $`x`$ and $`y`$ to zero, the following condition for the optimal time allocation becomes:
$`\begin{array}{c}\mathrm{\Gamma }_cx\left((1y)^2(S_h^{eq}S_c^{eq})+\sigma ^2(1y)(1/\tau _a+y/\tau _b)\right)=\hfill \\ \mathrm{\Gamma }_hy\left((1x)^2(S_h^{eq}S_c^{eq})\sigma ^2(1x)(x/\tau _a+1/\tau _b)\right)\hfill \end{array}`$ (62)
When $`\sigma =0`$, the previous frictionless result is retrieved. (Optimizing the entropy production $`\mathrm{\Delta }S^u`$ leads to an identical time allocation to Eq. (62)).
Eq. (62) can also be written in the following way:
$`\mathrm{\Gamma }_cx\left((1y)(1yx_{max})\right)=\mathrm{\Gamma }_hy\left((1x)(x_{max}x)\right)`$ (63)
where $`x_{max}`$ was defined in Eq. (38). The result is dependent on the time allocations of the ’adiabats’, through the dependence of $`x_{max}`$.
For the special case when $`\mathrm{\Gamma }_c=\mathrm{\Gamma }_h`$, the relation between the time allocations in contact with the hot and cold baths becomes:
$`x=x_{max}y`$ (64)
For the frictionless case, this result coincides with the former frictionless one$`x=y`$, meaning that equal time is allocated to contact with the cold and hot reservoirs. When friction is added this symmetry is broken, Eq. (64), to compensate for the additional heat generated by friction the time allocated to the cold branch, becomes larger than the time on the hot branch.
The Lagrangian for the heat-flow, $`𝒬_F`$, extracted from the cold reservoir is defined in parallel to the Lagrangian for the total work. Substituting $`\mathrm{\Gamma }_h`$ for $`\mathrm{\Gamma }_c`$, $`x`$ for $`y`$ and vice versa, also $`y_{max}`$ for $`x_{max}`$, where $`y_{max}`$ was defined in Eq. (44), one gets the optimal time allocation for the heat pump.
Optimization of power with respect to time allocation as a function of the cycle time, $`\tau `$ for different friction coefficients is shown in Fig. 6 (Left), together with the corresponding heat-flow (Middle) and the corresponding entropy production (Right). The left part shows that in the frictionless case the power obtains its maximum at zero cycle time with the value consistent with Eq. (71). When friction is introduced, the maximum power decreases and is shifted to longer cycle times. The figure also shows, that for short times the work done by the system is negative, and as the friction coefficient $`\sigma `$ increases, the boundary between positive and negative power shifts to longer cycle times. In the Middle of Fig. 6, the heat-flow corresponding to the optimal power on the left is shown. The shapes of the power and heat flow curves are similar. The heat-flow values are always positive and larger than the corresponding power values. The entropy production (Right) shows that unlike the power curves the friction changes significantly the shape of the curves. The entropy production rate for the case with friction sharply decreases. The parallel graphs for the heat pump are similar.
(b) Time allocation optimization between all branches of the refrigerator
Further optimization of the performance of the heat pump is possible by relaxing the assumption of constant time on the ’adiabats’. First the time allocation between the two ’adiabats’ is optimized, when $`\tau _a+\tau _b=\delta `$, where $`\delta `$ is a constant. Finally the time allocation between the ’adiabats’ and the heat exchange branches, is optimized. These results are compared to the recent analysis of Gordon et. all. .
From Eqs. (53) and (43) with constant time allocations along the heat exchange branches one gets for the cooling power:
$`𝒬_F=A_0A_1({\displaystyle \frac{y}{\tau _a}}+{\displaystyle \frac{1}{(\delta \tau _a)}}),`$ (65)
where $`A_0`$ and $`A_1`$ are constant functions of the parameters of the system. And on $`\delta `$, a double inequality is imposed $`\tau >\delta >`$ the larger of \[ $`(\tau \tau _{h,min});\tau _{b,min}`$\], see Eq. (42).
The optimal $`\tau _a`$ depends only on $`y`$ and on $`\delta `$. The optimal value of $`\tau _{a,opt}`$ becomes:
$`\tau _{a,opt}=\delta {\displaystyle \frac{y+\sqrt{y}}{1y}}`$ (66)
Further optimization by changing the the value of $`\delta `$, changes the cycle time $`\tau `$. This optimization step is done by numerical iteration. Typically the sum of the final optimal values of $`\tau _a`$ and $`\tau _b`$ is about twice their value before, and their ratio is about 0.7 of the value which was chosen initially.
The next step is to study the time allocation between the ’adiabats’ and the heat exchange branches when all other controls of the heat pump have optimal values. These controls include also the external fields of optimization which are described later.
For comparison with Gordon et. all. , the results of optimization are plotted in the $`1/𝒬_F`$ ,$`1/\eta `$ plane for a fixed cycle time $`\tau `$. The following example demonstrates the method followed: First an optimal starting value for $`𝒬_F`$ was found which determines the time allocation control parameters, $`\tau _c=0.44221,\tau _h=0.31779,\tau _a=0.0084,\tau _b=0.0116`$ with a total cycle time of $`\tau =0.78`$. Under such conditions $`𝒬_{F,max}=2.9158`$ ($`1/𝒬_{F,max}=0.34296`$).
Changing the time allocation between the ’adiabats’ and the heat exchange branches changes the balance between optimal cooling power and efficiency. Denoting the sum $`\tau _c`$+$`\tau _h`$ by $`\tau _{ch}`$, the ratio $`\tau _h/\tau _c`$ by $`r_{hc}`$, the sum $`\tau _a`$ \+ $`\tau _b`$ by $`\tau _{ab}`$, the ratio $`\tau _a/\tau _b`$ by $`r_{ab}`$, time is transfered from $`\tau _{ch}`$ by small steps to $`\tau _{ab}`$, while keeping the the ratios $`r_{hc}`$ and $`r_{ab}`$ constant. For each step the corresponding $`1/𝒬_F`$ and $`1/\eta `$, are calculated as in Fig. 7. The relation between the reciprocal efficiency and the reciprocal cooling power shows the tradeoff between losses due to friction and losses due to heat transfer. Following the curve in Fig. 7, starting from point $`𝐀`$ where the cooling power is optimal, resources represented by time allocation are transferred from the heat exchange branches to the ’adiabats’, reducing the friction losses. At point $`𝐁`$ an optimum is reached for the efficiency. This point has been found by Gordon et. al. to be the universal operating choice for commercial chillers. Point $`𝐁`$ represents the optimal compromise between maximum efficiency and cooling power.
Point $`𝐀`$ is located at the maximum cooling power. If more time is allocated to the heat exchange branches both $`1/𝒬_F`$ and $`1/\eta `$ will continue to increase as seen in the insert of Fig. 7.
(c) Optimization with respect to the fields.
The values of the fields $`\omega _a`$ and $`\omega _b`$ are control parameters of the engine. In a spin system these fields are equivalent to the value of the external magnetic field applied on the system. They directly influence the energy spacing of the TLS. The work function $`𝒲_{cycle}`$, or equivalently the power ($`𝒫`$) is optimized with respect to the fields, subject to the Carnot constraint:
$$\frac{\omega _a}{T_c}\frac{\omega _b}{T_h}$$
(67)
Optimal power is obtained by equating independently to zero the partial derivatives of $`𝒲_{cycle}`$, or of $`𝒫=𝒲_{cycle}/\tau `$ by varying $`\omega _a`$ and $`\omega _b`$. In addition the optimal solutions have to fulfill the inequality constraints in Eq. (67). As a result two transcendental equations in $`\omega _a`$ and $`\omega _b`$ are obtained which are solved numerically.
The two equations are:
$`\begin{array}{c}\hfill \frac{(1yx_{max})}{(\omega _b\omega _a)}(\mathrm{\Delta }S^{eq}+\sigma ^2/\tau _a)\mathrm{cosh}^2\left(\frac{\omega _a}{2k_BT_c}\right)=\frac{1y}{(4k_BT_c)}\\ \hfill \frac{(x_{max}x)}{(\omega _b\omega _a)}(\mathrm{\Delta }S^{eq}+\sigma ^2/\tau _a)\mathrm{cosh}^2\left(\frac{\omega _b}{2k_BT_h}\right)=\frac{1x}{(4k_BT_h)}\end{array}`$ (70)
Where $`\mathrm{\Delta }S^{eq}`$= S$`{}_{}{}^{eq}{}_{h}{}^{}`$\- S$`{}_{}{}^{eq}{}_{c}{}^{}`$ as defined in Eq. (36). Examining Eq. (70), and fixing the friction $`\sigma `$, it is found that $`\mathrm{\Delta }S^{eq}`$ is an extensive function of order zero (intensive ) with respect to the quartet of variables $`\omega _a,T_c,\omega _b,T_h`$. This means that scaling these parameters simultaneously will not change $`\mathrm{\Delta }S^{eq}`$. Also x<sub>max</sub>, and $`\mathrm{cosh}^2\left(\frac{\omega }{2k_BT}\right)`$ are extensive (order zero). The work function however, is extensive with order one (Eqs. (51) and (52)). This property will be exploited in paragraph III.
The optimization of power with respect to the fields is shown in Fig. 8 for the frictionless engine, as a function of the fields with fixed time allocation. A global maximum can be identified.
The heat pump optimization of $`𝒬_F`$ with respect to the fields is different and therefore will be presented in Section III.
The analysis for the optimization with respect to the fields for the entropy production $`\mathrm{\Delta }𝒮^u`$, is presented in appendix C. The optimal solution without friction($`\sigma =0`$) leads to $`\mathrm{\Delta }𝒮_{min}^u=0`$. When $`\sigma 0`$, the minimum value of $`\mathrm{\Delta }𝒮^u`$ is different from zero, and is achieved on the boundary of the region.
### E Global Optimization of the Heat Engine
Global optimization of the power means searching for the optimimum with respect to the control parameters cycle time, time allocation and the fields. An iterative procedure is used.
The procedure is initiated by setting the optimal time allocation from the corresponding Lagrangian, with $`\sigma =0`$. The power becomes a product of two functions, one depending only on time the other only on the fields, and therefore, the fields can be changed independently of time. The optimal fields for the above time allocation are then sought. For the frictionless case, the overall time on the adiabats tends to zero. The optimal field values become independent of time. The value $`𝒫=107.501`$ is the short time limit in accordance with the equation:
$`𝒫(\omega _b\omega _a)(S_h^{eq}S_c^{eq})(\mathrm{\Gamma }_c\mathrm{\Gamma }_h)/(\sqrt{\mathrm{\Gamma }_c}+\sqrt{\mathrm{\Gamma }_h})^2`$ (71)
These fields are inserted into the expression with friction $`\sigma 0`$, and the new optimal times and fields are computed. The iteration converges after two to three steps, as indicated by Table IV for $`\sigma =0.005`$. Notice that the location of the optimum is not very sensitive to the friction parameter.
## III Asymptotic Properties of the Heat Pump when the cold bath temperature approaches absolute zero.
The goal is to obtain an asymptotic upper bound on the cooling power when the heat pump is operating close to absolute zero temperature. This requires optimizing the performance of the heat pump with respect to all control parameters.
### A Optimization of the heat-flow $`𝒬_F`$ with respect to the fields and to the cooling power upper bound.
The heat-flow,$`𝒬_F`$ extracted from the cold reservoir now becomes the subject of interest:
$`𝒬_F=\omega _a(S_2S_1)/\tau `$ (72)
or from Eq. (43),
$`𝒬_F=(\omega _a/\tau )\left((S_2^{eq}S_1^{eq})F(x,y){\displaystyle \frac{\sigma ^2(1x)(y/\tau _a+1/\tau _b)}{(1xy)}}\right)`$ (73)
No global maximum for the $`𝒬_F`$ with respect to the fields is found. The derivative of $`𝒬_F`$ with respect to $`\omega _b`$ becomes:
$`{\displaystyle \frac{𝒬_F}{\omega _b}}={\displaystyle \frac{F(x,y)\omega _a}{\tau }}{\displaystyle \frac{1}{4k_BT_h\mathrm{cosh}^2\frac{\omega _b}{2k_BT_h}}}0`$ (74)
leading to the result that $`𝒬_F`$ is monotonic in $`\omega _b`$. Under such conditions, $`\omega _b`$ is set, and the optimum with respect to $`\omega _a`$ is sought for. The derivative of $`𝒬_F`$ with respect to $`\omega _a`$ becomes:
$`{\displaystyle \frac{𝒬_F}{\omega _a}}=(S_2^{eq}S_1^{eq}){\displaystyle \frac{\sigma ^2}{(1y)}}(y/\tau _a+1/\tau _b)\omega _a{\displaystyle \frac{1}{4k_BT_c\mathrm{cosh}^2\frac{\omega _a}{2k_BT_c}}}=0`$ (75)
Introducing from Eq. (75) the optimal value of $`(S_2^{eq}S_1^{eq})\frac{\sigma ^2}{(1y)}(y/\tau _a+1/\tau _b)`$, into Eq. (73), leads to the optimal cooling rate:
$`𝒬_F^{optimum}={\displaystyle \frac{F(x,y)\omega _a^2}{\tau }}{\displaystyle \frac{1}{4k_BT_c\mathrm{cosh}^2\frac{\omega _a}{2k_BT_c}}}={\displaystyle \frac{F(x,y)}{4k_B\tau }}\left({\displaystyle \frac{\omega _a}{T_c}}\right)^2{\displaystyle \frac{T_c}{\mathrm{cosh}^2\frac{\omega _a}{2k_BT_c}}}`$ (76)
Due to its extensivity, the ratio $`\frac{\omega _a}{T_c}`$ becomes a constant, while both $`\omega _a`$ and T<sub>c</sub> can approach zero.
From Eq.(76), an upper-bound for the cooling rate $`𝒬_F`$ is obtained:
$`𝒬_F^{optimum}{\displaystyle \frac{F(x,y)}{4k_B\tau }}\left({\displaystyle \frac{\omega _a}{T_c}}\right)^2T_c.`$ (77)
From Eqs. (77), when T<sub>c</sub> approaches zero, the cooling rate vanishes, at least linearly with temperature. This is a third law statement which shows that absolute zero cannot be reached since the rate of cooling vanishes as absolute zero is approached.
### B The asymptotic relation between the internal and external temperature on the cold branch
When the bath temperature tends to zero, the internal working fluid temperature has to follow. This becomes a linear relationship between $`T^{}`$ and $`T_c`$ as $`T_c`$ tends to zero.
Calculating the polarization at the end of the contact with the cold bath $`S_2`$:
$`S_2=S_2^{eq}{\displaystyle \frac{(S_2^{eq}S_1^{eq})x(1y)\sigma ^2x(1/\tau _b+y/\tau _a)}{(1xy)}}`$ (78)
Assuming the relation $`T_h`$ = $`\rho `$ $`T_c`$ as $`T_c`$ tends to zero, the exponents can be expanded to the first order to give:
$`S_2={\displaystyle \frac{T_c}{\omega _a}}{\displaystyle \frac{1xy_{max}+(\omega _a/\omega _b)\rho x(y_{max}y)}{(1xy)}}+1/2{\displaystyle \frac{x(\sigma ^2/\tau _a)(y_{max}y)}{(1xy)}}`$ (79)
Also, $`S_2`$ defines the internal temperature $`T^{}`$ through the relation: $`S_2=\frac{1}{2}\mathrm{tanh}\left(\frac{\omega _a}{2k_BT^{}}\right)`$. Expanding the hyperbolic tangent, one gets:
$`T^{}=T_c{\displaystyle \frac{1xy_{max}+\rho (\omega _a/\omega _b)x(y_{max}y)}{(1xy)}}{\displaystyle \frac{x\omega _a(\sigma ^2/\tau _a)(y_{max}y)}{(1xy)}}`$ (80)
proving that $`T_c`$ and $`T^{}`$ both tend asymptotically to zero. It should be noted that the term independent of $`T_c`$ depends on $`\omega _a`$, which also tends to zero as $`T_c`$ tends to zero ( (Eq. 77). Eq. (77) also shows that $`𝒬_F^{optimum}`$\*T<sub>c</sub> is a quadratic function of $`\omega _a`$, Cf. Fig. 9.
Eq. (77) represents an upper-bound to the rate of cooling. In order to determine how closely this limit be approached, a strategy of cooling must be devised, which re-optimizes the cooling power during the changing conditions when T<sub>c</sub> approaches zero.
### C Optimal cooling strategy
The goal is to follow an optimal cooling strategy, which exploits the properties of the equations and achieves the upper-bound for the rate of cooling, $`𝒬_F`$.
The properties of the equations employed are;
* i: The derivative with respect to $`\omega _a`$ of $`𝒬_F`$ ( Eq. (75)), is extensive of order zero in the ’quartet’ ($`\omega _a,\omega _b,T_c,T_h`$).
* ii: For $`\frac{𝒬_F}{\omega _a}`$ the extensivity holds also for the ’doublets’ ($`\omega _b,T_h`$) or ($`\omega _a,T_c`$). Scaling these variables by the same number, leaves Eq. (75) equal to zero, and the value of $`𝒬_F^{optimum}`$ does not change.
* iii: In spite of $`𝒬_F`$ being monotonic in $`\omega _b`$, $`𝒬_F^{optimum}`$ is independent of $`\omega _b`$ (and of T<sub>h</sub>), therefore $`𝒬_F`$ saturates as $`\omega _b`$ is increased.
From property (i) it follows, that once an optimal ’quartet’($`\omega _a,\omega _b,T_c,T_h`$) is created, it is possible to cool optimally with a set of quartets, which are scaled by a decreasing set $`r_n<1`$, $`lim_n\mathrm{}r_n=0`$. For this set the limit of the ratio $`\frac{\omega _a}{T_c}`$ is a non zero constant. Therefore in Eq. (76) $`\omega _a`$ and T<sub>C</sub> are optimal leading to:
$`𝒬_F^{optimum}={\displaystyle \frac{F(x,y)}{4k_B\tau _{optimal}}}\left({\displaystyle \frac{\omega _{a,optimal}}{T_{c,optimal}}}\right)^2{\displaystyle \frac{T_{c,optimal}}{\mathrm{cosh}^2\frac{\omega _{a,optimal}}{2k_BT_{c,optimal}}}}`$ (81)
In general, the hot bath temperature is constant, and the property (ii) is used to scale back the value of the optimal T<sub>h</sub> to the bath temperature. As a result, the optimal high field is also scaled.
Property (iii) will be exploited by changing only $`\omega _b`$ in the optimal quartets and checking for saturation. See Fig. 13 and the dashed curves of Fig. 9. Summarizing, for every ’quartet’ the upper-bound in Eq. (77) can be reached. The details of the cooling strategy can be found in Appendix E
Fig. 10 shows that the cooling strategy ( Tables V and VI ) can approach the upper bound leading to a linear relation of the optimal cooling power with temperature. With respect to the fields the optimal strategy leads to a decrease of the field $`\omega _a`$ which is in contact with the cold bath. This causes the internal temperature of the TLS T’ to be lower than the cold bath temperature T<sub>c</sub>. On the hot side the optimal solution requires as large an energy separation as possible $`\omega _a\mathrm{}`$ but this effect saturates.
The linear relation of the cooling rate with T<sub>c</sub> leads to a constant asymptotic entropy production as can be seen in the right of Fig. 10 ( Cf. Appendix C).
## IV Conclusion
The detailed study of the four stroke discrete heat engine with internal friction serves as a source of insight on the performance of refrigerators at temperatures which are very close to absolute zero. The next step is to find out if the behavior of the specific heat pump described in the study can be generalized. A comparison with other systems studied indicates that the conclusions drawn from the model are generic. As a heat engine the model shows the generic behavior of maximum power as a function of control parameters found in finite time thermodynamics . This is despite the fact that the heat transfer laws in the microscopic model of the working fluid are different from the macroscopic laws such as the Newtonian heat transfer law . When operated as a heat pump with friction, the present model shows the universal behavior observed for commercial chillers caused by a tradeoff between allocating resources to the ’adiabats’ or to the heat exchange branches.
Another question is whether the linear scaling of the optimal cooling power at low cold bath temperatures is a universal phenomenon. For low temperatures the results of the present model can be extended to a working fluid consisting of an ensemble of harmonic oscillators or any N-level systems. This is because at the limit of absolute zero only the two lowest energy levels are relevant. When examined, other models with different operating cycles show an identical behavior. For example the continuous model of a quantum heat engine based on reversing the operation of a laser shows this linear scaling phenomena. Another example is the Ericsson refrigeration cycle Cf. Eq. (23) in the study of Chen et al which shows the same asymptotic linear relationship.
A point of concern is the dependence of the heat transfer laws on temperature when absolute zero is approached. The kinetic parameters $`k_{}`$ and $`k_{}`$ represent an individual coupling of the two-level-system to the bath. Considering coefficients derived from gas phase collisions they settle to a constant asymptotic value as the temperature is lowered . The reason is that the slow approach velocity is compensated by the increase in the thermal De-Broglie wavelength.
There has been an ongoing interest in the meaning of the third law of thermodynamics . The issue at stake has been: is the third law an independent postulate or it is a consequence of the second law and the vanishing of the heat capacity. This study presents a dynamical interpretation of the third law. The absolute temperature cannot be reached because the maximum rate of cooling vanishes linearly at least with temperature.
###### Acknowledgements.
This research was supported by the US Navy under contract number N00014-91-J-1498. The authors want to thank Jeff Gordon for his continuous help, discussions and willingness to clarify many fine points. T.F. thanks Sylvio May for his help.
## A Analysis for the ’moving’ Cycles.
Insight into the origin of the behavior of the ’moving’ cycles is seen in Fig. 11, where the polarizations S<sub>1</sub>, S<sub>2</sub> are shown as monotonically decreasing functions of the time allocation on the cold bath. However, the envelope of S<sub>1</sub> for maximal power, namely for maximal S<sub>1</sub>-S<sub>2</sub> is worth noticing. It is a decreasing function for short cycle times, achieves a minimum at $`\tau _0`$, and starts to increase for $`\tau >\tau _0`$. Thus it is responsible for shifting the cycles to smaller polarization for short cycle times, and for the change of that trend for larger cycle times. The envelope of S<sub>2</sub> for maximal S<sub>1</sub>-S<sub>2</sub> is also a monotonically decreasing function of $`\tau _c`$, or equivalently of $`\tau `$, supporting the increase of S<sub>1</sub>-S<sub>2</sub>. The figure also shows, that for a short time allocation both S<sub>1</sub> and S<sub>2</sub> are close to the equilibrium polarization S$`{}_{}{}^{eq}{}_{h}{}^{}`$, When not enough time is allocated on the hot bath both the polarizations S<sub>1</sub> and S<sub>2</sub> approach S$`{}_{}{}^{eq}{}_{c}{}^{}`$.
## B The computation of $`\tau _0`$
The computation of $`\tau _0`$ Eq. (48) is not sufficient since it gives only the relation between the times spent on the cold and hot branches for zero work. The natural additional requirement is to seek for the optimal allocations, $`\tau _{c,0}`$ and $`\tau _{h,0}`$ using Eq. (63): $`\tau _0`$ = $`\tau _{c,0}`$ \+ $`\tau _{h,0}`$ \+ $`\tau _a`$ \+ $`\tau _b`$
Denoting by $`x_0`$ and $`y_0`$ the corresponding x and y values defined in Eq. (21), the following two equations for x<sub>0</sub> and y<sub>0</sub> are obtained:
$`y_0={\displaystyle \frac{(x_{max}x_0)R}{(x_{max}x_0)Rx_0}}`$ (B1)
and
$`\mathrm{\Gamma }_cx_0\left((1y_0)(1y_0x_{max})\right)=\mathrm{\Gamma }_hy_0\left((1x_0)(x_{max}x_0)\right)`$ (B2)
Where R is defined as:
$`R={\displaystyle \frac{\sigma ^2\omega _a(1/\tau _a+1/\tau _b)}{(\omega _b\omega _a)(S_h^{eq}S_c^{eq}+\sigma ^2/\tau _a)}}`$ (B3)
and x<sub>max</sub> was defined in Eq. (38) as:
$`x_{max}={\displaystyle \frac{(S_h^{eq}S_c^{eq})\sigma ^2/\tau _b}{(S_h^{eq}S_c^{eq})+\sigma ^2/\tau _a}}`$ (B4)
The quadratic equation to be solved for x<sub>0</sub> is,
$`AAx_0^2+BBx_0+CC=0`$ (B5)
Where AA = $`\mathrm{\Gamma }_h`$ (1 + R) BB = - ($`\mathrm{\Gamma }_h`$ ((1 + R) (x<sub>max</sub>-R) + x<sub>max</sub>) \+ $`\mathrm{\Gamma }_c`$ (1 + R - $`x_{max}`$)) and CC = $`\mathrm{\Gamma }_h`$ (x<sub>max</sub>-R) x<sub>max</sub>
## C Entropy production.
(1) Heat Engine.
$`\mathrm{\Delta }𝒮_{cyle1}^u=(𝒬_{AB}/T_h+𝒬_{CD}/T_c)`$ (C1)
Or from Table (I)
$`\mathrm{\Delta }𝒮_{cyle1}^u=(\omega _a/T_c\omega _b/T_h)(S_1S_2)+{\displaystyle \frac{\sigma ^2\omega _a}{T_c}}(1/\tau _a+1/\tau _b)`$ (C2)
The entropy production results are shown in Fig. 12. The left figure shows $`\mathrm{\Delta }𝒮^u`$ with increasing friction. The middle figure shows the corresponding cycles, while the right figure shows the corresponding power values.
The reciprocal behavior of the entropy production and the power is clear from Fig. 12. One also observes, that for the given cycle time the ’free’ time for the cycles with increasing $`\sigma `$ becomes more restricted. This follows from the dependence of $`\tau _{c,min}`$ on $`\sigma `$. See also Fig.( 3)
Introducing Eq. (36) into Eq. (C2). The entropy production becomes,
$`\mathrm{\Delta }𝒮_{cyle1}^u=(\omega _a/T_c\omega _b/T_h)(S_h^{eq}S_c^{eq})F(x,y)+\mathrm{\Delta }𝒮_{\sigma 1}^u`$ (C3)
where
$`\mathrm{\Delta }𝒮_{\sigma 1}^u=\sigma ^2{\displaystyle \frac{1}{(1xy)}}\left({\displaystyle \frac{\omega _a}{T_c}}(1x)(1/\tau _a+y/\tau _b)+{\displaystyle \frac{\omega _b}{T_h}}(1y)(x/\tau _a+1/\tau _b)\right)`$ (C4)
Notice, that $`\mathrm{\Delta }𝒮_{\sigma 1}^u`$ is always positive. For $`\sigma =0`$ Eq. (C4) reduces to the frictionless results .
(2) Heat Pump
The entropy production for the heat pump becomes:
$`\mathrm{\Delta }𝒮_{}^{u}{}_{ref}{}^{}=\left({\displaystyle \frac{\omega _b}{T_h}}{\displaystyle \frac{\omega _a}{T_c}}\right)(S_2S_1)+\sigma ^2{\displaystyle \frac{\omega _b}{T_h}}(1/\tau _a+1/\tau _b)`$ (C5)
$`=\left({\displaystyle \frac{\omega _b}{T_h}}{\displaystyle \frac{\omega _a}{T_c}}\right)(S_2^{eq}S_1^{eq})F(x,y)+`$ (C6)
$`\sigma ^2F(x,y)\left\{{\displaystyle \frac{\omega _b}{T_h}}{\displaystyle \frac{1}{1x}}({\displaystyle \frac{1}{\tau _a}}+{\displaystyle \frac{x}{\tau _b}})+{\displaystyle \frac{\omega _a}{T_c}}{\displaystyle \frac{1}{1y}}({\displaystyle \frac{1}{\tau _b}}+{\displaystyle \frac{y}{\tau _a}})\right\}`$ (C7)
The asymptotic entropy production as T<sub>c</sub> tends to zero can be calculated leading to
$`\mathrm{\Delta }𝒮_{ref}^u=F(x,y)[(\omega _b/(\rho \omega _a))(1\rho (\omega _a/\omega _b))^2+`$ (C8)
$`\sigma ^2({\displaystyle \frac{\omega _b}{\rho T_c}}{\displaystyle \frac{1}{(1x)}}(1/\tau _a+x/\tau _b)+{\displaystyle \frac{\omega _a}{T_c}}{\displaystyle \frac{1}{(1y)}}(1/\tau _b+y/\tau _a))]`$ (C9)
Since $`T_h=\rho T_c`$, the r.h.s. of Eq. (C9) tends to a constant, for each term depends on the constant ratios ($`\omega _b`$/T<sub>h</sub>), ($`\omega _a`$/T<sub>c</sub>) or on their ratio. This result is demonstrated on the right side of Fig. 10.
The optimization with respect to time allocation has the same result as for the heat engine. Therefore, only optimization with respect to the fields are presented;
Equating to zero the derivatives with respect to x an y of the entropy production, one gets two similar equation to the total work derivatives:
$`{\displaystyle \frac{(1yx_{max})}{(\omega _a/T_c\omega _bT_h)}}(\mathrm{\Delta }S^{eq}+\sigma ^2/\tau _a)\mathrm{cosh}^2\left({\displaystyle \frac{\omega _a}{2k_BT_c}}\right)+{\displaystyle \frac{1y}{(4k_BT_c)}}0`$ (C10)
$`{\displaystyle \frac{(x_{max}x)}{(\omega _a/T_c\omega _bT_h)}}(\mathrm{\Delta }S^{eq}+\sigma ^2/\tau _a)\mathrm{cosh}^2\left({\displaystyle \frac{\omega _b}{2k_BT_h}}\right)+{\displaystyle \frac{1x}{(4k_BT_h)}}0`$ (C11)
Where $`\mathrm{\Delta }S^{eq}`$, is $`S_h^{eq}`$\- $`S_c^{eq}`$.
Eqs. (C10) and (C11) show that the entropy production is a monotonic function in the allowed range, namely, for
$`{\displaystyle \frac{\omega _a}{T_c}}>{\displaystyle \frac{\omega _b}{T_h}}.`$ (C12)
To conclude the entropy production has a minimum value: $`\mathrm{\Delta }𝒮_{min}^u`$, will be
$`\mathrm{\Delta }𝒮_{min}^u={\displaystyle \frac{\omega _a}{T_c}}\sigma ^2(1/\tau _a+1/\tau _b)`$ (C13)
obtained on the boundary of the range.
## D The Total Work done on the System for the Heat Pump
The total work done on the system becomes,
$`𝒲_{cyle3}^{on}=(\omega _b\omega _a)(S_2S_1)+\sigma ^2\omega _b(1/\tau _a+1/\tau _b)`$ (D1)
or
$`𝒲_{cyle3}^{on}=(\omega _b\omega _a)(S_2^{eq}S_1^{eq})F(x,y)+W_{\sigma 3}`$ (D2)
where
$`𝒲_{\sigma 3}={\displaystyle \frac{\sigma ^2}{(1xy)}}\left(\omega _b(1y)(1/\tau _a+x/\tau _b)+\omega _a(1x)(y/\tau _a+1/\tau _b)\right)`$ (D3)
$`=\sigma ^2F(x,y)\left\{{\displaystyle \frac{\omega _b}{1x}}({\displaystyle \frac{1}{\tau _a}}+{\displaystyle \frac{x}{\tau _b}})+{\displaystyle \frac{\omega _a}{1y}}({\displaystyle \frac{1}{\tau _b}}+{\displaystyle \frac{y}{\tau _a}})\right\}`$ (D4)
Eq. (D1) can be interpreted as the work done on the working fluid see (Cf. Fig. 2), as the sum of three positive areas, $`(\omega _b\omega _a)(S_2S_1)`$, $`\sigma ^2\omega _b(1/\tau _a)`$ and $`\sigma ^2\omega _b(1/\tau _b)`$ with the corresponding corners, D,C,B<sup>1</sup>,A<sup>1</sup>, B,B<sup>1</sup>,S<sub>2</sub>,S<sub>3</sub> and A<sup>1</sup>,A,S<sub>4</sub>,S<sub>1</sub>.
## E The optimal cooling strategy close to the absolute zero temperature
The first step in the cooling strategy is to create the first optimal quartet;
* (0) The systems external parameters $`\sigma `$, $`\tau _a`$, $`\tau _b`$, $`\mathrm{\Gamma }_c`$ and $`\mathrm{\Gamma }_h`$ are set.
* (1) A decreasing set of $`\omega _b`$ is chosen.
* (2) A constant ratio ($`\rho `$) for T<sub>h</sub>/T<sub>c</sub>, is chosen which is the ratio of the initial bath temperatures.
* (3) For the above chosen values, the optimal values of $`\omega _a`$, T<sub>c</sub>, $`\tau `$, and its optimal allocations between the branches to give maximal Q<sub>F</sub> are found for each $`\omega _b`$ in the set in (1) , by solving numerically the following additional equation to Eq. (75), with the condition that T$`{}_{h}{}^{}=\rho `$T<sub>c</sub>:
$`{\displaystyle \frac{𝒬_F}{T_c}}={\displaystyle \frac{F(x,y)}{4\tau k_B}}\left({\displaystyle \frac{\omega _a}{T_c}}\right)^2\left({\displaystyle \frac{1}{\mathrm{cosh}^2\frac{\omega _a}{2k_BT_c}}}{\displaystyle \frac{\omega _b}{\rho \omega _a\mathrm{cosh}^2\frac{\omega _b}{2\rho k_BT_c}}}\right)=0`$ (E1)
The above strategy causes the decrease of $`T_h`$ together with the $`T_c`$. Nevertheless according to (ii) above, the doublet $`\omega _b`$ and $`T_h`$ can be rescaled to increase $`T_h`$ back to its original value. The solid curves of Fig. 9 are optimal in the in the above described sense. Increasing only the value of $`\omega _b`$ in the optimal quartet according to point (iii), leads to larger values of the cooling rate, but eventually the increase of $`𝒬_F`$ will slow down and saturate. See Fig. 13 and the dashed curves of Fig. 9.
Fig. 13 represents the saturation phenomenon on $`\omega _b`$. Three points from Fig. 9 are chosen, and all parameters are fixed, except $`\omega _b`$, which is allowed to increase.
In order to approach the upper-bound for $`𝒬_F`$ in Eq. (77), a decreasing set of $`\omega _a/T_c`$ is created, achieved in an optimal way:
First step: After having an optimal ’quartet’, T<sub>c</sub> and T<sub>h</sub>, are fixed. Then, by lowering $`\omega _b`$, one finds the corresponding optimal $`\omega _a`$ values. This procedure is checked globally, by also iterating the time allocations. The results of a typical example are shown in Table V.
Second step: Using again the property of extensivity, the cooling will be achieved by multiplying the rows of Table V by a decreasing sequence, e.g. by $`2^n`$ for the n-th row. Table VI describes the cooling strategy, checking also the non-divergence of the entropy production both for the frictionless case and the case with friction. The results are also summarized in Fig. 10.
Table V demonstrates, that the procedure shifts down to the Carnot bound. The ratio R = $`\frac{\omega _b}{T_h}`$/ $`\frac{\omega _a}{T_c}`$. was computed showing only small changes. |
warning/0003/hep-th0003102.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Spontaneous anti de-Sitter compactifications of $`D=10`$ chiral $`N=2`$ supergravity have received some attention recently, in the light of the conjectured AdS/CFT correspondence . In particular $`AdS_5\times M_5`$ solutions , where $`M_5`$ = compact space, have been used to glean information on the related 4-dimensional conformal field theory . The presence of a (complex) three-form field strength can trigger compactifications on $`AdS_3\times M_7`$, an example being the known solutions $`M_7=S^3\times T^4`$ and $`M_7=S^3\times S^3\times S^1`$, whose corresponding two dimensional conformal theory has been investigated in . Further ref.s on IIB compactifications on spheres can be found in .
In this Letter we present a general class of $`AdS_3\times M_7`$ solutions with $`M_7`$ = compact 7-dimensional coset space $`G/H`$, and derive the conditions for these solutions to be supersymmetric backgrounds. This class of compactifications may be relevant for the study of the two-dimensional CFT’s examined in ref.s .
The theory contains a complex anti-Weyl gravitino $`\psi _M`$ and a complex Weyl spinor $`\lambda `$. The bosonic fields are: the graviton $`g_{MN}`$, a complex antisymmetric tensor $`A_{MN}`$, a real antisymmetric tensor $`A_{MNRS}`$ (restricted by a self-duality condition) and a complex scalar $`\varphi `$. There is a global $`U(1)`$ symmetry that rotates the two supersymmetry charges into each other. According to a general mechanism in supergravity theories, the scalars can be interpreted as coordinates of noncompact coset spaces. Here the complex scalar $`\varphi `$ parametrizes the coset $`SU(1,1)/U(1)`$.
After setting the spinor fields to zero, the field equations read :
$`2R_{MN}=P_MP_N^{}+P_M^{}P_N+{\displaystyle \frac{1}{6}}F_M^{PQRS}F_{PQRSN}+`$
$`+{\displaystyle \frac{1}{8}}(G_M^{PQ}G_{PQN}^{}+G_M^{PQ}G_{PQN}{\displaystyle \frac{1}{6}}g_{MN}G^{PQR}G_{PQR}^{})`$ (1.1)
$`F_{M_1M_5}={\displaystyle \frac{1}{5!}}ϵ_{M_1M_5N_1N_5}F^{N_1N_5}`$ (1.2)
$`(^SiQ^S)G_{MNS}=P^SG_{MNS}^{}{\displaystyle \frac{2i}{3}}F_{MNPQR}G^{PQR}`$ (1.3)
$`(^M2iQ^M)P_M={\displaystyle \frac{1}{24}}G^{PQR}G_{PQR}`$ (1.4)
where $`R_{MN}R_{MSN}^S`$ and the curvature two-form is defined as $`R_M^SdB_M^S+B_N^SB_M^N`$; the vectorial quantities $`P_M`$ (complex) and $`Q_M`$ (real) are related to the scalar fields (and derivatives thereof), $`F_{PQRSN}`$ and $`G_{PQN}`$ to the field strengths of the four-form and of the two-form, and $``$ is the Lorenz covariant derivative. Here we have adopted the normalizations and conventions of ; note however a sign correction in (1.4), already found in and noted also in . Moreover the following Bianchi identities hold (a consequence of the field definitions):
$`(_{[M}2iQ_{[M})P_{N]}=0,_{[M}Q_{N]}=iP_{[M}P_{N]}^{}`$ (1.5)
$`(_{[M}iQ_{[M})G_{NRS]}=P_{[M}G_{NRS]}^{}`$ (1.6)
$`_{[N}F_{M_1..M_5]}={\displaystyle \frac{1}{8}}\text{Im}G_{[NM_1M_2}G_{M_3M_4M_5]}^{}`$ (1.7)
The supersymmetry variations of the bosonic fields are proportional to Fermi fields, and these vanish in the type of backgrounds we are considering. On the other hand, the supersymmetry variations of the fermionic fields are:
$`\delta \lambda =i\mathrm{\Gamma }^M\epsilon ^{}P_M{\displaystyle \frac{1}{24}}iG_{MNP}\mathrm{\Gamma }^{MNP}\epsilon `$ (1.8)
$`\delta \psi _M=(_M{\displaystyle \frac{i}{2}}Q_M)\epsilon +{\displaystyle \frac{i}{480}}F_{N_1N_5}\mathrm{\Gamma }^{N_1N_5}\mathrm{\Gamma }_M\epsilon +`$
$`+{\displaystyle \frac{1}{96}}(\mathrm{\Gamma }_M^{N_1N_3}G_{N_1N_3}9\mathrm{\Gamma }^{N_1N_2}G_{MN_1N_2})\epsilon ^{}`$ (1.9)
(in backgrounds with $`\psi =0,\lambda =0`$) . A solution is supersymmetric if there exist spinors $`\epsilon `$ for which these variations vanish.
## 2 The Ansatz for the $`AdS_3\times M_7`$ solutions
We use the index conventions
M,N,P…= 1-10
m,n,p…= 1-3 (run on $`AdS_3`$)
a,b,c…= 4-10 (run on $`M_7`$).
and the flat “mostly minus” D=10 metric $`\eta =(+,,;,,,,,,)`$. With the Ansatz:
$`g_{MN}=\text{metric of }AdS_3\times M_7`$
$`\text{fermions}=0;P_M=Q_M=F_{M_1M_5}=0`$
$`G_{mnp}=eϵ_{mnp},G_{abc}=gJ_{abc}`$ (2.1)
where $`e`$ and $`g`$ are complex constants and $`J_{abc}`$ a real constant antisymmetric tensor, the field eq.s (1.4) take the form:
$$6e^2g^2J^2=0,J^2J_{abc}J_{abc}>0$$
(2.2)
implying
$$g=\rho \frac{\sqrt{6}}{J}e,\rho =\pm 1$$
(2.3)
Substituting the Ansatz(2.1)and the relation (2.3) into the remaining field eq.s yields:
$`R_{mn}={\displaystyle \frac{1}{4}}|e|^2g_{mn}`$ (2.4)
$`R_{ab}={\displaystyle \frac{3}{4}}{\displaystyle \frac{J_{ab}}{J^2}}|e|^2,J_{ab}J_{acd}J_{bcd}`$ (2.5)
$`^cJ_{abc}=0,=M_7\text{covariant derivative}`$ (2.6)
Eq. (1.2) is trivially satisfied, while eq. (1.3) with free indices $`m,n`$ holds because $`ϵ_{mnp}`$ is a covariantly conserved tensor in $`AdS_3`$. Moreover it is immediate to check that the Bianchi identities (1.5)-(1.7) are satisfied. Thus our Ansatz is a solution of the classical 2b supergravity equations provided eq.s (2.4) \- (2.6) hold. The first equation fixes the $`AdS_3`$ radius. If there exist an $`M_7`$ \- covariantly conserved three-index antisymmetric tensor $`J_{abc}`$ in $`M_7`$ the third equation is satisfied, and the second equation becomes a condition on the Ricci tensor of $`M_7`$.
As we show in the following, such $`J_{abc}`$ always exist in nonsymmetric coset spaces $`G/H`$ ( and in various symmetric $`G/H`$). Moreover $`J_{ab}`$ is diagonal, allowing in most cases $`G/H`$ to solve eq. (2.5) after a vielbein rescaling.
## 3 $`G/H`$ geometry and the tensor $`J_{abc}`$
The structure constants of $`𝔾=+𝕂`$ are defined by
$`[H_i,H_j]=C_{ij}^kH_kH_i`$
$`[H_i,K_a]=C_{ia}^jH_j+C_{ia}^bK_bK_a𝕂`$
$`[K_a,K_b]=C_{ab}^jH_j+C_{ab}^cK_c`$ (3.1)
where the index conventions are obvious. As discussed in ref. (p. 251), whenever $`H`$ is compact or semisimple one can always find a basis of $`K_a`$ such that the structure constants $`C_{ia}^j`$ vanish. In that case the $`𝔾=+𝕂`$ split, or equivalently the coset space $`G/H`$ is said to be reductive. For this reason we will deal in this paper only with reductive coset spaces. Another important observation is that when $`G/H`$ is reductive the structure constants $`C_{ia}^b`$ can always be made antisymmetric in $`a,b`$ by an appropriate redefinition $`K_aN_a^bK_b`$ .
For later use we recall the expression of the $`G/H`$ Riemannian connection
$$B_b^a=\frac{1}{2}\left(\frac{r_br_c}{r_a}C_{bc}^a+\frac{r_ar_c}{r_b}\eta _{bg}C_{dc}^g\eta ^{ad}+\frac{r_ar_b}{r_c}\eta _{cg}C_{db}^g\eta ^{ad}\right)V^cC_{bi}^a\mathrm{\Omega }^i$$
(3.2)
where $`V^c`$ and $`\mathrm{\Omega }^i`$ are the $`K`$ and $`H`$ vielbeins respectively, and we have allowed for rescalings $`r_c`$ of $`V^c`$ in the isotropy irreducible subspaces of $`K`$, see ref.s . These subspaces correspond to the block-diagonal pieces of the matrices $`C_{bi}^a`$, so that
$$\frac{r_a}{r_b}C_{ia}^b=C_{ia}^b$$
(3.3)
The $`G/H`$ Riemann curvature is defined by $`R_b^adB_b^a+B_c^aB_b^cR_{bde}^aV^dV^e`$ and reads :
$$R_{bde}^a=\frac{1}{4}\frac{r_dr_e}{r_c}𝑪_{bc}^aC_{de}^c+\frac{1}{2}r_dr_eC_{bi}^aC_{de}^i+\frac{1}{8}𝑪_{cd}^a𝑪_{be}^c\frac{1}{8}𝑪_{ce}^a𝑪_{bd}^c$$
(3.4)
with
$$𝑪_{bc}^a\frac{r_br_c}{r_a}C_{bc}^a\frac{r_ar_c}{r_b}C_{ac}^b\frac{r_ar_b}{r_c}C_{ab}^c$$
(3.5)
Consider now the field eq. (2.6), i.e.:
$$B_{e[a}^dJ_{bc]d}\eta ^{ec}=0$$
(3.6)
where the connection 1-form components are defined by $`B_a^dB_{ca}^dV^c`$. One possible choice for a $`J_{abc}`$ satisfying (3.6) is given by:
$$J_{abc}=C_{abc}C_{ab}^G\gamma _{cG},\text{G runs on the group }G\text{}\gamma \text{ = Killing metric}$$
(3.7)
Indeed $`B_{e[a}^dC_{bc]d}\eta ^{ec}=0`$ holds for the following reason. Observe that the left hand side is an $`H`$-invariant tensor, since the connection components, the structure constants $`C_{abc}`$ and the Killing metric are all $`H`$-invariant tensors. By $`H`$-invariant tensor we mean, for example, that:
$$\delta B_{cb}^aC_{ic}^dB_{db}^aC_{id}^aB_{cb}^d+C_{ib}^dB_{cd}^a=0$$
(3.8)
i.e. the adjoint action of $`H`$ on $`B`$ vanishes. It is not difficult to prove that in (3.2) the term multiplying $`V^c`$ is $`H`$-invariant. In fact each of the three terms within parentheses in (3.2) is $`H`$-invariant, as one can show by using Jacobi identities and (3.3).
In general the only $`H`$-invariant tensor with two free indices is the Killing metric (except in the special case of $`S^2`$, where one has also $`ϵ_{ab}`$), and therefore $`B_{e[a}^dC_{bc]d}\eta ^{ec}`$, being antisymmetric in its free indices, has to vanish.
In conclusion, the choice (3.7) satisfies the field eq.s (2.6). Moreover $`J_{ab}J_{acd}J_{bcd}`$, being a symmetric $`H`$-invariant tensor, must be proportional to the Killing metric in the $`H`$-isotropy subspaces.
The structure constants $`C_{abc}`$ are not the most general solution to eq. (3.6). For example also Antisymm ($`C_{ab}^c`$), i.e. the antisymmetrization of $`C_{ab}^c`$ on its three indices, satisfies eq. (3.6). In fact there may exist $`J_{abc}`$ tensors satisfying (3.6) even for symmetric $`G/H`$; this happens obviously for $`S^3=SO(4)/SO(3)`$, where $`J_{abc}`$ is proportional to $`ϵ_{abc}`$, or less trivially in the case $`G/H=S^3\times P^2`$ discussed in Section 7.
## 4 Supersymmetry conditions
We adopt the following real representation of the $`D=3+7`$ gamma matrices:
$$𝚪_M=(\gamma _m1\text{I}_{8\times 8}\sigma _2,1\text{I}_{2\times 2}\mathrm{\Gamma }_a\sigma _1)$$
(4.1)
where the $`SO(1,2)`$ gamma matrices are:
$$\gamma _1=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\gamma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\gamma _3=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)$$
(4.2)
and the real $`SO(7)`$ gamma matrices are given by the octonion structure constants (totally antisymmetric):
$`(\mathrm{\Gamma }_a)_{bc}=a_{abc},(\mathrm{\Gamma }_a)_{b8}=\delta _{ab}`$ (4.3)
$`a_{abc}:[123]=[165]=[257]=[354]=[367]=[246]=[147]=1`$ (4.4)
The supersymmetry parameter $`\epsilon `$ has the same Weyl chirality as the gravitino $`\psi `$, i.e. $`𝚪_{11}\epsilon =\epsilon `$ (anti-Weyl). Any anti-Weyl $`SO(1,9)`$ spinor can be decomposed as :
$$\epsilon =c_N\xi ^N\eta ^N\left(\begin{array}{c}0\\ 1\end{array}\right)$$
(4.5)
where $`\xi ^N`$ and $`\eta ^N`$ are real $`SO(1,2)`$ and $`SO(7)`$ spinors, respectively, and $`c_N`$. Substituting in the $`\lambda `$ supersymmetry condition:
$$G_{MNR}𝚪^{MNR}\epsilon =0(eϵ_{mnr}𝚪^{mnr}+gJ_{abc}𝚪^{abc})\epsilon =0$$
(4.6)
yields:
$$c_N\xi ^N(6e1\text{I}+gJ_{abc}\mathrm{\Gamma }^{abc})\eta ^N=0,$$
(4.7)
which can hold only if there exist $`SO(7)`$ spinors satisfying:
$$(1\text{I}+\rho \frac{J_{abc}}{\sqrt{6}J}\mathrm{\Gamma }^{abc})\eta =0$$
(4.8)
where we have also used (2.3).
Consider now the $`\psi `$ supersymmetry condition:
$`(d+{\displaystyle \frac{1}{4}}B^{mn}𝚪_{mn}+{\displaystyle \frac{1}{4}}B^{ab}𝚪_{ab})\epsilon +`$
$`+{\displaystyle \frac{1}{96}}V^m(G_{abc}𝚪_m𝚪^{abc}9eϵ_{mnr}𝚪^{nr})\epsilon ^{}+`$
$`+{\displaystyle \frac{1}{96}}V^d(G_{abc}𝚪_d^{abc}+eϵ_{mnr}𝚪_d𝚪^{mnr}9G_{dab}𝚪^{ab})\epsilon ^{}=0`$ (4.9)
$`V^m`$ and $`V^d`$ are the $`AdS_3`$ and $`M_7`$ vielbeins respectively. Substituting the $`𝚪`$-matrix and $`\epsilon `$ decomposition leads to:
$`c(_m\xi +{\displaystyle \frac{1}{4}}B_m^{rs}\gamma _{rs}\xi ){\displaystyle \frac{1}{4}}ic^{}e\gamma _m\xi =0`$ (4.10)
$`c(_c\eta +{\displaystyle \frac{1}{4}}B_c^{ab}\mathrm{\Gamma }_{ab}\eta ){\displaystyle \frac{1}{8}}c^{}G_{cab}\mathrm{\Gamma }^{ab}\eta =0`$ (4.11)
where we have dropped the index $`N`$ since the conditions can be satisfied only if they hold separately for every $`N`$. Moreover use of (4.8) has been necessary in order to achieve the factorization of (4.9) into (4.10) and (4.11). The modulus of $`c`$ is irrelevant in (4.10) and (4.11), and we can set $`c=\mathrm{exp}(i\phi )`$. Then the integrability condition for (4.10) is:
$$[\frac{1}{4}R_{rs}^{mn}\gamma _{mn}(\frac{1}{4})^2\mathrm{exp}(4i\phi )e^2\gamma _{rs}]\xi =0$$
(4.12)
The field equations (2.4) tell us that the $`AdS_3`$ curvature is:
$$R_{rs}^{mn}=\frac{1}{4}|e|^2\delta _{rs}^{mn}$$
(4.13)
so that the integrability condition can be satisfied if and only if the phase of $`c`$ is such that:
$$|e|^2=e^2\mathrm{exp}(4i\phi )e=\alpha |e|\mathrm{exp}(2i\phi ),\alpha =\pm 1$$
(4.14)
cf. ref. . Using this relation and (2.3) into (4.11) yields finally:
$`_m\xi +{\displaystyle \frac{1}{4}}B_m^{rs}\gamma _{rs}\xi {\displaystyle \frac{1}{4}}i\alpha |e|\gamma _m\xi =0`$ (4.15)
$`_c\eta +{\displaystyle \frac{1}{4}}B_c^{ab}\mathrm{\Gamma }_{ab}\eta {\displaystyle \frac{\sqrt{6}}{8}}\alpha \rho |e|{\displaystyle \frac{J_{cab}}{J}}\mathrm{\Gamma }^{ab}\eta =0`$ (4.16)
The integrability condition for (4.16) reads:
$$(R_{cd}^{ab}+\frac{3}{2}|e|^2\frac{J_{cae}J_{dbe}}{J^2})\mathrm{\Gamma }_{ab}\eta =0$$
(4.17)
The $`AdS_3\times M_7`$ solution preserves $`N`$ supersymmetries if and only if there exist $`N`$ independent $`SO(7)`$ spinors $`\eta `$ satisfying simultaneously eq.s (4.8) and (4.16).
In next Section we test our formulae in the case of the two known supersymmetric solutions, corresponding to $`G/H=S^3\times T^4`$ and $`G/H=S^3\times S^3\times S^1`$.
## 5 The supersymmetric solutions $`AdS_3\times S^3\times T^4`$ and $`AdS_3\times S^3\times S^3\times S^1`$
We’ll treat the two solutions simultaneously. The $`J_{abc}`$ tensors are simply the $`ϵ_{abc}`$ Levi-Civita tensors in the $`S^3`$ directions, so that Einstein field equations (2.5) are respectively:
$`R_{ab}={\displaystyle \frac{1}{4}}|e|^2\delta _{ab},a,b=1,2,3`$
$`R_{ab}=0,a,b=4,5,6,7`$ (5.1)
and
$`R_{ab}={\displaystyle \frac{1}{8}}|e|^2\delta _{ab},a,b=1,2,\mathrm{..6}`$
$`R_{ab}=0,a,b=7`$ (5.2)
fixing the radii of the $`S^3`$ spheres.
In the real gamma matrix representation of Section 4, the $`\delta \lambda =0`$ supersymmetry condition is satisfied by any linear combination of the four independent spinors $`\eta _1,\eta _2,\eta _3,\eta _8`$, the 8-dimensional spinor $`\eta _a`$ having the a-th component as only nonvanishing component. This holds for both solutions. On the other hand, the $`\delta \psi =0`$ supersymmetry condition is satisfied by 8 independent real spinors $`\eta _1^\pm ,\eta _2^\pm ,\eta _3^\pm ,\eta _8^\pm `$, the <sup>±</sup> referring to the sign $`\alpha `$ of (4.16). These spinors depend on the $`M_7`$ coordinates. Again this holds for both solutions. The two compactifications have then $`N=8`$ supersymmetries, or 16 real conserved supercharges (since $`\xi `$ has two real components).
## 6 Other $`AdS_3\times G/H`$ solutions
A complete classification of all $`AdS_3\times G/H`$ solutions based on the Ansatz (2.1) is postponed to a later publication. This requires to find the most general $`J_{abc}`$ that solves eq. (3.6) for each 7-dimensional $`G/H`$.
Here we give some selected examples, choosing some particular $`J`$ tensors. In fact, all the 7-dimensional $`G/H`$ cosets classified in are solutions of the IIB field equations, after suitable rescalings of the coset vielbeins. But the $`G/H`$ list of IIB solutions is actually larger than the one relevant for $`D=11`$ supergravity compactified on $`AdS_4\times G/H`$. Indeed in the IIB case the field equations do not force $`G/H`$ to be an Einstein space, but rather to be “isotropy Einstein”, i.e. with a Ricci tensor $`R_{ab}`$ proportional to $`\delta _{ab}`$ in each isotropy irreducible subspace of $`G/H`$. The proportionality constant can also vanish: when this happens in one-dimensional subspaces $`S^1`$ factors are allowed (they were excluded in the list of ).
### 6.1 $`AdS_3\times M^{pqr}`$
The $`M^{pqr}`$ spaces have been studied in detail in ref.s . They have three isotropy irreducible subspaces, allowing three independent rescalings $`a,b,c`$ of the vielbeins corresponding to the (1,2), 3, (4,5,6,7) directions and preserving the $`SU(3)\times SU(2)\times U(1)`$ isometry. Their Ricci tensor is given by :
$`R_{mn}={\displaystyle \frac{1}{4}}b^2(2{\displaystyle \frac{b^2}{c^2}}q^2)\delta _{mn},m,n=1,2`$
$`R_{33}={\displaystyle \frac{9a^4p^2+2b^4q^2}{8c^2}}`$
$`R_{AB}={\displaystyle \frac{3}{16}}a^2(43{\displaystyle \frac{a^2}{c^2}}p^2)\delta _{AB},A,B=4,5,6,7`$ (6.1)
For $`q0`$, $`p0`$ we can redefine:
$$a=\frac{q}{p}\gamma \sqrt{\frac{2\alpha }{3}},b=\gamma \sqrt{2\beta },c=q\gamma $$
(6.2)
so that the Ricci tensor becomes:
$`R_{mn}=\gamma ^2\beta (1\beta )\delta _{mn},m,n=1,2`$
$`R_{33}=\gamma ^2(\beta ^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{q^2}{p^2}}\alpha ^2)`$
$`R_{AB}={\displaystyle \frac{1}{2}}\gamma ^2\alpha (1{\displaystyle \frac{1}{2}}\alpha ){\displaystyle \frac{q^2}{p^2}}\delta _{AB}`$ (6.3)
An explicit check reveals that taking $`J_{abc}`$ to be the Levi-Civita tensor in the directions 1,2,3 (as in the case of the $`S^3\times T^4`$ solution) and otherwise zero satisfies the condition (3.6). Then the field equations are as in (5.1), and are satisfied by the Ricci tensor in (6.3) when the rescalings are:
$$\alpha =2,\beta =\frac{1\pm \sqrt{116\frac{q^2}{p^2}}}{4},\gamma =\frac{1}{4\beta (1\beta )}|e|^2$$
(6.4)
requiring $`p/q4`$. The particular case $`q=0`$, corresponding to $`S^2\times S^5`$, is also a solution: the rescalings are then:
$$a^2=\frac{1}{6}|e|^2,b^2=\frac{1}{2}|e|^2,c^2=\frac{1}{8}p^2|e|^2$$
(6.5)
The $`\delta \lambda `$ supersymmetry condition is satisfied by the same spinors $`\eta _1,\eta _2,\eta _3,\eta _8`$ discussed in the previous Section. However these spinors do not satisfy the $`\delta \psi `$ condition, and therefore these solutions are not supersymmetric.
There are other possible choices for $`J_{abc}`$. For example $`J_{abc}=`$ Antisymm($`C_{ab}^c`$) leads to field equations that can still be solved by a set of different rescalings $`\alpha ,\beta ,\gamma `$. In this case the four directions 4,5,6,7 are not Ricci flat.
### 6.2 $`AdS_3\times N^{010}`$
The $`N^{010}`$ coset spaces are a special case in the class of $`N^{pqr}`$ coset spaces studied in ref.s in the context of $`D=11`$ supergravity compactifications. They can be realized as the quotient:
$$N^{010}=\frac{SU(3)\times SU(2)}{SU(2)\times U(1)}$$
(6.6)
where the $`SU(2)`$ in the denominator is diagonally embedded in $`G=SU(3)\times SU(2)`$. In this formulation the full isometry of $`N^{010}`$ comes from the left action of $`G`$ . The $`N^{010}`$ geometry has been studied in detail in , and the Ricci tensor is given by:
$`R_{ab}=\left(\alpha ^2+{\displaystyle \frac{1}{32}}{\displaystyle \frac{\beta ^4}{\alpha ^2}}\right)\delta _{ab}`$ (6.7)
$`R_{AB}={\displaystyle \frac{3}{4}}\beta ^2\left(1{\displaystyle \frac{1}{16}}{\displaystyle \frac{\beta ^2}{\alpha ^2}}\right)\delta _{AB}`$ (6.8)
Again one can check explicitly that the same $`J_{abc}=ϵ_{abc},a,b,c=1,2,3`$ used in the previous solutions satisfies the field equations (3.6). Then $`AdS_3\times N^{010}`$ is a solution provided the rescalings are fixed to the values:
$$\alpha =\pm \frac{1}{6}|e|,\beta =\pm \frac{2}{3}|e|$$
(6.9)
As in the previous case supersymmetry is absent because (4.16) has no solutions.
In a similar way we easily find that for every 7-dimensional coset space $`G/H`$ of the classification there exist a set of constants $`J_{abc}`$ such that $`AdS_3\times G/H`$ is a solution for IIB supergravity, with the exception of the round 7-sphere $`SO(8)/SO(7)`$. Also, every $`G/H\times S^1`$ space with $`G/H`$ = nonsymmetric 6-dimensional coset is a solution of the IIB equations for appropriate rescalings and $`J_{abc}=C_{abc}`$.
What remains to be done is to determine in each case the most general $`J_{abc}`$, and check whether there are instances in which both supersymmetry conditions can be satisfied simultaneously, yielding supersymmetric solutions as in Section 5.
We expect that, in order to make contact with the $`d=2`$ superconformal theories discussed in , we need to extend our Ansatz to a nonconstant scalar field.
## 7 Asymptotic supersymmetry: $`S^3\times (P^2T^4)`$
We study here a particular class of solutions , characterized by a continuous parameter $`\sigma `$. Consider the cosets $`G/H=M^{010}`$, a special case of the $`M^{pqr}`$ class with $`S^3\times P^2`$ topology . In general a rescaling of the $`G`$ structure constants given by:
$$C_{G_1G_2}^{G_3}\frac{r_{G_1}r_{G_2}}{r_{G_3}}C_{G_1G_2}^{G_3}$$
(7.1)
still defines a Lie algebra. Take all $`r`$ = 1 except those in the $`4,5,6,7`$ coset directions, for which $`r=\sigma `$. Then the Ricci tensor becomes:
$`R_{mn}={\displaystyle \frac{1}{4}}b^2(2{\displaystyle \frac{b^2}{c^2}})\delta _{mn},m,n=1,2`$
$`R_{33}={\displaystyle \frac{b^4}{4c^2}}`$
$`R_{AB}={\displaystyle \frac{3}{4}}a^2\sigma ^2\delta _{AB},A,B=4,5,6,7`$ (7.2)
The limit $`\sigma 0`$ corresponds to a group contraction yielding $`S^3\times T^4`$: in fact it amounts to sending the radius of $`P^2`$ to infinity. For any $`\sigma `$, i.e. for any $`P^2`$-radius, $`AdS_3\times S^3\times P^2`$ is a solution of the IIB equations if we choose the $`J_{abc}`$ tensor to be:
$$J_{123}=1,J_{345}=J_{367}=\sigma $$
(7.3)
and the rescalings:
$$a^2=\frac{|e|^2}{3(1+2\sigma ^2)},b^2=\pm c|e|,c=\frac{\sigma ^2+1}{2\sigma ^2+1}|e|$$
(7.4)
One can see easily that the $`\delta \lambda `$ supersymmetry condition is satisfied by just one of these solutions, corresponding to $`\sigma =0`$, i.e. the $`N=8`$ supersymmetric $`AdS_3\times S^3\times T^4`$ compactification. All the other values of $`\sigma `$ break supersymmetry. Thus we obtain a class of continuously connected solutions, parametrized by $`\sigma `$, all of them nonsupersymmetric except in the limit $`\sigma =0`$. For $`\sigma 0`$, $`S^3`$ is a squashed three-sphere, cf. eq.s (7.2).
Acknowledgements
It is a pleasure to acknowledge useful discussions with Pietro Fré and Igor Pesando. |
warning/0003/nucl-th0003030.html | ar5iv | text | # How unique is the Asymptotic Normalisation Coefficient (ANC) method?
## I Introduction
Considerable effort, both theoretical and experimental, has been devoted in the last few years to the analysis of nuclear capture reactions. At the energies relevant for Astrophysics, (p,$`\gamma `$) or ($`\alpha ,\gamma `$) reactions have very low cross section values due to the Coulomb barrier repulsion. Thus, in many cases, the only access to the low energy region is through model dependent extrapolations of the higher energy data. In addition to this experimental limitation, many reactions of astrophysical interest involve radioactive beams which cannot be performed using conventional experimental techniques. The Coulomb dissociation and the asymptotic normalisation coefficients (ANCs) extracted from transfers have been put forward recently as alternative methods to obtain information about the astrophysical S-factors. As recognised by the physics community, while very appealing, these methods need to be subject to severe tests in order to assess their validity . The aim of this work is to check upon the validity of the ANC method.
Given the limited sets of data for peripheral transfer reactions, the results presented here may not be conclusive. However, we hope that this work will underline the present difficulties in validating the method and motivate further measurements.
We firstly describe the ANC method (section II). Then, we analyse in detail the different reactions that will be used in the present work (section III). Particular attention will be paid to the characteristics of the data. Finally we present a discussion of the results and conclusions in section IV.
## II A systematic study on proton transfer reactions
The ANC method for the transfer reaction
$`A+aB+b(a=b+x,B=A+x),`$ (1)
relies on two assumptions. Firstly, the reaction mechanism used to describe the transfer mechanism should give direct information of the nuclear overlap integrals $`A|B`$, $`a|b`$. The differential cross section is given by
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}={\displaystyle \frac{\mu _\mathrm{i}\mu _\mathrm{f}}{4\pi ^2\mathrm{}^4}}{\displaystyle \frac{k_b}{k_a}}{\displaystyle \frac{1}{(2J_A+1)(2J_a+1)}}{\displaystyle |T_{fi}|^2},`$ (2)
with $`\mu _\mathrm{i}`$, $`\mu _\mathrm{f}`$ the reduced masses for the initial ($`Aa`$) and final ($`bB`$) channels and $`k_a,k_b`$ the incident and outgoing momenta in the centre-of-mass frame. The DWBA reaction mechanism has been used to analyse the differential cross section for the transfer reaction. The transition amplitude for the transfer reaction process in the post form is
$`T_{fi}={\displaystyle \mathrm{\Psi }_f^{()}_{AB}|V_{xb}+V_{bA}U_{bB}|_{ab}\mathrm{\Psi }_i^{(+)}}.`$ (3)
In this equation $`\mathrm{\Psi }_f^{()}`$ and $`\mathrm{\Psi }_i^{(+)}`$ are the distorted waves in the final and initial channels respectively. $`_{AB}`$ and $`_{ab}`$ are the nuclear overlap integrals $`A|B`$ and $`a|b`$. The remnant term, $`V_{bA}U_{bB}`$, (where $`V_{bA}`$ is the interaction between the projectile core and the target A and $`U_{bB}`$ the optical potential for the outgoing channel), is usually small and may be neglected but will be included in our calculations.
Secondly, according to the ANC method, the transfer reaction should be peripheral so that the asymptotic part of the overlap integrals gives the dominant contribution to the transition amplitude. For example, outside the range $`R_N`$ of the $`A`$-$`x`$ nuclear interaction, the overlap integral $`A|B`$ becomes
$`_{AB\mathrm{}j}C_{AB\mathrm{}j}W_{\eta \mathrm{}+\frac{1}{2}}(2\kappa r)/rrR_N`$ (4)
where $`W_{\eta \mathrm{}+\frac{1}{2}}(2\kappa r)`$ is the Whittaker function, $`\eta =Z_AZ_xe^2\mu /\kappa `$ the Sommerfeld parameter, $`\kappa =\sqrt{2\mu ϵ}/\mathrm{}`$, $`\mu `$ and $`ϵ`$ the reduced mass and the binding energy for the ($`A`$$`x`$) system. $`C_{AB\mathrm{}j}`$ is the the asymptotic normalisation coefficient (ANC) for the overlap function $`A|B`$, related to the asymptotic normalisation of the single particle (s.p.) wave function $`b_{AC\mathrm{}j}`$ and a spectroscopic factor $`𝒮`$ by
$`C_{AB\mathrm{}j}=𝒮^{1/2}b_{AB\mathrm{}j}.`$ (5)
The ANC $`C_{AB\mathrm{}j}`$ defines the vertex for the virtual transitions $`BA+x`$ as shown in fig.(1). It has been shown that as long as the reaction is peripheral, the ANC is independent of the details of the s.p. parameters used to describe the nucleus B ground state. That is, the effect of different s.p. parameters (which result in different s.p. asymptotic normalisations) is compensated by the deduced experimental spectroscopic factors such that the ANCs become independent of the s.p. model.
For proton transfer reactions, the extracted ANCs gives an alternative method of determining the zero energy cross section for the capture reaction $`A+pB+\gamma `$ or alternatively S$`{}_{1A}{}^{}(0)`$ , providing of course that the overlap integral $`a|b`$ is known. The spectroscopic factor is obtained (by a $`\chi ^2`$ fit) from the ratio between the data and the DWBA calculation in the forward angle region and defined here as $`𝒮_{\mathrm{exp}}`$. To simplify notation we shall omit the angular momenta quantum numbers from the ANCs.
By choosing appropriate beam energies and scattering angles such that the transfer reaction remains peripheral this method is expected to provide a unique, structure model independent ANC.
The ANC method was firstly applied for extracting the S<sub>17</sub>-factor from the study of the reaction <sup>7</sup>Be(d,n)<sup>8</sup>B . The peripheral character of the reaction and the dependence on the optical potential for the incoming and outgoing channels have been recently studied . It was shown that for the DWBA analysis, the optical potentials for the entrance and outgoing channels need to be known in order to minimise uncertainties on the extracted S-factors .
The ANC method was also applied for extracting the S<sub>17</sub>-factor from the study of the <sup>10</sup>B(<sup>7</sup>Be,<sup>8</sup>B)<sup>9</sup>Be reaction . The transfer differential cross section was measured with high accuracy using an 84 MeV <sup>7</sup>Be radioactive beam, in the forward angle region, to ensure its peripheral character. The optical potentials for the incoming and outgoing channels were derived from folding model calculations and validated by the elastic data. From the measured transfer differential cross section the asymptotic normalisation for the virtual transitions <sup>8</sup>B $`^7`$Be + p and thus the S<sub>17</sub>(0) was extracted assuming that the ANC for the <sup>10</sup>B $`^9`$Be + p vertex was known. This ANC was determined in the same way from the analysis of <sup>9</sup>Be(<sup>10</sup>B,<sup>9</sup>Be)<sup>10</sup>B at an incident energy of 100 MeV .
Due to the increasing interest on this method it is timely to perform tests, to ensure its applicability. A first test of the ANC method was made in where the proton transfer reaction <sup>16</sup>O(<sup>3</sup>He,d)<sup>17</sup>F was analysed. It was shown in that work that the deduced S–factor for the capture reaction <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F agreed well with the capture data. In the present work further tests are performed.
Given the ANC $`C_{AB\mathrm{}j}`$, the question we address here is: how unique is this value, deduced from different proton transfer reactions or from the same reaction but at different energies, assuming that the peripheral character is satisfied ? For the present analysis we choose the case of the ANC for the <sup>10</sup>B $`^9`$Be + p vertex here called C<sub>19</sub>. This choice was motivated by the accurate forward angle data for the transfer reaction <sup>9</sup>Be(<sup>10</sup>B,<sup>9</sup>Be)<sup>10</sup>B, measured at a laboratory energy of 100 MeV, together with good knowledge of the optical potentials . A literature search was then performed to find proton transfer data at low energy from which independent values for C<sub>19</sub> could be extracted. In order to reduce the uncertainties associated with the lack of knowledge of the optical potential for the incoming and outgoing channel, our search was restricted to cases where elastic scattering data was available, whenever it was possible for both the incoming and outgoing channels. With these requirements in mind, the set of reactions used to study the uniqueness property are shown in table (VI).
A spherical two-body model is used for describing the ground state of each (B = A + p) system. We take a Woods-Saxon potential with radius of 1.25 fm, diffuseness a=0.65 fm, and depth adjusted to give the appropriate binding energy ($`ϵ`$(p + <sup>9</sup>B) = 6.5858 MeV, $`ϵ`$(p + <sup>12</sup>C) = 1.93435 MeV, and $`ϵ`$(p + d) = 5.4935 MeV). A spin-orbit term with the same geometry parameters and depth of 2.06 MeV was included. The <sup>10</sup>B g.s. was then described as $`p_{3/2}`$ proton coupled to the <sup>9</sup>Be(3/2<sup>-</sup>) core ($`p_{3/2}3/2^{}`$), the <sup>13</sup>N g.s. as ($`p_{1/2}0^+`$) and the <sup>3</sup>He g.s. as ($`s_{1/2}1^+`$). The two-body $`p^9`$Be s.p. model is not expected to provide a complete description for the <sup>10</sup>B system. In fact, the core ground state is close to the $`\alpha \alpha n`$ threshold and other terms may contribute significantly to the wave function. Within the present reaction mechanism framework, the incompleteness of the two-body model in describing the composite nucleus is taken into account through the extracted spectroscopic factor $`𝒮_{\mathrm{exp}}`$.
In order to extract the ANC $`C_{19}^2`$ from reaction ($`𝒜`$), we proceed in the same way as in . The transfer reaction ($``$) provides similarly information on the $`C_{19}^2`$. The $`{}_{}{}^{10}\mathrm{B}(\mathrm{d},^3\mathrm{He})^9\mathrm{Be}`$ reaction ($`𝒞`$) can be expressed in terms of the product $`C_{19}^2C_{12}^2`$ where $`C_{12}`$ is the ANC for the vertex $`{}_{}{}^{3}\mathrm{He}d+p`$ given in . As for reaction ($`𝒟`$), the DWBA cross section can be expressed in terms of the product $`C_{19}^2C_{112}^2`$ where $`C_{112}`$ is the ANC constant for the transition <sup>13</sup>N $``$ p + <sup>12</sup>C, that was extracted from the transfer reaction ($``$). In all cases the calculations were performed using FRESCO .
## III Results
The experimental analysis
The elastic and transfer experimental differential cross sections, for all the reactions we are considering, are shown in figs.(2-7). For each transfer reaction, we firstly determine a set of optical potential parameters that fit the elastic channels. In doing the elastic fit we take into account the whole angular range available. The starting parameters for this fit were taken from sets found in the literature at a nearby projectile energy. These were chosen to have considerable differences in order to truly evaluate the uncertainties on the ANCs due to the choice of optical potentials. The parameters for these optical potentials are collected in the Appendix.
The ANC is directly related to $`𝒮_{\mathrm{exp}}`$ (see eq.(5)). For a given pair of optical potentials (entrance and exit channels) $`𝒮_{\mathrm{exp}}`$ is the normalisation of the forward angle DWBA cross section coming from a $`\chi ^2`$ fit to the transfer data. The quality of the fit (the accuracy with which the DWBA predicted angular distribution is able to reproduce the angular dependence of the data) is quantified by $`\chi ^2=\frac{1}{N_{\mathrm{exp}}}_i\left(\frac{\sigma _{\mathrm{exp}}(i)𝒮_{\mathrm{exp}}\sigma _{\mathrm{Theo}}(i)}{\mathrm{\Delta }\sigma _{\mathrm{exp}}(i)}\right)^2`$ with $`N_{\mathrm{exp}}`$ the number of experimental points. These optical potentials are presented in tables (VI-VI). To evaluate the effect of different choices of optical potentials on the calculated transfer differential cross section, we also show in these tables the corresponding $`\chi ^2`$ values for the transfer.
Since the aim of this method is to extract an overall normalisation of the transfer data, it is not only essential to have data with low statistical errors but, more importantly, low systematic errors. The uncertainty in the target thickness is a large contributor to the systematic error, except for the data in where special attention was paid to this issue. For the (d,n) reactionm the neutron efficiency uncertainty is also quite significant. Other typical errors that may arise are due to beam collection or error in solid angle, but are much lower than those mentioned above. The systematic errors for the set of reactions are collected in table (VI), according to the information in the literature. It is evident from table (VI) that only the normalisation error of the data from has the desired low value.
Another source of error could come from the angular range from which data is being considered. We must ensure that the peripheral character of the transfer reaction is satisfied. A large number of experimental points in the forward angle region (where the transfer is clearly peripheral) is desirable, but in most cases non-existent. For this reason, we have taken a forward angle subset of the data: the first 7 points. Even then some sets have angular ranges up to $`40^{}`$ (see table VI). In two cases a smaller set of data points had to be chosen. The worst example we have considered is for the (d,n) reactions.
Finally, there will be errors on the derived ANCs arising from uncertainties on the optical potentials since the elastic data does not totally probe the interaction. The ANC errors shown in table(VI) and quoted in the text are associated only with the optical potential uncertainties.
An overall panorama of the errors involved when using this method is given in fig.(8). It is clear from what has been presented in this section that our results should be interpreted as indicators until further measurements are available. They underline the need for further experimental work, before definite conclusions on the uniqueness property of the ANC method can be drawn.
$`{}_{}{}^{9}\mathrm{Be}(^{10}\mathrm{B},^9\mathrm{Be})^{10}\mathrm{B}`$
The reaction <sup>9</sup>Be(<sup>10</sup>B,<sup>9</sup>Be)<sup>10</sup>B is particularly adequate for extracting the ANC since it has the same vertex for the incoming and outgoing channels. The experimental spectroscopic factor $`𝒮_{\mathrm{exp}}`$ should then be proportional to $`C_{19}^4`$. The elastic and transfer data at E<sub>lab</sub>=100 MeV was taken from . We take four sets of Wood-Saxon optical model potentials for the incoming elastic channel: the first two obtained by a fitting procedure to elastic data shown in fig.(2a) and the others taken from Mukhamedzhanov et al. . For this particular reaction, the starting parameter set used in fitting the elastic data was taken from Comer at 40 MeV. The fits to the elastic scattering are shown in fig.(2a). The calculated DWBA transfer cross section renormalised by the spectroscopic factor, fig.(2b), reproduces quite well the transfer data at small angles.
We note that, the calculated ANC is not strongly dependent on the details of the optical potentials pinned by the elastic data. Even if unrealistically shallow potentials are used, the ANC hardly changes. Thus, in this case, the uncertainties associated with the choice of the optical potentials are very small.
The calculated reaction cross section as a function of the partial wave, fig.(9a), clearly shows that this reaction is peripheral. In fact, the cross section at E<sub>lab</sub>=100 MeV peaks around L=24 corresponding to an impact parameter of 7.29 fm which is significantly greater than the sum of the <sup>9</sup>Be and <sup>10</sup>B interaction radius . For that reason the reaction cross section only becomes significantly smaller for a cuttoff radius much bigger than the interaction radius. We obtained for the ANC, $`C_{19}^2`$= 4.9 $`\pm `$ 0.25 fm<sup>-1</sup>, where, as mentioned before, the error is associated with an optical potential uncertainty. This is in good agreement with the value obtained in .
$`{}_{}{}^{9}\mathrm{Be}(\mathrm{d},\mathrm{n})^{10}\mathrm{B}`$
The reaction <sup>9</sup>Be(d,n)<sup>10</sup>B provides direct information on the $`C_{19}^2`$ as the vertex for the deuteron is well known. We performed calculations for 2 different deuteron laboratory energies: 7 and 15 MeV.
At 7 MeV we used transfer data from Park and elastic data for <sup>9</sup>Be(d,d) at 6.3 MeV from Djaloeis . For the incoming channel, optical model potential parameters were obtained from fitting the data shown in fig.(3 a). For the outgoing channel, the potential parameters were taken from Dave and Gould . As follows from the fig.(3 b), the calculated transfer differential cross section describes quite well the data. We extracted an ANC of $`C_{19}^2`$ = 4.8 $`\pm `$ 0.35 fm<sup>-1</sup> which is in good agreement with our previous result obtained from the analysis of reaction $`𝒜`$.
At 15 MeV we also used transfer data from Park and elastic <sup>9</sup>Be(d,d) at 15 MeV from Armstrong published by . Four entrance potential parameter sets were obtained fitting the data. The outgoing parameter set was taken from at the appropriate energy. As can be seen from fig.(3 c), and fig.(3 d) the 4 parameter sets used describe quite well the elastic data, but none is able to reproduce satisfactory the transfer data in the low angle region ($`\theta 20^{}`$). As a result, the calculated ANCs depend crucially on both the input parameter set type, surface or volume (about 20 %), and on the low angle region chosen to minimise $`\chi ^2`$ in order to obtain $`𝒮_{\mathrm{exp}}`$ (vide table(VI)). The derived value for the ANC at this energy is $`C_{19}^2`$ = 6.09 $`\pm `$ 0.54 fm<sup>-1</sup> which is higher than that found by . As can be concluded from fig.(11a) and fig.(11c), while the reaction is peripheral at 7 MeV, this is no longer the case for 15 MeV. Thus, this data is not useful for the purpose of this work.
$`{}_{}{}^{10}\mathrm{B}(\mathrm{d},^3\mathrm{He})^9\mathrm{Be}`$
The proton pickup reaction <sup>10</sup>B(d,<sup>3</sup>He)<sup>9</sup>Be has two different vertices, <sup>10</sup>B $`^9`$Be + p and <sup>3</sup>He $``$ d + p, and therefore the experimental spectroscopic factor will be proportional to the ANCs product $`C_{19}^2C_{12}^2=𝒮_{\mathrm{exp}}b_{19}^2b_{12}^2`$.
The transfer and elastic data for this reaction at E<sub>lab</sub>=11.8 MeV was taken from . For the entrance channel, one parameter set was obtained fitting the elastic data shown in fig.(4a). For the exit channel, we used three parameter sets from literature: the first from , the second from and the third from . The description of transfer data shown in fig.(4b) is very reasonable, specially for the low angle region. The analysis of the reaction cross section as a function of the partial wave number shows that this transfer reaction is peripheral fig.(9b). By renormalising the calculated DWBA differential cross section from the data we obtained $`C_{19}^2C_{12}^2=19.17\pm 1.82`$ fm<sup>-2</sup>. For the <sup>3</sup>He $``$ d + p vertex, we used the value taken from $`C_{12}^2`$ = 3.9 $`\pm `$ 0.06 fm<sup>-1</sup>. Consequently, we get $`C_{19}^2`$ = 4.92 $`\pm `$ 0.54 fm<sup>-1</sup>, in good agreement with the result of .
$`{}_{}{}^{12}\mathrm{C}(^{10}\mathrm{B},^9\mathrm{Be})^{13}\mathrm{N}`$
We proceed in our systematics by looking at other proton stripping reactions involving <sup>9</sup>Be. A candidate for which measured data was found is <sup>12</sup>C(<sup>10</sup>B,<sup>9</sup>Be)<sup>13</sup>N at 100 MeV . The elastic scattering fig.(5a) was taken from . We also take the same data for the outgoing channel due to the absence of experimental measurements for this channel. When fitting the elastic data for the incoming channel, we obtain fouor parameter sets. The optical potential parameters for the outgoing elastic channel were taken to be the same but with an appropriate radius as discussed in the Appendix. The calculated DWBA cross section describes quite well the transfer data shown in fig.(5b) specially for sets 2, 3 and 4.
As this reaction has two different vertices for the two composite nuclei, $`{}_{}{}^{10}\mathrm{B}^9\mathrm{Be}+\mathrm{p}`$ and <sup>13</sup>N $`^{12}`$C + p, the experimental spectroscopic factor $`𝒮_{\mathrm{exp}}`$ will be proportional to the ANCs product, $`C_{19}^2C_{112}^2=𝒮_{exp}b_{19}^2b_{112}^2`$. We obtained $`C_{19}^2C_{112}^2=7.4\pm 0.5\mathrm{fm}^2`$.
$`{}_{}{}^{12}\mathrm{C}(\mathrm{d},\mathrm{n})^{13}\mathrm{N}`$
In order to extract $`C_{19}^2`$ from the results obtained with the last reaction, it is necessary to extract $`C_{112}^2`$ from another independent reaction. We chose the <sup>12</sup>C(d,n)<sup>13</sup>N reaction at two different energies: 9 and 12.4 MeV.
At 9 MeV we used transfer data available from two different sources, Davis et al. and Schelin et al. . We take 3 sets of potential parameters from fitting the entrance channel elastic data of shown in fig.(6a), and one set for the exit channel data of for <sup>13</sup>C(n,n) at 10 MeV shown in fig.(6b). As shown in fig.(6c) the shape of the calculated DWBA differential cross section adjusts much better to Davis’s data than to Schelin’s in the low angle region, $`\theta 35^{}`$. Not surprisingly the derived spectroscopic factors from the two data sets differ by 30 $`\%`$. We found $`C_{112}^2=2.56\pm 0.37`$ for and $`C_{112}^2=3.31\pm 0.45`$ for . Again, these errors are associated only with the optical potential uncertainty. Using the results for the ANC product from reaction <sup>12</sup>C(<sup>10</sup>B,<sup>9</sup>Be)<sup>13</sup>N we obtain respectively $`C_{19}^2`$ = 2.98 $`\pm `$ 0.63 fm<sup>-1</sup> and $`C_{19}^2`$ = 2.30 $`\pm `$ 0.47 fm<sup>-1</sup>.
An energy average differential cross section data, at 12.4 MeV, is given in Schelin’s work (using 13.0 MeV data and 11.8 MeV data of Mutchler ). For the entrance channel we used the Matusevich experimental points at 13.6 MeV fig.(7a). The exit channel data was taken from for <sup>13</sup>C(n,n) at 12 MeV fig.(7b). Three parameter set fits were obtained for the deuteron potential and one for the neutron potential.
As follows from fig.(7c), a good agreement in the low angle region ($`\theta 40^{}`$) is achieved between the data and the calculated cross section. The calculated spectroscopic factors lead us to $`C_{112}^2`$ = 1.65 $`\pm `$ 0.2 fm<sup>-1</sup>. Using the ANC product from the reaction for <sup>12</sup>C(<sup>10</sup>B,<sup>9</sup>Be)<sup>13</sup>N, one obtains $`C_{19}^2`$ = 4.6 $`\pm `$ 0.9 fm<sup>-1</sup>.
We note from figures (11b) and (11d) that while the reaction $`{}_{}{}^{12}\mathrm{C}(\mathrm{d},\mathrm{n})^{13}\mathrm{N}`$ is peripheral at 9 MeV, the situation is rather unclear at 12.4 MeV. Although the impact parameter is greater than the <sup>12</sup>C interaction radius , the reaction cross section drops significantly for a cutoff radius R<sub>cut</sub> = 1 fm.
The experimental situation concerning this reaction is rather unsatisfactory due to different available data at 9 MeV. In the works of Davis and Shelin , a contribution due to compound nucleus formation to the cross section is estimated using the Hauser-Feshbach statistical model. For the $`{}_{}{}^{12}\mathrm{C}(\mathrm{d},\mathrm{n})^{13}\mathrm{N}`$ reaction at 12.4 MeV, the calculated compound nucleus cross section lead to an overestimation of the cross section at large angles. The Hauser-Feshbach model is then clearly inadequate in this case. Arbitrary reduction factors can be found in producing a large uncertainty on the derived spectroscopic factors and ANCs. In order to extract a meaningful ANC factor from a transfer reaction, the reaction mechanism should be properly understood.
Since no information on the uncertainty of the absolute cross section for reaction $`𝒟`$ is given the C<sub>19</sub> thus extracted should not be used to validate the ANC method. Agravating the situation are the differences between the <sup>12</sup>C(d,n) data sets suggesting that the C$`_{112}`$ are not sufficiently reliable to be taken into account.
## IV Conclusions
We have determined the ANC for the <sup>10</sup>B $`^9`$Be + p using a set of proton-transfer reactions at different energies. The calculated ANCs from the different reactions reproduced in fig.(8) clearly reveals the present experimental situation if one wants to check the validity of the ANC method. The sum of the contributions of the statistical, optical potential and systematic uncertainties is in most cases quite large. The graph evidently shows that, from a particular set of transfer reactions (those that are clearly peripheral and have quotable normalisation errors), the uniqueness property of the ANC is satisfied. However, with the present data, we cannot undoubtedly conclude if this property is fulfilled.
More data for both, the transfer and elastic channels, with good resolution and carefully normalised cross sections in the forward angular region, is crucial if we want to unambiguously check the uniqueness of the ANCs. The elastic data is essential to reduce the optical parameter uncertainties. In the measurements special attention should be paid to minimise the uncertainty on the target thickness. For (d,n) reaction it is important to reduce as much as possible the neutron efficiency error, given that this may be a large source of uncertainty.
Even though the DWBA method is widely used, care should be taken to fully understand the reaction mechanisms before extracting the ANCs. Early studies on the deuteron breakup effects on the differential cross section indicate that DWBA analysis may not be a useful tool to study deuteron transfer reactions. However, even nowadays, these reactions are still used to extract ANCs . More data on deuteron transfer reactions is necessary in order to have a better understanding of the mechanisms and to check if they can be used to extract the ANCs. Generally, further tests on the ANC method, focusing on the reaction mechanism, should also be performed.
###### Acknowledgements.
This work was supported by Fundação de Ciência e Tecnologia (Portugal) through grant Praxis PCEX/C/FIS/4/96. We would like to thank L. Trache for providing us with the experimental results for the <sup>9</sup>Be(<sup>10</sup>B, <sup>9</sup>Be)<sup>10</sup>B reaction and elastic data.
## Appendix
We collect in tables (VI-VI) the optical potential parameters obtained by fitting the elastic channels. The potentials are calculated using the following expressions:
$`\mathrm{Real}\mathrm{central}:`$ $`U_R`$ $`=V{\displaystyle \frac{f(r)}{1+f(r)}},`$ (6)
$`\mathrm{Imaginary}\mathrm{central}\mathrm{volume}:`$ $`U_I`$ $`=W{\displaystyle \frac{f(r)}{1+f(r)}},`$ (7)
$`\mathrm{Imaginary}\mathrm{central}\mathrm{surface}:`$ $`U_W`$ $`=\mathrm{\hspace{0.17em}4}W_d{\displaystyle \frac{f(r)}{(1+f(r))^2}},`$ (8)
$`\mathrm{Spin}\mathrm{Orbit}:`$ $`U_{SO}`$ $`={\displaystyle \frac{4}{ra_{SO}}}V_{SO}{\displaystyle \frac{f(r)}{(1+f(r))^2}}\stackrel{}{l}.\stackrel{}{s},`$ (9)
with $`f(r)=e^{\frac{rR}{a_i}}`$ and $`R=r_iA_T^{1/3}`$ except the <sup>12</sup>C(<sup>10</sup>B,<sup>9</sup>Be)<sup>13</sup>N case where $`R=r_i(A_P^{1/3}+A_T^{1/3})`$. For all set of optical potentials, we use $`r_i=r_0`$ for real central, etc. |
warning/0003/hep-th0003278.html | ar5iv | text | # Contents
## 1 Introduction and Summary
The world-volume theory on the D-brane of a bosonic string theory contains a tachyonic mode. It has been conjectured that the tachyon potential has a non-trivial extremum where the potential energy of the tachyon exactly cancels the tension of the D-brane, and that this configuration represents the closed string vacuum without any D-brane. It has been further conjectured that various tachyonic lump solutions on the D-brane world-volume represent D-branes of lower dimensions. These conjectures and their generalisations to superstring theories have been tested by various methods,. However, since the exact effective action for the tachyon field is not known, there is no direct proof of these conjectures.
In this paper we point out that in the $`p`$-adic string theory introduced in (see for a review) we can explicitly check these conjectures. It should be emphasised that although the $`p`$-adic ‘string’ is an exotic object, the spacetime it describes is the familiar one<sup>3</sup><sup>3</sup>3A different type of $`p`$-adic string was considered in . . In the $`p`$-adic open string theory, which in modern language can be regarded as the world-volume theory of a space-filling D-brane, the exact classical action of the tachyon field and various solutions of the equations of motion are known. Among the known non-trivial solutions is a translationally invariant solution with the property that it is a local minimum of the potential, and that the propagator of the tachyon field describing fluctuations around this background has no physical pole. Thus this configuration has no physical open string excitations, and is naturally identified with the vacuum without a D-brane. The exact tachyon equation of motion of the $`p`$-adic string theory also has classical lump solutions for all codimension $`1`$, which approach the vacuum solution far away from the core of the soliton. If the original open string theory is formulated in $`(d1,1)`$ dimensional space-time<sup>4</sup><sup>4</sup>4 There is as yet no compelling reason for a critical dimension in $`p`$-adic string theory, but the so called adelic formula relating the product of four tachyon amplitudes in $`p`$-adic strings for all primes $`p`$ to that in the bosonic string suggests that they all have the same critical dimension $`d=26`$., then such a lump solution of codimension $`(dq1)`$ describes a solitonic $`q`$-brane. We show that the world-volume theory on the solitonic $`q`$-brane agrees with the expected world-volume theory on a Dirichlet $`q`$-brane in the $`p`$-adic string theory, to the extent that we can compare them with the present knowledge. This provides strong evidence that these lump solutions can be identified as lower dimensional D-branes.
The paper is organised as follows. In section 2 we summarise the exact effective action of the tachyon in the $`p`$-adic string theory, the known solutions of the equation of motion derived from the action and their properties. In section 3 we analyse the world-volume theory of the solitonic $`q`$-brane, and in section 4 we compare this with the world-volume theory of a Dirichlet $`q`$-brane. Section 5 contains some comments on further extension of this work, and ends with speculation on its possible application to the study of tachyon condensation in ordinary bosonic string theory.
## 2 Solitonic $`q`$-branes of $`p`$-adic string theory
In ref. $`p`$-adic string theory was defined as follows. Consider the expressions for various amplitudes in ordinary bosonic open string theory, written as integrals over the boundary of the world-sheet which is the real line R. Now replace the integrals over R by integrals over the $`p`$-adic field $`𝐐_p`$ with appropriate measure, and the norms of the functions in the integrand by the $`p`$-adic norms. These rules were subsequently derived from a local action defined on the “world-sheet” of the $`p`$-adic string. Using $`p`$-adic analysis, it is possible to compute $`N`$ tachyon amplitudes at tree-level for all $`N3`$.
This leads to an exact action for the open string tachyon in $`d`$ dimensional $`p`$-adic string theory. This action is given in ref.
$`S`$ $`=`$ $`{\displaystyle d^dx}`$ (2.1)
$`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \frac{p^2}{p1}}{\displaystyle d^dx\left[\frac{1}{2}\varphi p^{\frac{1}{2}\mathrm{}}\varphi +\frac{1}{p+1}\varphi ^{p+1}\right]},`$
where $`\mathrm{}`$ denotes the $`d`$ dimensional Laplacian, $`\varphi `$ is the tachyon field (after a rescaling and a shift), $`g`$ is the open string coupling constant, and $`p`$ is an arbitrary prime number. We are using metric with signature $`(,+,+\mathrm{}+)`$. If we denote by $`(2\pi \alpha _p^{})^1`$ the ‘tension of the $`p`$-adic string’ as defined in ref., then our choice of units correspond to
$$\alpha _p^{}=\frac{\mathrm{ln}p}{2\pi }.$$
(2.2)
We have added a constant term to the Lagrangian density $``$ so that it vanishes at $`\varphi =0`$. Fig.1 shows the qualitative features of the tachyon potential for different values of $`p`$.
The equation of motion derived from this action is,
$$p^{\frac{1}{2}\mathrm{}}\varphi =\varphi ^p.$$
(2.3)
Different known solutions of this equation are as follows:
* The configuration $`\varphi =1`$ is the original vacuum around which we quantised the string<sup>5</sup><sup>5</sup>5For $`p2`$, there is also an equivalent solution corresponding to $`\varphi =1`$. Since the action is symmetric under $`\varphi \varphi `$, we shall restrict our analysis to solutions with positive $`\varphi `$.. We shall call this the D-$`(d1)`$-brane solution. The energy density associated with this configuration, which can be identified as the tension $`T_{d1}`$ of the D-$`(d1)`$-brane configuration, is given by
$$T_{d1}=(\varphi =1)=\frac{1}{2g^2}\frac{p^2}{p+1}.$$
(2.4)
* The configuration $`\varphi =0`$ denotes a configuration around which there is no perturbative physical excitation. We shall identify this with the vacuum configuration. By definition we have taken the energy density of this vacuum to be zero.
* The configuration:
$$\varphi (x)=f(x^{q+1})f(x^{q+2})\mathrm{}f(x^{d1})F^{(dq1)}(x^{q+1},\mathrm{},x^{d1}),$$
(2.5)
with
$$f(\eta )p^{\frac{1}{2(p1)}}\mathrm{exp}\left(\frac{1}{2}\frac{p1}{p\mathrm{ln}p}\eta ^2\right),$$
(2.6)
describes a soliton solution with energy density localised around the hyperplane $`x^{q+1}=\mathrm{}=x^{d1}=0`$. This follows from the identity:
$$p^{\frac{1}{2}_\eta ^2}f(\eta )=\left(f(\eta )\right)^p.$$
(2.7)
We shall call (2.5), with $`f`$ as in (2.6), the solitonic $`q`$-brane solution. Let us denote by $`x_{}=(x^{q+1},\mathrm{},x^{d1})`$ the coordinates transverse to the brane and by $`x_{}=(x^0,\mathrm{},x^q)`$ those tangential to it. The energy density per unit $`q`$-volume of this brane, which can be identified as its tension $`T_q`$, is given by
$$T_q=d^{dq1}x_{}(\varphi =F^{(dq1)}(x_{}))=\frac{1}{2g_q^2}\frac{p^2}{p+1},$$
(2.8)
where,
$$g_q=g\left[\frac{p^21}{2\pi p^{2p/(2p1)}\mathrm{ln}p}\right]^{(dq1)/4}.$$
(2.9)
From eqs.(2.4),(2.8) and (2.9) we see that the ratio of the tension of a $`q`$-brane to a $`(q1)`$-brane is
$$\frac{T_q}{T_{q1}}=\left[\frac{2\pi p^{\frac{2p}{p1}}\mathrm{ln}p}{p^21}\right]^{\frac{1}{2}}=\frac{\sqrt{p^21}}{p^{\frac{p}{p1}}}\frac{1}{2\pi \sqrt{\alpha _p^{}}}.$$
(2.10)
In the above equation we have used dimensional analysis and (2.2) to restore factors of $`\alpha _p^{}`$. Note that the ratio (2.10) is independent of $`q`$. This is a feature of the D-branes in ordinary bosonic string theory, and suggests that the solitonic $`q`$-branes of $`p`$-adic string theory should have interpretation as D-branes. This also suggests that the self-dual radius $`R_{sd}`$ of the $`p`$-adic string theory, where the tension $`2\pi R_{sd}T_q`$ of a wrapped $`q`$-brane is equal to the tension $`T_{q1}`$ of a $`(q1)`$-brane, is given by
$$R_{sd}=\frac{p^{p/(p1)}}{\sqrt{p^21}}\sqrt{\alpha _p^{}}.$$
(2.11)
Note that as $`p\mathrm{}`$, this approaches the formula for the self-dual radius in ordinary bosonic string theory.
## 3 World-volume theory on the solitonic $`q`$-branes
Let us now consider a configuration of the type
$$\varphi (x)=F^{(dq1)}(x_{})\psi (x_{}),$$
(3.1)
with $`F^{(dq1)}(x_{})`$ as defined in (2.5),(2.6). For $`\psi =1`$ this describes the solitonic $`q`$-brane. Fluctuations of $`\psi `$ around 1 denote fluctuations of $`\varphi `$ localised on the soliton; thus $`\psi (x_{})`$ can be regarded as one of the fields on its world-volume. We shall call this the tachyon field on the solitonic $`q`$-brane world-volume<sup>6</sup><sup>6</sup>6In the linearised approximation this tachyonic mode was discussed in ref... Substituting (3.1) into (2.3) and using (2.7) we get
$$p^{\frac{1}{2}\mathrm{}_{}}\psi =\psi ^p,$$
(3.2)
where $`\mathrm{}_{}`$ denotes the $`(q+1)`$ dimensional Laplacian involving the world-volume coordinates $`x_{}`$ of the $`q`$-brane. The action involving $`\psi `$ can be obtained by substituting (3.1) into (2.1):
$`S_q(\psi )`$ $`=`$ $`S\left(\varphi =F^{(dq1)}(x_{})\psi (x_{})\right)`$ (3.3)
$`=`$ $`{\displaystyle \frac{1}{g_q^2}}{\displaystyle \frac{p^2}{p1}}{\displaystyle d^{q+1}x_{}\left[\frac{1}{2}\psi p^{\frac{1}{2}\mathrm{}_{}}\psi +\frac{1}{p+1}\psi ^{p+1}\right]},`$
where $`g_q`$ has been defined in eqn.(2.9).
Note that a solution of (3.2) gives an exact solution of the full equation of motion (2.3). Thus eq.(3.2) describes the dynamics of the mode $`\psi `$ on the $`q`$-brane world-volume exactly. This does not mean that there are no other modes on the $`q`$-brane world-volume; rather what this implies is that it is possible to obtain a consistent truncation of the world-volume theory of the $`q`$-brane by setting all the modes except $`\psi `$ to zero. In terms of scattering amplitudes this means that the tree level S-matrix on the $`q`$-brane world-volume, involving only external tachyon states, can be calculated exactly from the action (3.3).
Of the various other (infinite number of) modes living on the $`q`$-brane world-volume are the $`(dq1)`$ massless modes $`\xi ^i`$ associated with translations of the brane in the $`(dq1)`$ directions $`x_{}`$ transverse to the brane. Inclusion of these modes correspond to deformation of $`\varphi `$ of the form
$$\varphi (x)=F^{(dq1)}(x_{})\psi (x_{})+_{x_{}^i}F^{(dq1)}(x_{})\xi ^i(x_{})+\mathrm{}.$$
(3.4)
Substituting this in eq.(2.3), and comparing the coefficients of the independent functions $`\left(F(x_{})\right)^p`$ and $`\left(F(x_{})\right)^{p1}_{x_{}^i}F(x_{})`$ on both sides, we get the following equations of motion:
$`p^{\frac{1}{2}\mathrm{}_{}}\psi `$ $`=`$ $`\psi ^p+𝒪(\xi ^2)`$
$`p^{\frac{1}{2}\mathrm{}_{}}\xi ^i`$ $`=`$ $`\psi ^{p1}\xi ^i+𝒪(\xi ^2).`$ (3.5)
The above equations can be derived from the effective action:
$`S_q(\psi ,\xi ^i)`$ $`=`$ $`{\displaystyle \frac{1}{g_q^2}}{\displaystyle \frac{p^2}{p1}}{\displaystyle }d^{q+1}x_{}[{\displaystyle \frac{1}{2}}\psi p^{\frac{1}{2}\mathrm{}_{}}\psi +{\displaystyle \frac{1}{p+1}}\psi ^{p+1}`$ (3.6)
$`C\{{\displaystyle \frac{1}{2}}\xi ^ip^{\frac{1}{2}\mathrm{}_{}}\xi ^i{\displaystyle \frac{1}{2}}\psi ^{p1}\xi ^i\xi ^i\}+𝒪(\xi ^3)].`$
$`C`$ is a normalisation constant whose value is not important to this order, as it can be changed by rescaling $`\xi ^i`$.
$`\psi =1`$ corresponds to the solitonic $`q`$-brane solution. For computing amplitudes involving the world-volume fields on the $`q`$-brane, we define the shifted field $`\sigma `$ and rescaled fields $`\chi ^i`$ through the relation
$$\psi =1+\frac{g_q\sigma }{p},\xi ^i=\frac{g_q}{\sqrt{pC}}\chi ^i,$$
(3.7)
and expand the action (3.6) in powers of $`\sigma `$ and $`\chi ^i`$. This gives
$`S_q`$ $`=`$ $`{\displaystyle \frac{p}{p1}}{\displaystyle }d^{q+1}x_{}[{\displaystyle \frac{1}{2}}\sigma p^{\frac{1}{2}\mathrm{}_{}1}\sigma +{\displaystyle \frac{1}{g_q^2}}{\displaystyle \frac{p}{p+1}}(1+{\displaystyle \frac{g_q\sigma }{p}})^{p+1}{\displaystyle \frac{\sigma }{g_q}}{\displaystyle \frac{p}{2g_q^2}}`$ (3.8)
$`{\displaystyle \frac{1}{2}}\chi ^ip^{\frac{1}{2}\mathrm{}_{}}\chi ^i+{\displaystyle \frac{1}{2}}(1+{\displaystyle \frac{g_q\sigma }{p}})^{p1}\chi ^i\chi ^i+𝒪(\chi ^3)].`$
There is no linear term in $`\sigma `$ in action (3.8) reflecting the fact that $`\sigma =0`$ is a solution of the equation of motion. The momentum space $`\sigma `$ and the $`\chi ^i`$ propagators computed from this action are given by:
$`\mathrm{\Delta }_{\sigma \sigma }(k)`$ $`=`$ $`i{\displaystyle \frac{p1}{p}}{\displaystyle \frac{1}{p^{\frac{1}{2}k^21}1}},`$
$`\mathrm{\Delta }_{\chi ^i\chi ^j}(k)`$ $`=`$ $`i{\displaystyle \frac{p1}{p}}{\displaystyle \frac{1}{p^{\frac{1}{2}k^2}1}}\delta _{ij}.`$ (3.9)
The residues at the poles in the $`\sigma `$ and the $`\chi ^i`$ propagators (at $`k^2=2`$ and $`k^2=0`$ respectively) have the same values. This will help us compare the amplitudes involving $`\sigma `$’s and $`\chi ^i`$’s in the external legs.
## 4 Comparison with the world-volume theory on a Dirichlet $`q`$-brane
We shall now compare the world-volume action on the solitonic $`q`$-brane with that on the D-$`q`$-brane. We are immediately faced with the question whether it is possible to have Dirichlet branes in $`p`$-adic string theory. Fortunately, a ‘world-sheet’ approach to $`p`$-adic string has been developed in ref.. According to these authors, this ‘world-sheet’ is a so called Bruhat-Tits tree — a Bethe lattice with $`p+1`$ nearest neighbours — the ‘boundary’ of which is the field of $`p`$-adic numbers $`𝐐_p`$. The generalisation of the Polyakov action is the lattice discretisation of the action for free scalar fields corresponding to the target space coordinates. Now one can either choose Neumann or Dirichlet boundary conditions as in the case of ordinary strings. It was shown in that Neumann boundary conditions leads to the tachyon amplitudes postulated in .
While this may be the proper way to define D-branes in $`p`$-adic string theory, we shall content ourselves with the continuation of the relevant formulæ from ordinary bosonic string theory. Thus, for our purposes the world-volume theory of a $`p`$-adic D-$`q`$-brane is defined by taking the expressions for various amplitudes for an ordinary D-$`q`$-brane, written as integrals over world-sheet coordinates of the appropriate vertex operators, and then replacing the integrals over real line by integrals over the $`p`$-adic field, with all the norms appearing in the integrand replaced by $`p`$-adic norms as in ref.. In principle one should be able to derive these rules from the world-sheet description in ref..
For amplitudes involving the external tachyons, described by the vertex operators of the type $`e^{ikX_{}}`$ with momentum $`k`$ restricted to lie along the world-volume of the D-$`q`$-brane, the computation of the amplitude is identical to the one described in ref.. Thus, following the analysis there, these S-matrix elements can be obtained from an effective action of the form:
$$\widehat{S}_q(\psi )=\frac{1}{\widehat{g}_q^2}\frac{p^2}{p1}d^{q+1}x_{}\left[\frac{1}{2}\psi p^{\frac{1}{2}\mathrm{}_{}}\psi +\frac{1}{p+1}\psi ^{p+1}\right],$$
(4.1)
where the tachyon field $`\psi `$ is shifted so that $`\psi =1`$ describes the D-$`q`$-brane, and $`\widehat{g}_q`$ denotes the coupling constant which characterises the strength of the interaction in the world-volume theory of the D-$`q`$-brane. Comparing this with (3.3) we see that the world-volume actions for the tachyon fields on the solitonic $`q`$-brane and the Dirichlet $`q`$-brane agree exactly if we choose:
$$\widehat{g}_q=g_q.$$
(4.2)
At present there is no independent derivation of $`\widehat{g}_q`$ in terms of $`g`$, and hence we cannot verify eqn.(4.2) independently<sup>7</sup><sup>7</sup>7This is related to the problem of computing the tension of the D-$`q`$-brane independently.. But assuming (4.2) to be true, we have a complete agreement between the world-volume theories involving the tachyon fields on the D-$`q`$-brane and the solitonic $`q`$-brane.
Next we shall compare the amplitude $`\sigma \sigma \chi ^i\chi ^j`$ on the solitonic $`q`$-brane and the D-$`q`$-brane. First let us compute this on the solitonic $`q`$-brane using the action (3.8). The four Feynmann diagrams contributing to it have been shown in Fig.2. These can be easily evaluated, and the answer is:
$$A_{ij}(k_1,k_2,k_3,k_4)=\delta _{ij}g_q^2\left[\frac{p2}{p}+\frac{p1}{p}\left\{\frac{1}{p^{k_1k_2+1}1}+\frac{1}{p^{k_1k_3+1}1}+\frac{1}{p^{k_1k_4+1}1}\right\}\right],$$
(4.3)
with the contribution to the four terms in the right hand side of (4.3) coming from the Feynman digrams (a), (b), (c) and (d) respectively in Fig.2. In deriving (4.3) we have used the mass-shell conditions
$$k_1^2=k_4^2=2,k_2^2=k_3^2=0.$$
(4.4)
Let us now evaluate the same amplitude on the D-$`q`$-brane. The vertex operator associated with the mode $`\chi ^i`$ on the D-$`q`$-brane is given by<sup>8</sup><sup>8</sup>8The notation $`X`$ is schematic, as care is needed to define the correct vertex operator that corresponds to the analogous one for ordinary string. $`X_{}^ie^{ikX_{}}`$. Inserting the $`\chi ^i`$ and $`\chi ^j`$ vertex operators carrying momenta $`k_2`$ and $`k_3`$ at 0 and 1 respectively, and the two $`\sigma `$ vertex operators carrying momenta $`k_1`$ and $`k_4`$ at $`x`$ and $`\mathrm{}`$ respectively, we can express the amplitude as:
$$\widehat{A}_{ij}(k_1,k_2,k_3,k_4)=\delta _{ij}\widehat{g}_q^2_{𝐐_p}𝑑x|x|^{k_1k_2}|1x|^{k_1k_3}.$$
(4.5)
Here $`||`$ denotes the $`p`$-adic norm and integral over $`x`$ is over the $`p`$-adic field. This is precisely the integral evaluated in . Using the identity
$$k_1k_2+k_1k_3+2=k_1k_4,$$
(4.6)
we can express this amplitude as
$$\widehat{A}_{ij}(k_1,k_2,k_3,k_4)=\delta _{ij}(\widehat{g}_q)^2\left[\frac{p2}{p}+\frac{p1}{p}\left\{\frac{1}{p^{k_1k_2+1}1}+\frac{1}{p^{k_1k_3+1}1}+\frac{1}{p^{k_1k_4+1}1}\right\}\right].$$
(4.7)
This agrees precisely with eq.(4.3) for $`\widehat{g}_q=g_q`$.
In fact, it is possible to give a general argument showing that an amplitude with two external $`\chi `$ fields and $`N`$ external $`\sigma `$ fields for arbitrary $`N`$, computed from the action (3.8), agrees with the corresponding amplitude on a Dirichlet $`q`$-brane. To see this, let us consider the situation where we start with the action (2.1) with $`g`$ replaced by another coupling constant $`\overline{g}`$, and compactify<sup>9</sup><sup>9</sup>9Alternatively, we can work with the uncompactified theory, but just examine those modes of $`\varphi `$ which carry either 0 or $`\pm \sqrt{2}`$ units of momentum in $`(dq1)`$ of the directions. $`(dq1)`$ directions on circles of radii $`1/\sqrt{2}`$. Let $`u^i`$ denote the compact coordinates and $`z^\mu `$ the non-compact ones, and consider an expansion of the field $`\varphi `$ of the form:
$$\varphi (x)=\stackrel{~}{\psi }(z)+\sqrt{\frac{C}{p}}\underset{i=1}{\overset{dq1}{}}\stackrel{~}{\xi }^i(z)\left(\sqrt{2}\mathrm{cos}(\sqrt{2}u^i)\right)+\mathrm{}.$$
(4.8)
We have restricted $`\varphi `$ to be even under $`u^iu^i`$ for each $`i`$; this gives a consistent truncation of the theory at the tree level. The dots stand for higher momentum modes which will not be required for our analysis. Substituting this into (2.1) (with $`g`$ replaced by $`\overline{g}`$) we get the action:
$`{\displaystyle \frac{1}{\overline{g}^2}}{\displaystyle \frac{p^2}{p1}}\left({\displaystyle \frac{2\pi }{\sqrt{2}}}\right)^{dq1}{\displaystyle }d^{q+1}z[{\displaystyle \frac{1}{2}}\stackrel{~}{\psi }p^{\frac{1}{2}\mathrm{}_z}\stackrel{~}{\psi }+{\displaystyle \frac{1}{p+1}}\stackrel{~}{\psi }^{p+1}`$
$`C\{{\displaystyle \frac{1}{2}}\stackrel{~}{\xi }^ip^{\frac{1}{2}\mathrm{}_z}\stackrel{~}{\xi }^i{\displaystyle \frac{1}{2}}\stackrel{~}{\psi }^{p1}\stackrel{~}{\xi }^i\stackrel{~}{\xi }^i\}+𝒪(\stackrel{~}{\xi }^3)+\mathrm{}],`$ (4.9)
If we identify
$$g_q^2=\overline{g}^2\left(\frac{\sqrt{2}}{2\pi }\right)^{dq1},$$
(4.10)
this action looks identical to the one in (3.6) with the fields $`\psi `$, $`\xi ^i`$ replaced by $`\stackrel{~}{\psi }`$, $`\stackrel{~}{\xi }^i`$ and the identification $`x_{}z`$. In particular we can define the analogues of eqs.(3.7)
$$\stackrel{~}{\psi }=1+\frac{g_q\stackrel{~}{\sigma }}{p},\stackrel{~}{\xi }^i=\frac{g_q}{\sqrt{pC}}\stackrel{~}{\chi }^i,$$
(4.11)
and compute the S-matrix elements around the vacuum $`\stackrel{~}{\psi }=1`$ by expanding (4) in a power series in $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{\chi }^i`$. Similarity of (4) and (3.6) (and hence (3.8)) shows that the S-matrix elements computed from the action (4) around the $`\psi =1`$ background are identical to those computed from (3.8) around the $`\stackrel{~}{\psi }=1`$ background. In particular the S-matrix element involving a $`\stackrel{~}{\chi }^i`$, a $`\stackrel{~}{\chi }^j`$, and an arbitrary number of $`\stackrel{~}{\sigma }`$ quanta for (4) is identical to the S-matrix element involving $`\chi ^i`$, $`\chi ^j`$ and an arbitrary number of $`\sigma `$ quanta in (3.8)<sup>10</sup><sup>10</sup>10Since the similarity of (3.6) and (4) holds only to quadratic order in $`\chi `$ ($`\stackrel{~}{\chi }`$), we can only make this claim for two or less external $`\chi `$ ($`\stackrel{~}{\chi }`$) particles..
On the other hand, the S-matrix elements computed from (4) have direct string theory interpretation, as the action is obtained by compactifying a $`p`$-adic string theory. In particular the amplitude $`\stackrel{~}{\chi }^i\stackrel{~}{\chi }^j\stackrel{~}{\sigma }^N`$ is given in terms of correlation functions of $`\stackrel{~}{\chi }^i`$, $`\stackrel{~}{\chi }^j`$ and $`N`$ $`\stackrel{~}{\sigma }`$ vertex operators on the upper half plane. The vertex operator for $`\stackrel{~}{\sigma }`$ is proportional to $`e^{ik.Z}`$, whereas that for $`\stackrel{~}{\chi }^i`$ is given by $`\sqrt{2}\mathrm{cos}(\sqrt{2}U^i)e^{ik.Z}`$. Comparing this with the corresponding computation for the D-$`q`$-brane we see that the $`\stackrel{~}{\sigma }`$ vertex operator is identical to the $`\sigma `$ vertex operator with $`X_{}`$ replaced by $`Z`$. The $`\chi ^i`$ vertex operator on the D-$`q`$-brane, given by $`X_{}^ie^{ik.X_{}}`$, looks different from the $`\stackrel{~}{\chi }^i`$ vertex operator even after we identify $`X_{}`$ with $`Z`$. However if we note that on the boundary of the upper half plane the two point functions $`\sqrt{2}\mathrm{cos}(\sqrt{2}U^i(x_1))\sqrt{2}\mathrm{cos}(\sqrt{2}U^j(x_2))`$ and $`X_{}^i(x_1)X_{}^j(x_2)`$ are identical, both being equal to $`\delta _{ij}|x_1x_2|^2`$, we can conclude that these particular amplitudes in the compactified string theory are indeed identical to those on the D-$`q`$-brane.
To summarise, we have shown that the amplitudes $`\chi ^i\chi ^j\sigma ^N`$ computed from (3.8) are identical to the corresponding amplitudes in the compactified string theory, which in turn are identical to the corresponding ones on the D-$`q`$-brane. This establishes the desired result. In presenting this argument we have not been careful about the overall normalisation factors, but the equality already established for the amplitudes $`\sigma ^N`$ and $`\chi ^i\chi ^j\sigma \sigma `$ in the two theories guarantees that the overall normalisation factors also agree in the two theories.
This provides strong evidence that the solitonic $`q`$-branes of the $`p`$-adic string theory should be identified with Dirichlet $`q`$-branes.
It will be interesting to systematically extend this comparison to S-matrix elements involving more than two external $`\chi ^i`$ states, and also to S-matrix elements involving higher level states. It is not easy to establish this in all generality, however we can consider a subset of the massive modes on the $`q`$-brane and show the agreement between the S-matrix elements on the solitonic $`q`$-brane and D-$`q`$-brane with at most two of these states on the external leg.
We start with the solitonic $`q`$-brane, and consider a generalisation of the expansion (3.4):
$$\varphi (x)=F^{(dq1)}(x_{})\psi (x_{})+\underset{r=1}{\overset{dq1}{}}\underset{\{i_1,\mathrm{}i_r\}}{}^{}_{x_{}^{i_1}}\mathrm{}_{x_{}^{i_r}}F^{(dq1)}(x_{})\xi ^{i_1\mathrm{}i_r}(x_{})+\mathrm{},$$
(4.12)
where $`^{}`$ above denotes sum over those indices $`\{i_1,\mathrm{}i_r\}`$ for which no two in the set are equal. In this case
$$_{x_{}^{i_1}}\mathrm{}_{x_{}^{i_r}}\left(F^{(dq1)}(x_{})\right)^p=p^r\left(F^{(dq1)}(x_{})\right)^{p1}_{x_{}^{i_1}}\mathrm{}_{x_{}^{i_r}}F^{(dq1)}(x_{}).$$
(4.13)
Substituting (4.12) into the equation of motion (2.3), and using eq.(4.13) we get
$`p^{\frac{1}{2}\mathrm{}_{}}\psi `$ $`=`$ $`\psi ^p+𝒪(\xi ^2)`$
$`p^{\frac{1}{2}\mathrm{}_{}}\xi ^{i_1\mathrm{}i_r}`$ $`=`$ $`p^{1r}\psi ^{p1}\xi ^{i_1\mathrm{}i_r}+𝒪(\xi ^2).`$ (4.14)
The action involving these fields is given by
$`S_q(\psi ,\xi )`$ $`=`$ $`{\displaystyle \frac{1}{g_q^2}}{\displaystyle \frac{p^2}{p1}}{\displaystyle }d^{q+1}x_{}[{\displaystyle \frac{1}{2}}\psi p^{\frac{1}{2}\mathrm{}_{}}\psi +{\displaystyle \frac{1}{p+1}}\psi ^{p+1}`$
$`{\displaystyle \underset{r=1}{\overset{dq1}{}}}{\displaystyle \underset{\{i_1,\mathrm{}i_r\}}{}^{}}C_r\{{\displaystyle \frac{1}{2}}\xi ^{i_1\mathrm{}i_r}p^{\frac{1}{2}\mathrm{}_{}+r1}\xi ^{i_1\mathrm{}i_r}{\displaystyle \frac{1}{2}}\psi ^{p1}\xi ^{i_1\mathrm{}i_r}\xi ^{i_1\mathrm{}i_r}\}+𝒪(\xi ^3)].`$
$`C_r`$ is a normalisation constant which can be absorbed into the definition of $`\xi ^{i_1\mathrm{}i_r}`$. From this we see that the mass-shell constraint for $`\xi ^{i_1\mathrm{}i_r}`$, in the $`\psi =1`$ background, is
$$k^2=2(1r).$$
(4.16)
In order to compare this with the world-volume theory on the D-$`q`$-brane, we need to first identify the vertex operator corresponding to the mode $`\xi ^{i_1\mathrm{}i_r}`$. We take this to be
$$V_{i_1\mathrm{}i_r}=X_{}^{i_1}\mathrm{}X_{}^{i_r}e^{ik.X_{}}.$$
(4.17)
This describes a physical state satisfying the same mass shell constraint as eq.(4.16) as long as the indices in the set $`\{i_1,\mathrm{}i_r\}`$ are all different.
We shall now compare the S-matrix elements computed from the action (4) with two or less external $`\xi `$-legs to that computed directly in the D-$`q`$-brane. For two external tachyons and two external $`\xi `$ we can do this explicitly and verify that it agrees with the corresponding computation on the D-$`q`$-brane. The computation is identical to the one discussed earlier. For arbitrary number of external tachyon legs, one can generalise the argument given for external $`\chi ^i`$-legs. The key ingredient of this argument is that the two point function of $`X_{}^{i_1}\mathrm{}X_{}^{i_r}e^{ik.X_{}}`$ and $`X_{}^{i_1^{}}\mathrm{}X_{}^{i_r^{}}e^{ik^{}.X_{}}`$ for the D-$`q`$-brane is identical to that between the vertex operators $`(\sqrt{2})^r\mathrm{cos}(\sqrt{2}U^{i_1})\mathrm{}\mathrm{cos}(\sqrt{2}U^{i_r})e^{ik.Z}`$ and $`(\sqrt{2})^r\mathrm{cos}(\sqrt{2}U^{i_1^{}})\mathrm{}\mathrm{cos}(\sqrt{2}U^{i_r^{}})e^{ik^{}.Z}`$ of the compactified string theory.
## 5 Comments
* We have shown that one can get a consistent truncation of the world-volume theory on a D-$`q`$-brane in $`p`$-adic string theory by keeping only the tachyonic mode. Thus by examining the tree level tachyon amplitudes in the world-volume theory we shall not discover the existence of the other modes. This suggests that there may be other (massless and massive) modes living on the world-volume of the space-filling D-$`(d1)`$-brane as well, inspite of the fact that there are no poles in the tachyon S-matrix elements corresponding to these states. Indeed, ref. attempted to generalise the $`p`$-adic string amplitudes to external vector states. If these modes are present they will give rise to new degrees of freedom on the solitonic $`q`$-brane, and will have to be taken into account in comparing the world-volume theory on the D-$`q`$-brane with that on the solitonic $`q`$-brane.
* It will also be of interest to compute the tension of a Dirichlet $`q`$-brane in the $`p`$-adic string theory independently, and compare with eq.(2.8) describing the tension of a solitonic $`q`$-brane. This will require careful analysis of the cylinder amplitude, and a proper understanding of the closed string sector of the theory.
* It has been shown in ref. that it is possible to assign Chan-Paton factors to the open string states of a $`p`$-padic string theory. This shows the existence of multiple D-$`(d1)`$-branes. Furthermore if there are massless gauge fields in the spectrum of open strings in the $`p`$-adic string theory, and if there is a T-duality transformation relating the D-$`(d1)`$-brane to D-$`q`$-brane, then by switching on Wilson lines corresponding to the gauge fields followed by a T-duality transformation, we can produce static configuration of D-$`q`$-branes separated in space. It will be interesting to examine if the equation of motion (2.3) admit such solutions. Ideas developed in ref. may be useful in this context.
* It is natural to ask if this analysis has any relevance to the ordinary bosonic string theory. Firstly, we would like to point out that even if the tachyon potential in the $`p`$-adic string theory is totally unrelated to that in the ordinary bosonic string theory, it can be regarded as a toy model which nicely illustrates the features expected of the full bosonic string field theory action. Besides, there is evidence of close relationship between tachyon amplitudes in the $`p`$-adic and ordinary bosonic string theory. Thus one might hope that the full tachyon effective action in bosonic string field theory is related in some way to the tachyon effective action in $`p`$-adic string theory.
In this direction, we cannot resist the temptation to point out some apparent similarities between the equation of motion (2.3) and that in the open bosonic string field theory. To lowest order in the level truncation scheme, the tachyon equation of motion in open bosonic string field theory may be written as
$$\left[\left(\alpha ^{}\mathrm{}+1\right)e^{c\alpha ^{}\mathrm{}}2\right]\varphi =\overline{g}\varphi ^2,$$
(5.1)
where $`c=\mathrm{ln}(3^3/4^2)`$, $`\overline{g}`$ is open string coupling constant after suitable normalisation, and $`\varphi `$ is related to the original tachyon field $`T`$ by a field redefinition $`T=e^{c\alpha ^{}\mathrm{}/2}\varphi +\overline{g}^1`$, so that $`\varphi =0`$ is the vacuum without any D-brane, and $`\varphi =1/\overline{g}`$ denotes the D-brane. If we drop the first and the third terms on the left hand side of eq.(5.1) by hand, then this equation, after suitable rescaling of $`x`$ and $`\varphi `$, reduces to eq.(2.3) for $`p=2`$. Of course there is no justification for dropping these terms, so we shall not pursue this matter any further; but it is not inconceivable that some exact relation between bosonic string field theory and $`p`$-adic string theory will be discovered in the future.
Acknowledgement: We wish to thank S. Mukhi for useful discussions. |
warning/0003/cond-mat0003099.html | ar5iv | text | # Exact solution of generalized Schulz-Shastry type models
## I Introduction
Integrability has been considered to be one of the most striking properties of a model for some time. However, the interest was confined to an abstract level since integrable quantum systems typically live in one spatial dimension (but they can be mapped onto classical $`1+1`$ dimensional systems ). In the meantime, one-dimensional systems have become experimental reality (Quantum Hall bars , polymers , charge density wave systems ) and thus integrable models have gained immediate relevance for real physical problems. The Coordinate Bethe Ansatz (CBA) solvability is the first milestone in proving integrability. The (Quantum) Inverse Scattering Method together with the algebraic Bethe Ansatz complete the CBA procedure to construct integrable models. However, the interrelation between solvability by algebraic Bethe ansatz and CBA is still controversial. But solvability by algebraic Bethe Ansatz is a strong hint for solvability by CBA within an appropriately chosen basis. In the case of the quantum impurity problems, for instance, the standard CBA procedure with plane waves fails, in contrast to using a basis consisting of solitary waves . The reason for this is that the bulk theory already describes fully interacting electrons and there is no (free) single particle description of the system. As a result the scattering with the impurity is “diffractive” in the plane-wave basis. Instead, it turns out to be factorizable if the many-electron wave function is written in terms of kinks and anti–kinks.
A breakthrough in the theory of integrable models was the solution of the Hubbard Model (HM) obtained by E. H. Lieb and F. Y. Wu . The HM is a model including on-site Coulomb interaction for electrons moving in an atomic lattice. It is believed to capture important features of high-$`T_C`$ superconductivity and has a metal-insulator transition.
In a recent paper, H. J. Schulz and B. S. Shastry found a new class of solvable one-dimensional Hubbard and XXZ type models. The modification of the original HM and XXZ model consisted in a configuration dependent unitary factor in the hopping term. This can be interpreted as an interaction of the charged particles with a gauge-field, generated by the density of particles. The structure of the unitary factor was $`\mathrm{exp}[\mathrm{i}\widehat{N}]`$, where $`\widehat{N}`$ is a mono-linear functional of particle-number operators; we term such models “single particle correlated hopping ($`1`$-CH)” models.
The idea behind Schulz’ and Shastry’s approach is finding a basis (through a unitary transformation of the original Fock basis) in which the model takes the form of the original Hubbard or XXZ model up to boundary twists which do not affect their solvability . We point out that this is equivalent to equipping the plane waves entering the CBA with phase factors canceling exactly the configuration dependent gauge fields in the hopping term. In the present paper we will generalize such an idea to consider hopping in which $`\widehat{N}`$ is a multilinear functional of particle-number operators. We shall call the resulting models $`n`$–CH models. We will answer the question which $`n`$–CH-Hubbard/XXZ models can be mapped unitarily onto a corresponding uncorrelated but twisted model. This finally proves solvability of the model. Conversely, a non-removable correlated hopping destroys solvability by CBA, since the $`S`$ matrix becomes configuration dependent (see Appendix A).
The paper is laid out as follows. In section II, we will study the most general form of $`1`$–CH HMs, hence including Shastry–Schulz models. Most of the features of the general problem already occur here. Section III accounts for conserved quantum numbers other than the spin orientation of electrons. Multi-chain models where the chains interact with each other exclusively gauge-like, will be considered. The central results obtained in sections II and III will be used to discuss higher correlated hopping in section IV. The $`2`$-CH will be treated in detail, enlightening the approach to general $`n`$-CH. Finally, conclusions will be drawn in section V.
A special class of unitary transformations was used to remove the gauge-like correlation terms from the hopping. In Appendix C the effect of the complementary class of unitary transformations will be elaborated, always restricting on automorphic mappings on the class of twisted $`n`$-CH models of Hubbard- or XXZ-kind. Two propositions needed in this section are proven in Appendix D (see also Ref. ).
## II Single-particle correlated hopping
In this section we discuss a simple generalization of Schulz-Shastry models. Such models have the following Hubbard-type Hamiltonian
$`H`$ $`=`$ $`t{\displaystyle \underset{j,\sigma }{}}\{c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}(\mathrm{i}\gamma _j(\sigma ))\times `$ (3)
$`\times \mathrm{exp}\left[\mathrm{i}{\displaystyle \underset{l}{}}(\alpha _{j,l}(\sigma )n_{l,\sigma }+A_{j,l}(\sigma )n_{l,\sigma })\right]+\mathrm{h}.\mathrm{c}.\}+`$
$`+V{\displaystyle \underset{i}{}}n_{i,}n_{i,}.`$
where $`\{c_{j,\sigma },c_{l,\sigma ^{}}^{}\}=\delta _{\sigma ,\sigma ^{}}\delta _{j,l}`$, $`\{c_{j,\sigma },c_{l,\sigma ^{}}\}=0`$, and $`n_{l,\sigma }:=c_{l,\sigma }^{}c_{l,\sigma }`$. The parameters $`t`$, $`V`$ are the hopping amplitude and the Coulomb repulsion respectively. The class of models (3) is a generalization of Schulz–Shastry models, since a) the parameter $`A`$ occurs, which means correlation between particles with the same spin orientation, and b) the spin- and coordinate dependence of $`A,\alpha ,\gamma `$ is unrestricted here .
We point out that $`A_{j,j+1}`$ can be set to an arbitrary value for all $`j`$, without affecting the physics of the model (since $`n_{l,\sigma }\{0,1\}c_{j+1,\sigma }^{}n_{j+1,\sigma }0`$). We name parameters like $`A_{j,j+1}`$ “irrelevant”. A similar argument holds for the parameter $`A_{j,j}`$: contributions from $`n_{j,\sigma }`$ arise only if $`n_{j,\sigma }=1`$ because of $`c_{j,\sigma }n_{j,\sigma }c_{j,\sigma }`$. Hence, this term can be included in the parameter $`\gamma _j(\sigma )`$. Parameters like $`A_{j,j}`$ will be called “subrelevant” throughout this paper. Irrelevant as well as subrelevant parameters appear as soon as phase–correlations among particles having the same spin orientation as the hopping particle are involved.
It is worthwhile noting that Hamiltonian (3) is not diagonalizable by direct CBA since the scattering matrix is configuration dependent (see Appendix A). This destroys the factorizability of a many-particle $`S`$ matrix into two-particle $`S`$ matrices. Thus, we first remove the phases in the hopping term of (3) by a unitary transformation and then, we can diagonalize the transformed Hamiltonian by CBA.
The ansatz for the unitary transformation is achieved through the operator
$$U:=\mathrm{exp}\left[\mathrm{i}(\xi _{l,m}^{\mu ,\nu }n_{l,\mu }n_{m,\nu }+\zeta _{l,\mu }n_{l,\mu })\right]=:\mathrm{exp}(\mathrm{i}S),$$
(4)
(we use the sum convention) where $`\xi _{i,j},\zeta _{l,m}\text{ }\mathrm{R}`$ are unknown variables which have to be fixed for cancelling the unitary prefactor in the hopping term of (3). Since an antisymmetric part in the parameter $`\xi _{i,j}`$ vanishes after summation, it can be defined fully symmetric: $`\xi _{l,m}^{\mu ,\nu }=\xi _{m,l}^{\nu ,\mu }`$. We can further choose $`\xi _{m,m}^{\mu ,\mu }=0`$ (a non-zero $`\xi _{m,m}^{\mu ,\mu }`$ can be included in the parameter $`\zeta _{l,m}`$).
We locally transform the Hamiltonian by $`U`$: $`c_{j,\sigma }\stackrel{U}{}Uc_{j,\sigma }U^1.`$ Number operators remain unchanged, but the hopping term is altered
$`c_{j+1,\sigma }^{}c_{j,\sigma }`$ $`\stackrel{U}{}`$ $`c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}\left[2\mathrm{i}\left(\xi _{j+1,m}^{\sigma ,\mu }\xi _{j,m}^{\sigma ,\mu }\right)n_{m,\mu }\right]`$
$`\mathrm{exp}\left[\mathrm{i}\left(\zeta _{j+1,\sigma }\zeta _{j,\sigma }2\xi _{j,j+1}^{\sigma ,\sigma }\right)\right].`$
The unitary factors in the hopping term of (3) are removed if
$`2(\xi _{j,m}^{\sigma ,\sigma }\xi _{j+1,m}^{\sigma ,\sigma })`$ $`\stackrel{!}{=}`$ $`\alpha _{j,m}(\sigma )`$ (5)
$`(\zeta _{j,\sigma }\zeta _{j+1,\sigma }+2\xi _{j,j+1}^{\sigma ,\sigma })`$ $`\stackrel{!}{=}`$ $`\gamma _j(\sigma ),`$ (6)
and
$$2(\xi _{j,m}^{\sigma ,\sigma }\xi _{j+1,m}^{\sigma ,\sigma })\stackrel{!}{=}A_{j,m}(\sigma );m\{1,\mathrm{},L\}\{j,j+1\}.$$
(7)
For periodic boundary conditions (PBC), Eqs. (5)–(7) for $`j=L`$ represent the jump across the boundary. Admitting boundary phases in company with the boundary’s crossing, Eqs. (5)–(7) are modified as
$`2(\xi _{L,m}^{\sigma ,\sigma }\xi _{1,m}^{\sigma ,\sigma })`$ $`\stackrel{!}{=}`$ $`\alpha _{L,m}(\sigma )\varphi _{}^{(1)}(\sigma ),`$ (8)
$`(\zeta _{L,\sigma }\zeta _{1,\sigma }+2\xi _{L,1}^{\sigma ,\sigma })`$ $`\stackrel{!}{=}`$ $`\gamma _L(\sigma )\varphi (\sigma ).`$ (9)
and
$$2(\xi _{L,m}^{\sigma ,\sigma }\xi _{1,m}^{\sigma ,\sigma })\stackrel{!}{=}A_{L,m}(\sigma )\varphi _{}^{(1)}(\sigma );m\{2,\mathrm{},L1\},$$
(10)
In fact, the relations (5)–(7) and (8)–(10) constitute a system of recursive relations for $`\xi _{j,m}^{\sigma ,\sigma ^{}}`$ and $`\zeta _{j,\sigma }`$.
We will now discuss the exclusions in Eq. (7). The corresponding part of the transformed hopping term for $`m=j`$ is (since $`\xi _{j,j}^{\sigma ,\sigma }=0`$) $`c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}\left[\mathrm{i}\left(2\xi _{j+1,j}^{\sigma ,\sigma }+A_{j,j}(\sigma )\right)n_{j,\sigma }\right].`$ This term is non-zero only if $`n_{j,\sigma }=1`$. Hence, this “correlation” factor is equivalent to $`c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}\left[\mathrm{i}\left(2\xi _{j+1,j}^{\sigma ,\sigma }+A_{j,j}(\sigma )\right)\right].`$ For $`m=j+1`$ it is $`c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}\left[\mathrm{i}\left(2\xi _{j,j+1}^{\sigma ,\sigma }A_{j,j+1}(\sigma )\right)n_{j+1,\sigma }\right].`$ This term is non-zero only if $`n_{j+1,\sigma }=0`$. Hence, this “correlation” factor is equivalent to $`c_{j+1,\sigma }^{}c_{j,\sigma }`$ irrespective of what value $`A_{j,j+1}(\sigma )2\xi _{j,j+1}^{\sigma ,\sigma }`$ may take. As a consequence, there is no condition for $`m=j+1`$ in Eq. (7), and $`A_{j,j}(\sigma )+2\xi _{j,j+1}^{\sigma ,\sigma }`$ enters as a modification of Eq. (6) for $`\gamma _j`$ (see Appendix B).
The system of equations (5)–(6) cannot be solved for arbitrary $`\alpha `$ and $`A`$ (the number of parameters on the right-hand side is larger than on the left-hand side). Together with the symmetry of $`\xi _{i,j}`$, Eqs. (5) and (7) define the effect of an increase in both site indices of $`\xi _{i,j}`$; namely an increase of $`\xi `$ by the related $`\alpha `$ respectively $`A`$. Starting from an initial parameter, say $`\xi _{1,1}:=0`$, every $`\xi _{i,j}`$ is then defined by passing from $`(1,1)`$ to $`(i,j)`$ and summing up the contributions from the recursive relation. $`\xi _{i,j}`$ is well defined if this procedure is path independent. This is equivalent to demanding that contributions from closed loops in the $`(i,j)`$ plane vanish. To verify this, it is sufficient facing the smallest possible loops: $`(j,m)(j+1,m)(j+1,m+1)(j,m+1)(j,m)`$ (which we will henceforth call “elementary”). Applying Eq. (5) respectively Eq. (7), we obtain
$$\alpha _{j,m+1}(\sigma )\alpha _{j,m}(\sigma )=\alpha _{m,j+1}(\sigma )\alpha _{m,j}(\sigma ),$$
(11)
and
$$A_{j,m+1}(\sigma )A_{j,m}(\sigma )=A_{m,j+1}(\sigma )A_{m,j}(\sigma )$$
(12)
for $`mj,j\pm 1`$. We call the conditions (11), (12)) “closedness conditions”. The recursive relation and the closedness condition can be written in a more compact and clear form in terms of the discrete gradient, defined by $`_xf(x):=f(x+1)f(x)`$. The recursive relations (5) and (7) then read
$`2_x\xi ^{\sigma ,\sigma }(x,y)2_x\xi ^{\sigma ,\sigma }(y,x)`$ $`=`$ $`\alpha (x,y;\sigma )`$ (13)
$`2_x\xi ^{\sigma ,\sigma }(x,y)2_x\xi ^{\sigma ,\sigma }(y,x)`$ $`=`$ $`A(x,y;\sigma )`$ (14)
and the closedness conditions take the form
$`_y\alpha (x,y;\sigma )`$ $`=`$ $`_x\alpha (y,x;\sigma )`$ (15)
$`_yA(x,y;\sigma )`$ $`=`$ $`_xA(y,x;\sigma );xy.`$ (16)
With these conditions being fulfilled , the correlations from the hopping term can be removed and the rotated model is finally known to be solvable by CBA.
For open boundary conditions the correlated hopping can be “gauged away” completely, yielding the HM without any boundary phases.
Instead, PBC lead to the HM with twisted boundary conditions. Periodicity implies that the parameters $`\xi _{i,j}`$, $`\alpha `$, $`A`$ and $`\gamma `$ are periodic in their site-indices with period $`L`$. The boundary phase is determined by hopping from site $`L`$ to site $`L+1\widehat{=}\mathrm{\hspace{0.17em}1}`$; such a phase enters as written in Eqs. (8)–(10) for the relations $`\xi _{1,m}\xi _{L,m}`$ and $`\zeta _{1,m}\zeta _{L,m}`$ and then for $`A_{1,m}A_{L,m}`$, $`\alpha _{1,m}\alpha _{L,m}`$. Equations (5)–(10) lead to
$`2(\xi _{1,m}^{\sigma ,\sigma }\xi _{L,m}^{\sigma ,\sigma })`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L1}{}}}\alpha _{j,m}(\sigma )\stackrel{!}{=}`$ (18)
$`\stackrel{!}{=}\alpha _{L,m}(\sigma )+\varphi _{}^{(1)}(\sigma ),`$
$`2(\xi _{1,m}^{\sigma ,\sigma }\xi _{L,m}^{\sigma ,\sigma })`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jm,m1}}{\overset{L1}{}}}A_{j,m}(\sigma )+`$ (21)
$`+2(\xi _{m1,m}^{\sigma ,\sigma }\xi _{m,m+1}^{\sigma ,\sigma })`$
$`\stackrel{!}{=}A_{L,m}(\sigma )+\varphi _{}^{(1)}(\sigma ),`$
$`\zeta _{1,\sigma }\zeta _{L,\sigma }`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L1}{}}}(\gamma _j(\sigma )+A_{j,j})`$ (22)
$`\stackrel{!}{=}`$ $`\gamma _L(\sigma )A_{L,L}+\varphi (\sigma ),`$ (23)
where $`\varphi `$ denotes the boundary phases. They can be determined (without solving for the $`\xi `$’s explicitly) from Eqs. (18) to (23) as
$`\varphi _{}^{(1)}(\sigma )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\alpha _{j,m}(\sigma ),`$ (24)
$`\varphi _{}^{(1)}(\sigma )`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jm1,m}}{\overset{L}{}}}A_{j,m}(\sigma )+`$ (26)
$`+A_{m,m1}(\sigma )+A_{m1,m+1}(\sigma ),`$
$`\varphi (\sigma )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\left(\gamma _j(\sigma )+A_{j,j}(\sigma )\right).`$ (27)
Eqs. (11) and (12) also ensure that the phases in Eqs. (24) and (26) are $`m`$-independent. Summarizing, iff Eqs. (11) and (12) are fulfilled, the Hamiltonian (3) can be mapped by $`U`$ (Eq. (4)) onto the usual HM with modified boundary conditions. The boundary twists are given by
$$\mathrm{\Phi }_\sigma :=\varphi (\sigma )+\varphi _{}^{(1)}(\sigma )N_\sigma +\varphi _{}^{(1)}(\sigma )(N_\sigma 1).$$
(28)
We emphasize that here, the factor $`(N_\sigma 1)`$ appears instead of $`N_\sigma `$. The reason is that the recursive relation (7) does not exist for $`m=j,j+1`$. But since one particle with spin $`\sigma `$ has to be on site $`j`$ and none on site $`j+1`$ (or vice versa: see the corresponding hopping term in (3)), one particle less accounts for the phase.
### A Translational invariant models
In this section we assume a translational invariant model, and hence translational invariance of the parameters $`\alpha `$ and $`A`$: $`\alpha _{j,l}(\sigma )=\alpha _{jl}(\sigma )`$, $`A_{j,l}(\sigma )=A_{jl}(\sigma )`$.
Restricting ourselves to translational invariant unitary transformations ($`\xi _{j,k}=\xi _{jk}`$), the former elementary loops become multiple loops (see Fig. 2). Then, the elementary loops are $`(j,m)(j+1,m)(j+1,m+1)`$ and the closedness condition only consists of the following terms from Eqs. (11) and (12) $`\alpha _j(\sigma )=\alpha _{(j+1)}(\sigma )`$ for all $`j`$, $`A_j(\sigma )=A_{(j+1)}(\sigma )`$ for all $`j0`$. Also here, for the parameter $`A`$ not all elementary loops are viable. But the only exceptional loop emerges from the case corresponding to $`m=j`$ in Eq. (12) (that is a square on the diagonal – the fat dotted line in Fig. 1). But as already mentioned, the (excluded) condition for $`m=j`$ in Eq. (12) is trivially fulfilled. This is confirmed by calculating the effect of the jump from $`0`$ to $`2`$ in the exceptional loop in Fig. 1, which here constitutes an elementary loop (comparing with Fig. 2 it turns out that the points $`0`$ and $`2`$ are identified since translational invariant parameters are assumed). As a consequence, the feature of additional phases imported by subrelevant parameters and hence non-viable loops are absent here.
But we want to stress that even here Eqs. (11), (12) are necessary and sufficient conditions for solvability – for periodic as well as for open chains. This is because a translational invariant model need not necessarily be gauged away by a unitary transformation with translational invariant parameters.
Now, we can connect to Ref. . We obtain the boundary phases in Ref. through the identifications $`\alpha _{j,m}\alpha _{jm}`$ and $`\alpha _j(\sigma ):=\sigma \alpha _j`$. In addition, $`A_{j,m}(\sigma )=\gamma _j(\sigma )0`$ and $`\beta _{i,j}=\sigma /2\xi _{i,j}^{\sigma ,\sigma }`$, giving an antisymmetric $`\beta _{i,j}`$. However, in order to guarantee closedness and hence solvability of the recursive relations,
$$\alpha _{jm}\alpha _{jm1}=(\alpha _{mj}\alpha _{mj1})$$
has to hold.
A special type of such models is studied in Ref. with the only non-zero values $`\alpha _0=\alpha _1=\eta `$, where $`\eta `$ is a real parameter. The closedness conditions are all fulfilled and therefore the correlation factors in the hopping can be gauged away giving exactly the phases found in Ref. . However, we could not reproduce the results following the path suggested in Ref. (see the general discussion in appendix A).
## III Multi-chain models
Multi-chain models can be treated analogously. The chain-variable plays the role of the spin-variable now. The only difference is that in general the chain variable ranges in a set different from $`\{1/2,+1/2\}`$. Since the method we employ to diagonalize the Hamiltonian stays the same, we only state the results. For multi-chain models we assume intra-chain hopping only.
At first, we deal with spin-less fermions and let the hopping-term of (3) take the form
$$a_{c,j+1}^{}a_{c,j}\mathrm{exp}(\mathrm{i}\gamma _j(c))\mathrm{exp}\left[\mathrm{i}A_{j,l}^{c,d}n_{d,l}\right]+\mathrm{h}.\mathrm{c}.,$$
(29)
where $`c`$ and $`d`$ are chain indices, which run in $`\{1,\mathrm{},C\}`$ with $`C`$ being the number of chains . The unitary transformation will be
$$U:=\mathrm{exp}\left[\mathrm{i}(\xi _{l,m}^{\mu ,\nu }n_{\mu ,l}n_{\nu ,m}+\zeta _{\mu ,l}n_{\mu ,l})\right],$$
(30)
where $`\mu ,\nu `$ are now site indices (instead of spin indices of Sec. II), running in $`\{1,\mathrm{},C\}`$ (instead of $`\{1/2,1/2\}`$). In this case, the recursive relations are
$`2(\xi _{j,m}^{c,d}\xi _{j+1,m}^{c,d})`$ $`\stackrel{!}{=}`$ $`A_{j,m}^{c,d};cd,`$ (31)
$`2(\xi _{j,m}^{c,c}\xi _{j+1,m}^{c,c})`$ $`\stackrel{!}{=}`$ $`A_{j,m}^{c,c};mj,j+1,`$ (32)
$`(\zeta _{c,j}\zeta _{c,j+1}+2\xi _{j,j+1}^{c,c})`$ $`\stackrel{!}{=}`$ $`\gamma _{c,j}.`$ (33)
The closedness conditions read
$`A_{j,m+1}^{c,d}A_{j,m}^{c,d}`$ $`=`$ $`A_{m,j+1}^{d,c}A_{m,j}^{d,c};(cd)mj\pm 1,`$ (34)
$`A_j^{c,d}`$ $`=`$ $`A_{(j+1)}^{d,c};cdj0,1`$ (35)
for all $`c,d,j,m`$ respecting the noted exceptions. Equation (35) applies to the translational invariant ansatz. The boundary phases finally become
$$\mathrm{\Phi }_c=\varphi _c+\underset{d=1}{\overset{C}{}}\varphi _c^{(1)}(d)(N_d\delta _{c,d}),$$
(36)
where
$`\varphi _c^{(1)}(d)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}A_{j,m}^{c,d};cd,`$ (37)
$`\varphi _c^{(1)}(c)`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jm1,m}}{\overset{L}{}}}A_{j,m}^{c,c}+A_{m,m1}^{c,c}+A_{m1,m+1}^{c,c},`$ (38)
$`\varphi _c`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\left(\gamma _{c,j}+A_{j,j}^{c,c}\right).`$ (39)
To compare this with Ref. , one has to replace $`\xi _{l,m}^{\mu ,\nu }`$ with $`\frac{1}{2}\xi _{\mu ,\nu }B_{\mu ,\nu }(lm)`$ and additionally $`A_{j,m}^{c,d}`$ with $`\xi _{c,d}A_{c,d}(jm)`$. There, $`\xi _{\mu ,\nu }`$ was antisymmetric and hence, $`B_{\mu ,\nu }(lm)`$ has to be antisymmetric, too: $`B_{\mu ,\nu }(m)=B_{\nu ,\mu }(m)`$. The solvability conditions (35) then transport into
$`\begin{array}{c}\xi _{c,d}A_{c,d}(jm)=\xi _{d,c}A_{d,c}(mj1)\\ \\ A_{c,d}(jm)=A_{d,c}(mj1)\end{array}`$
for all $`m,j,c,d`$. This is equivalent to $`A_{c,d}(m)=A_{d,c}(m1)`$ for all $`m,c,d`$ as obtained in Ref. . This condition is sufficient for CBA solvability but not necessary. Necessary and sufficient is $`A_{c,d}(m)A_{d,c}(m1)A_{c,d}(m1)+A_{d,c}(m)=0`$ for all $`m`$.
If spin has to be included, nothing changes but the number of indices. The hopping term becomes
$$a_{c,j+1;\sigma }^{}a_{c,j;\sigma }\mathrm{exp}(\mathrm{i}\gamma _j(c))\mathrm{exp}\left[\mathrm{i}A_{j,l}^{c,d}(\sigma ,\rho )n_{d,l;\rho }\right]+\mathrm{h}.\mathrm{c}.,$$
(40)
where $`c,d`$ are the chain indices again and $`\sigma ,\rho `$ are spin indices. The unitary transformation takes the form
$$U:=\mathrm{exp}\left[\mathrm{i}(\xi _{l,m}^{\mu ,\nu }(\rho ,\tau )n_{\mu ,l;\rho }n_{\nu ,m;\tau }+\zeta _{\mu ,l;\rho }n_{\mu ,l;\rho })\right],$$
(41)
where $`\mu ,\nu `$ are the chain indices and $`\rho ,\tau `$ are spin indices. The recursion relations read
$`2(\xi _{j,m}^{c,d}(\sigma ,\rho )\xi _{j+1,m}^{c,d}(\sigma ,\rho ))`$ $`\stackrel{!}{=}`$ $`A_{j,m}^{c,d}(\sigma ,\rho );`$ (42)
$`cd\sigma \rho m`$ $``$ $`j,j+1,`$ (43)
$`(\zeta _{c,j}(\sigma )\zeta _{c,j+1}(\sigma )+2\xi _{j,j+1}^{c,c}(\sigma ,\sigma ))`$ $`\stackrel{!}{=}`$ $`\gamma _{c,j}(\sigma ),`$ (44)
and the closedness conditions are
$`A_{j,m+1}^{c,d}(\sigma ,\rho )A_{j,m}^{c,d}(\sigma ,\rho )`$ $`=`$ $`A_{m,j+1}^{d,c}(\rho ,\sigma )A_{m,j}^{d,c}(\rho ,\sigma );`$ (45)
$`cd\sigma \rho m`$ $``$ $`j\pm 1,`$ (46)
for all $`c,d,j,m,\sigma ,\rho `$ respecting the exceptions . The boundary phases become
$$\mathrm{\Phi }_{c,\sigma }=\varphi _{c,\sigma }+\underset{d=1}{\overset{C}{}}\underset{\rho \{,\}}{}\varphi _{c,\sigma }^{(1)}(d,\rho )(N_{d,\rho }\delta _{c,d}\delta _{\rho ,\sigma }),$$
(47)
where
$`\varphi _{c,\sigma }^{(1)}(d,\rho )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}A_{j,m}^{c,d}(\sigma ,\rho );cd\sigma \rho `$ (48)
$`\varphi _{c,\sigma }^{(1)}(c,\sigma )`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jm1,m}}{\overset{L}{}}}A_{j,m}^{c,c}(\sigma ,\sigma )+`$ (50)
$`+A_{m,m1}^{c,c}(\sigma ,\sigma )+A_{m1,m+1}^{c,c}(\sigma ,\sigma ),`$
$`\varphi _{c,\sigma }`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\left(\gamma _{c,j}(\sigma )+A_{j,j}^{c,c}(\sigma ,\sigma )\right).`$ (51)
The overlap of Ref. with the work presented in the Refs. was discussed by the authors and was argued to not exist. Now it is shown that the correlation in Refs. cannot be gauged away, since the correlation parameters violate Eq. (45). However, the Hamiltonian having the slightly but essentially different phases
$$\begin{array}{c}\mathrm{exp}\{\mathrm{i}\pi \phi _\sigma [n_{l+1,m,\sigma }n_{l1,m,\sigma }+\hfill \\ +n_{l+1,m+1,\sigma }n_{l1,m+1,\sigma }]+\frac{2\pi \gamma _\sigma }{N_a}\}\hfill \end{array}$$
(52)
instead of
$$\mathrm{exp}\left\{\mathrm{i}\pi \phi _\sigma \left[n_{l+1,m,\sigma }n_{l1,m,\sigma }\right]+\frac{2\pi \gamma _\sigma }{N_a}\right\}$$
(53)
yields the BE found in Reference.
The models studied in Ref. are equivalent to multi-chain models of spin-less fermions. In the first model, the only non-zero model parameters are $`A_{j,j}^{m,m+1}=A_{j,j+1}^{m+1,m}=4\mathrm{\Theta }_{m,m+1}`$ and $`\gamma _j(m)=2(\mathrm{\Theta }_{m+1,m}\mathrm{\Theta }_{m,m1})`$, where we used the notation in Ref. : $`j`$ and $`m`$ are the site and chain index respectively. The second model is represented by the parameters $`A_{j,j+1}^{m+1,m}=A_{j,j}^{m,m+1}=\mathrm{\Theta }_{m,m+1}+\alpha _{m,m+1}`$ and $`A_{j,j}^{m+1,m}=A_{j,j+1}^{m,m+1}=\mathrm{\Theta }_{m,m+1}\alpha _{m,m+1}`$. The closedness conditions are fulfilled iff $`\mathrm{\Theta }_{m,m+1}=\mathrm{\Theta }_{m+1,m}`$ and $`\alpha _{m,m+1}=\alpha _{m+1,m}`$.
## IV Two-particle and higher correlated hopping
The procedure developed in the previous section can be extended to consider higher (than one-particle) correlated hopping. First, we face explicitly 2–particle correlated hopping (2–CH). Then we will sketch how to deal with the general case of $`n`$–CH.
2–CH corresponds to the occurrence of a term like $`\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma )n_{l,\lambda }n_{m,\mu }`$ in the exponential factor of the hopping term of (3). The 2-CH Hubbard-type Hamiltonian is
$`H`$ $`=`$ $`t{\displaystyle \underset{j,\sigma }{}}\{c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}(\mathrm{i}\gamma _j(\sigma ))\times `$ (56)
$`\times \mathrm{exp}\left[\mathrm{i}\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma )n_{l,\lambda }n_{m,\mu }\right]\times `$
$`\times \mathrm{exp}\left[\mathrm{i}\alpha _m^\mu (j,\sigma )n_{m,\mu }\right]+\mathrm{h}.\mathrm{c}.\}+V{\displaystyle }_in_{i,}n_{i,}.`$
Without loss of generality the parameters $`\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma )`$ can be chosen symmetric in the index pairs $`(l,\lambda )`$ and $`(m,\mu )`$ and vanishing if these index pairs coincide (see Ref. ).
The parameters $`\overline{\alpha }_{j+1,m}^{\sigma ,\mu }(j,\sigma )`$ and $`\alpha _{j+1}^\sigma (j,\sigma )`$ are irrelevant for all $`j,m,\sigma ,\mu `$; the effect of the subrelevant parameters on the lower correlated ones will be discussed later on in the present section.
We first remove the phases in the hopping term of (56) by a unitary transformation. Then, we diagonalize the transformed Hamiltonian by CBA in computing the boundary phases. The 2–CH demands an exponent $`\overline{\xi }_{l,m,r}^{\lambda ,\mu ,\rho }n_{l,\lambda }n_{m,\mu }n_{n,\rho }`$ in the unitary transformation $`U`$:
$`U`$ $`:=`$ $`\mathrm{exp}[\mathrm{i}(\overline{\xi }_{l,m,r}^{\lambda ,\mu ,\rho }n_{l,\lambda }n_{m,\mu }n_{r,\rho }+`$ (58)
$`+\xi _{l,m}^{\lambda ,\mu }n_{l,\lambda }n_{m,\mu }+\zeta _{m,\mu }n_{m,\mu })].`$
Both $`\xi `$ and $`\overline{\xi }`$ are totally symmetric and vanish if any two pairs of parameters coincide.
The hopping term is transformed into
$`c_{j+1,\sigma }^{}c_{j,\sigma }`$ $`\stackrel{U}{}`$ $`c_{j+1,\sigma }^{}c_{j,\sigma }\times `$
$`\times \mathrm{exp}[3\mathrm{i}(\overline{\xi }_{j+1,l,m}^{\sigma ,\lambda ,\mu }\overline{\xi }_{j,l,m}^{\sigma ,\lambda ,\mu })n_{l,\lambda }n_{m,\mu }+`$
$`+2\mathrm{i}\left(\xi _{j+1,m}^{\sigma ,\mu }\xi _{j,m}^{\sigma ,\mu }3\overline{\xi }_{j,j+1,m}^{\sigma ,\sigma ,\mu }\right)n_{m,\mu }+`$
$`+\mathrm{i}(\zeta _{j+1,\sigma }\zeta _{j,\sigma }2\xi _{j,j+1}^{\sigma ,\sigma })],`$
whereas the Coulomb interaction term remains unchanged. This leads to the recursive relations (compare with (5)– (7))
$`3(\overline{\xi }_{j,l,m}^{\sigma ,\lambda ,\mu }\overline{\xi }_{j+1,l,m}^{\sigma ,\lambda ,\mu })`$ $`\stackrel{!}{=}`$ $`\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma ),`$ (59)
$`(ml\mu \lambda )\text{and}(l,mj,j+1`$ $``$ $`\lambda ,\mu \sigma )`$ (60)
$`2(\xi _{j,m}^{\sigma ,\mu }\xi _{j+1,m}^{\sigma ,\mu }+3\overline{\xi }_{j,j+1,m}^{\sigma ,\sigma ,\mu })`$ $`\stackrel{!}{=}`$ $`\alpha _m^\mu (j,\sigma );`$ (61)
$`mj,j+1\mu `$ $``$ $`\sigma `$ (62)
$`(\zeta _{j,\sigma }\zeta _{j+1,\sigma }+2\xi _{j,j+1}^{\sigma ,\sigma })`$ $`\stackrel{!}{=}`$ $`\gamma _j(\sigma )`$ (63)
for the parameters in $`U`$.
We point out that, in the present case, two kinds of elementary loops exist because of the variety of indices in $`\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma )`$. Namely: $`(j;l,m)(j+1;l,m)(j+1;l+1,m)(j;l+1,m)(j;l,m)`$ and $`(j;l,m)(j+1;l,m)(j+1;l,m+1)(j;l,m+1)(j;l,m)`$. However, due to the symmetry of the $`\overline{\alpha }`$ both loops give the same closedness condition for $`\overline{\alpha }`$, which is
$$\overline{\alpha }_{l+1,m}^{\lambda ,\mu }(j,\sigma )\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma )=\overline{\alpha }_{j+1,m}^{\sigma ,\mu }(l,\lambda )\overline{\alpha }_{j,m}^{\sigma ,\mu }(l,\lambda )$$
(64)
for $`lj,j\pm 1\lambda \sigma `$ and $`mj,j+1(l,l+1)\mu \lambda `$.
The corresponding boundary phases are
$`\varphi _\sigma ^{(2)}(\lambda ,\mu )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\overline{\alpha }_{l,m}^{\lambda ,\mu }(j,\sigma );\lambda ,\mu \sigma ,`$ (65)
$`\varphi _\sigma ^{(2)}(\sigma ,\mu )`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jl,l1}}{\overset{L}{}}}\overline{\alpha }_{l,m}^{\sigma ,\mu }(j,\sigma )+`$ (68)
$`+\overline{\alpha }_{l1,m}^{\sigma ,\mu }(l,\sigma )+\overline{\alpha }_{l+1,m}^{\sigma ,\mu }(l1,\sigma );`$
$`\mu \sigma ,`$
$`\varphi _\sigma ^{(2)}(\sigma ,\sigma )`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jl,l\pm 1}}{\overset{L}{}}}\overline{\alpha }_{l,l+1}^{\sigma ,\sigma }(j,\sigma )+`$ (71)
$`+\overline{\alpha }_{l1,l}^{\sigma ,\sigma }(l+1,\sigma )+\overline{\alpha }_{l1,l+2}^{\sigma ,\sigma }(l,\sigma )+`$
$`+\overline{\alpha }_{l+1,l+2}^{\sigma ,\sigma }(l1,\sigma ).`$
where $`lm,m\pm 1`$ can be chosen arbitrarily. The result turns out to be independent of this choice. For less than three sites , Eq. (71) is ill-defined. It reflects a physical limitation: for the boundary phase $`\varphi _\sigma ^{(2)}(\sigma ,\sigma )`$ to occur, at least three particles with the same spin orientation have to exist; this is possible only if at least three sites are available.
Now we will discuss the effect of the subrelevant parameters $`\overline{\alpha }_{j,m}^{\sigma ,\mu }(j,\sigma )`$ and $`\alpha _j^\sigma (j,\sigma )`$. The 2-CH subrelevant part of the exponent in the hopping term is $`2(\overline{\alpha }_{j,m}^{\sigma ,\mu }(j,\sigma )+3\overline{\xi }_{j,j+1,m}^{\sigma ,\sigma ,\mu })`$. As discussed in the previous section, this term does not vanish in general because the recursive relations do not cover the index grid of $`\overline{\xi }`$ completely. It contributes to the 1-CH part instead. It has to be added to the right-hand side of Eq. (61). As a consequence, the parameter $`\overline{\xi }`$ drops out and the recursive relation reads
$$2(\xi _{j,m}^{\sigma ,\mu }\xi _{j+1,m}^{\sigma ,\mu })\stackrel{!}{=}\beta _m^\mu (j,\sigma );mj,j+1\mu \sigma ,$$
(72)
where $`\beta _m^\mu (j,\sigma )`$ is defined as
$$\beta _m^\mu (j,\sigma ):=\alpha _m^\mu (j,\sigma )+2\overline{\alpha }_{j,m}^{\sigma ,\mu }(j,\sigma ).$$
(73)
Doing the same for the 1-CH subrelevant part in the hopping (concerning Eqs. (63) and (72) now), $`\zeta `$ drops out and in Eq. (63), $`\gamma _j(\sigma )`$ will be substituted by $`\stackrel{~}{\gamma }_j(\sigma )`$
$`\stackrel{~}{\gamma }_j(\sigma )`$ $`:=`$ $`\gamma _j(\sigma )+\alpha _j^\sigma (j,\sigma )+2\overline{\alpha }_{j,j}^{\sigma ,\sigma }(j,\sigma )`$ (74)
$`=`$ $`\gamma _j(\sigma )+\alpha _j^\sigma (j,\sigma ).`$ (75)
Using Eq. (72), the second set of closedness conditions are obtained
$$\beta _{m+1}^\mu (j,\sigma )\beta _m^\mu (j,\sigma )=\beta _{j+1}^\sigma (m,\mu )\beta _j^\sigma (m,\mu )$$
(76)
and the boundary phases are
$`\mathrm{\Phi }_\sigma `$ $`=`$ $`\varphi _\sigma +{\displaystyle \underset{\lambda }{}}\varphi _\sigma ^{(1)}(\lambda )(N_\lambda \delta _{\lambda ,\sigma })+`$ (78)
$`+{\displaystyle \underset{\lambda ,\mu }{}}\varphi _\sigma ^{(2)}(\lambda ,\mu )(N_\lambda \delta _{\lambda ,\sigma })(N_\mu \delta _{\mu ,\sigma }),`$
where
$`\varphi _\sigma ^{(1)}(\mu )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\beta _m^\mu (j,\sigma );\mu \sigma ,`$ (79)
$`\varphi _\sigma ^{(1)}(\sigma )`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{jm1,m}}{\overset{L}{}}}\beta _m^\sigma (j,\sigma )+`$ (81)
$`\beta _{m1}^\sigma (m,\sigma )+\beta _{m+1}^\sigma (m1,\sigma ),`$
$`\varphi _\sigma `$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\stackrel{~}{\gamma }_j(\sigma ).`$ (82)
The analysis presented above can be generalized to consider $`n`$-CH.
Let $`\stackrel{(n)}{\alpha },\stackrel{(n)}{\beta }`$ be the parameters of the $`n`$-CH part of the hopping:
$`\stackrel{(n)}{\beta }_{\{\}}^{\{𝒮\}}(j,\sigma )`$ $`:=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+k}{k}}\right)\stackrel{(n+k)}{\alpha }_{j,\mathrm{},j,\{\}}^{\sigma ,\mathrm{},\sigma ,\{𝒮\}}(j,\sigma )`$ (83)
$`=`$ $`\stackrel{(n)}{\alpha }_{\{\}}^{\{𝒮\}}(j,\sigma )+(n+1)\stackrel{(n+1)}{\alpha }_{j\{\}}^{\sigma \{𝒮\}}(j,\sigma ),`$ (84)
$`\stackrel{(n)}{\beta }_{j,\mathrm{}}^{\sigma ,\mathrm{}}(j,\sigma )`$ $`:=`$ $`\stackrel{(n)}{\beta }_{j,j,\mathrm{}}^{\sigma ,\sigma ,\mathrm{}}(.,.):=0.`$ (85)
(for multi-chain models the spin index is a multi index spin/chain). The dots ($`\mathrm{}`$) stand for an arbitrary series of indices except they appear in between two equal indices ($`\sigma \mathrm{}\sigma `$). In this case, the void is meant to be filled up with $`k`$ times $`\sigma `$. $`\{𝒮\}`$ is a set of spin- or color indices, $`\{\}`$ a set of coordinate indices. The sum in Eq. (83) of course is finite. The variables are confined since $`CL1`$ is the highest possible correlation level if $`C`$ is the number of inner degrees of freedom ($`C=2`$ for spin $`1/2`$). Hence, $`kCL1n`$ and for $`\overline{n}`$-CH with $`\overline{n}CL1`$, the sum in Eq. (83) already stops at $`\overline{k}=\overline{n}n`$. Furthermore, as mentioned above, parameters with coinciding index pairs can be assumed to be zero. This leads to Eq. (84). Eq. (85) only reminds us that all subrelevant parts are removed if $`\beta `$ is used instead of $`\alpha `$.
The closedness conditions take the form
$$\beta _{l+1,\mathrm{}}^{\lambda ,\mathrm{}}(j,\sigma )\beta _{l,\mathrm{}}^{\lambda ,\mathrm{}}(j,\sigma )=\beta _{j+1,\mathrm{}}^{\sigma ,\mathrm{}}(l,\lambda )\beta _{j,\mathrm{}}^{\sigma ,\mathrm{}}(l,\lambda ).$$
(86)
The $`n`$-CH boundary phase is given by
$$\mathrm{\Phi }_\sigma ^{(n)}=\mathrm{\Phi }_\sigma ^{(n1)}+\underset{\lambda ,\mathrm{}}{}\varphi _\sigma ^{(n)}(\lambda ,\mathrm{})(N_\lambda \delta _{\lambda ,\sigma })\mathrm{},$$
(87)
where the phases $`\varphi ^{(n)}`$ are
$`\varphi _\sigma ^{(n)}(\{𝒮\})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\underset{\{\}}{\overset{\{𝒮\}}{\stackrel{(n)}{\beta }}}(j,\sigma );`$ (89)
$`\sigma \{𝒮\},`$
$`\varphi _\sigma ^{(n)}(\sigma ,\mathrm{},\sigma ,\{𝒮^{}\})`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{\genfrac{}{}{0pt}{}{j=1}{jl_i,l_i1}}{i=1..k}}{\overset{L}{}}}\underset{l_1,\mathrm{},l_k,\{^{}\}}{\overset{\sigma ,\mathrm{},\sigma ,\{𝒮^{}\}}{\stackrel{(n)}{\beta }}}(j,\sigma )+`$ (90)
$`+{\displaystyle \underset{i=1}{\overset{k}{}}}`$ $`[`$ $`\underset{\mathrm{},l_i1,\mathrm{},\{^{}\}}{\overset{\sigma ,\mathrm{},\sigma ,\{𝒮^{}\}}{\stackrel{(n)}{\beta }}}(l_i,\sigma )+`$ (91)
$`+`$ $`\underset{\mathrm{},l_i+1,\mathrm{},\{^{}\}}{\overset{\sigma ,\mathrm{},\sigma ,\{𝒮^{}\}}{\stackrel{(n)}{\beta }}}(l_i1,\sigma )];`$ (93)
$`\sigma \{𝒮^{}\}.`$
The pair of index sets $`\{𝒮\},\left|\{𝒮\}\right|=n`$ and $`\{\},\left|\{\}\right|=n`$, respectively $`\{𝒮^{}\},\left|\{𝒮^{}\}\right|=nk`$ and $`\{^{}\},\left|\{^{}\}\right|=nk`$ must not have coinciding index pairs and in Eq. (93) no two $`l_i`$ must be identical or neighbored. Thus, Eq. (93) holds for $`L2k`$. This validity range can be maximally enlarged using the closedness conditions (see Eq. (86))
$`\varphi _\sigma ^{(n)}(\sigma ,\mathrm{},\sigma ,\{𝒮^{}\})`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j=1}{j[l,l+k]}}{\overset{L}{}}}\underset{l+1,\mathrm{},l+k,\{^{}\}}{\overset{\sigma ,\mathrm{},\sigma ,\{𝒮^{}\}}{\stackrel{(n)}{\beta }}}(j,\sigma )+`$ (94)
$`+`$ $`{\displaystyle \underset{i=0}{\overset{k}{}}}\underset{\{\}_{i,k},\{^{}\}}{\overset{\sigma ,\mathrm{},\sigma ,\{𝒮^{}\}}{\stackrel{(n)}{\beta }}}(l+i,\sigma ),`$ (95)
$`\{\}_{i,k}=[l,l+k+1]`$ $``$ $`\{l+i,l+i+1\},`$ (96)
$`\sigma `$ $``$ $`\{𝒮^{}\}.`$ (97)
This formula holds for $`Lk+1`$, which is a limit set by physics – analogous to the $`2`$-CH case. The result is $`l`$-independent. We point out that $`\mathrm{exp}[\mathrm{i}F(\{n\})]`$ with $`F(\{n\})`$ being an arbitrary functional of number operators is not the most general unitary operator in Fock space. The most general is $`\mathrm{exp}[\mathrm{i}G(\{c^{},c\})]`$ where $`G(\{c^{},c\})`$ constitutes a Hermitean functional of the complete set of creation/annihilation operators. In the Appendix C we will show that the class of unitary operators discussed so far is large enough as far as removals of phases in a nearest neighbor hopping term are concerned.
## V Summary and conclusions
In summary, we found a complete characterization of Coordinate Bethe Ansatz (CBA) solvable Hubbard-type Hamiltonians with unitary correlated hopping for fermions. A necessary and sufficient criterion for such a Hamiltonian being solvable by CBA was formulated (see Eqs. (11), (12), and (64)).
In contrast to what is suggested in Ref. , we find that these models are not CBA solvable in the ordinary plane waves basis. Indeed, in such a basis the scattering matrix is configuration dependent (see Appendix A) thus describing diffractive scattering. The particles interact non–trivially even for vanishing Coulomb interaction. Such a situation resembles the case of the impurity problem in the sense that also there, the free picture already contains some residual interaction due to the impurity. Solvability of the models is recovered if the correlations from the hopping terms can be gauged away. This is equivalent to equipping the plane-wave basis with additional density dependent phases. Only in this modified basis can the correlated hopping be absorbed in a boundary term. For the models considered in Refs. no such basis exists and hence, they are not solvable by Bethe ansatz as they stand. They can however be repaired by modifying slightly their hopping term, as done at the end of chapter III.
The boundary twists for solvable models with periodic boundary conditions are given explicitly in this work. The corresponding Bethe equations are known from Ref. adapting the boundary phases for a spin degree of freedom only
$`\mathrm{e}^{\mathrm{i}p_jL}=\mathrm{e}^{\mathrm{i}\mathrm{\Phi }_{}}{\displaystyle \underset{a=1}{\overset{N_{}}{}}}{\displaystyle \frac{i(\mathrm{sin}p_j\zeta _a)\frac{V}{4t}}{i(\mathrm{sin}p_j\zeta _a)+\frac{V}{4t}}},`$ (98)
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{b=1}{ba}}{\overset{N_{}}{}}}{\displaystyle \frac{i(\zeta _a\zeta _b)+\frac{V}{2t}}{i(\zeta _a\zeta _b)\frac{V}{2t}}}=\mathrm{e}^{\mathrm{i}(\mathrm{\Phi }_{}\mathrm{\Phi }_{})}{\displaystyle \underset{l=1}{\overset{N}{}}}{\displaystyle \frac{i(\mathrm{sin}p_l\zeta _a)\frac{V}{4t}}{i(\mathrm{sin}p_l\zeta _a)+\frac{V}{4t}}}.`$ (99)
Therein, $`p_j`$ are the quasi-momenta in the plane waves used in the Bethe ansatz (see Eq. (A1)), $`\zeta _a`$ are spin rapidities, $`t`$ and $`V`$ are the Hubbard model parameters and $`\mathrm{\Phi }_\sigma `$ are the boundary phases, which have been determined in this paper. The connection to Ref. shows that in the solvable cases, the correlation in the hopping term of the Hamiltonian (see, for instance, (3)) is equivalent to applying a magnetic flux to the system. However, such a flux is generated by the particles themselves (in particular, it is not an external magnetic flux). Ground state properties can be deduced from those calculated in Refs. . Even for absent Coulomb interaction, the many particle energy is not a sum of single particle energies. The effect in the energy density is of first order in $`1/L`$ and in $`|\mathrm{\Phi }_{}\mathrm{\Phi }_{}|`$ in the thermodynamic limit. So one can argue that correlated hopping accounts for a non trivial interaction between the particles even for vanishing Coulomb interaction.
The results obtained here can be applied to models for particles with deformed exchange statistics(DES). This is done via a mapping from DES to CH models, whereas special DES models have been discussed recently using direct Bethe ansatz . The details will appear in a forthcoming paper.
###### Acknowledgements.
Motivating and fruitful discussions with D. Braak, M. Dzierzawa, M. Rasetti, P. Schwab, and B.S. Shastry are gratefully acknowledged besides the support through the Graduiertenkolleg “Nonlinear Problems in Analysis, Geometry, and Physics” (GRK 283), financed by the German Science Foundation (DFG) and the State of Bavaria; this work was also partly supported by the SFB 484.
## A Configuration dependent $`S`$ matrix from standard CBA in presence of correlated hopping
In this appendix we explain why correlated hopping destroys solvability by direct CBA for Hubbard- and XXZ-type models.
The CBA is an ansatz for the wave function in a so-called fundamental region , in which the interaction term doesn’t contribute. This fundamental region has to exist at first. For the Hubbard model, this is guaranteed by particle-hole symmetries. This gives the energy in terms of the distinct quasi-momenta. The wave function has to be defined uniquely on the intersection lines of the fundamental regions, where for the Hubbard model also the interaction enters. Both demands yield the $`S`$ matrix, which represents the effect of the interaction, a scattering of two particles. For Bethe ansatz solvability, the interaction must not have an effect beyond permuting the quasi-momenta of the particles. This demands that the $`S`$ matrix fulfills the Yang-Baxter equation. Now the boundary conditions, if compatible with the $`S`$ matrix, lead to additional conditions, fixing the eigenfunctions constructed by Bethe ansatz up to normalization. As a consequence, any additional condition destroys Bethe ansatz solvability. Every correlation in the hopping further restricts the parameters in the Bethe ansatz.
Let us assume a given CH Hubbard-type model can be transformed iso-spectrally to another CH Hubbard-type model, for which fundamental regions exist. This should mean that neither interaction nor correlated hopping contributes. All correlations already contribute, when no interaction is yet to be included. The Bethe Ansatz will have the form
$$\psi (x_1,\mathrm{},x_N)=:\underset{\pi S_N}{}A(\pi |\pi ^{})\mathrm{e}^{\mathrm{i}\underset{k=1}{\overset{N}{}}x_{\pi ^{}(k)}p_{\pi (k)}}$$
(A1)
where the permutation $`\pi ^{}`$ is chosen such that an appropriate order in $`(x_1,\mathrm{},x_N)`$ is achieved. Here, the chosen order which defines the fundamental regions, will be: $`x_1\mathrm{}x_N`$.
### 1 “Pinned” correlations
At first assume a correlation which appears if a particle sits at a special site. Let particle number $`j`$ sit on this site. The effect of this is a shift in all momenta by the correlation strength $`\phi `$, except that momentum of the designated particle. This is seen by projecting the Schrödinger equation on the specified configuration. The particle causing the correlation with the others feels no correlation from the hopping term applied to it, since it will be transported to that site by the hopping. Thus the resulting term for the energy is
$$E=2t\underset{lj}{\overset{N}{}}\mathrm{cos}(p_{\pi (l)}+\phi )2t\mathrm{cos}(p_{\pi (j)}).$$
(A2)
In this equation, $`\pi `$ is the momentum permutation from the Bethe ansatz. This energy is neither independent of the permutation as it ought to be, nor does it coincide with the original energy formula $`E=2t_j\mathrm{cos}(p_j)`$. Note that $`\phi `$ could even depend on the spin orientation of the considered particles.
### 2 Relative correlations
Assuming the hopping term to be
$$c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}(\mathrm{i}\phi n_{j+\mathrm{\Delta },\mu })+\mathrm{h}.\mathrm{c}.,$$
(A3)
where $`\sigma `$ and $`\mu `$ are spin indices. Further assume one single particle at coordinate $`j`$ to be affected by this correlation. The corresponding condition from the Schrödinger equation is
$$\begin{array}{c}\psi (j1,\sigma ;j^{},\mu )\left(\mathrm{e}^{\mathrm{i}\phi n_{j1+\mathrm{\Delta },\mu }}1\right)\hfill \\ +\psi (j+1,\sigma ;j^{},\mu )\left(\mathrm{e}^{\mathrm{i}\phi n_{j+\mathrm{\Delta },\mu }}1\right)=0,\hfill \end{array}$$
(A4)
where we omitted the spin index from the argument of the wavefunction. Two cases can appear independently from each other:
* a particle with spin orientation $`\mu `$ sits at site $`j^{}=j1+\mathrm{\Delta }`$
* a particle with spin orientation $`\mu `$ sits at site $`j^{}=j+\mathrm{\Delta }`$
leading to $`\psi (j1,\sigma ;j1+\mathrm{\Delta },\mu )=0`$ in the first case and $`\psi (j+1,\sigma ;j+\mathrm{\Delta },\mu )=0`$ in the second case. These constraints appear in addition to the usual “continuity” condition arising from $`\psi (j,\sigma ;j,\mu )=\psi (j,\sigma ;j,\mu )`$. If $`\sigma =\mu `$, at the most two of these three conditions can coincide if $`\mathrm{\Delta }=0`$ or $`\mathrm{\Delta }=1`$, which however constitute a subrelevant and irrelevant correlation, respectively. For $`\mu =\sigma `$, one of the correlation terms then coincides with the on-site Coulomb interaction. The other, however, is still remaining. The only way out is to modify the Bethe ansatz slightly . In Ref. , the fundamental regions are without intersection for different kinds of particles (different spin orientations). Consequently, for different kinds of particles the “continuity” condition is absent and hence substituted by one equation coming e.g. from correlated hopping. The second equation has then to coincide with the interaction. This indicates that the hopping term
$$c_{j+1,\sigma }^{}c_{j,\sigma }\mathrm{exp}(\mathrm{i}\phi n_{j+\mathrm{\Delta },\sigma })+\mathrm{h}.\mathrm{c}.;\mathrm{\Delta }\{0,1\}$$
(A5)
is the only possible relevant correlated hopping which is tolerable by direct Bethe ansatz for Hubbard-type Hamiltonians (for XXZ-type models no relevant CH is treatable by direct Bethe ansatz). But then, the obtained $`S`$ matrix still had to fulfill the Yang-Baxter equation. In Ref. , this kind of hopping term is studied. There, $`\phi `$ is purely imaginary and $`\mathrm{\Delta }(\sigma ):=(\sigma +1)/2\{0,1\}`$ was chosen. This Hamiltonian is shown to be solvable by a modified Bethe ansatz in the absence of Coulomb interaction, $`V=0`$. For $`V0`$, the resulting $`S`$ matrix no longer fulfills the Yang-Baxter equation.
To summarize, except for a very special type of relative correlation, namely with particles at distance zero or one towards raising site number, one obtains independent additional sets of conditions which the wave function had to fulfill. This obstructs the direct Bethe ansatz, since the $`S`$ matrix becomes configuration dependent. It then no longer factorizes into two-particle $`S`$ matrices. Regarding the exceptional CH mentioned above, a configuration independent $`S`$ matrix is obtained by applying a variant of Bethe ansatz . But for finite Coulomb-interaction strength $`V`$, this $`S`$ matrix does not fulfill the Yang-Baxter equation. As a consequence, no CH XXZ-type model and no interacting CH Hubbard-type model is tractable by direct Bethe ansatz for finite $`V`$.
## B A note on subrelevant parameters
In the following, we will discuss the contribution of the subrelevant parameters. Recalling the definitions, irrelevant parameters like $`A_{j1,j}(\sigma )`$ do not at all affect the physics of the model, whereas subrelevant parameters like $`A_{j,j}(\sigma )`$ contribute to the uncorrelated part of the hopping. For this reason no recursive relations come from them. It is worthwhile noting that Eq. (12) is trivially fulfilled for $`m=j`$. For this case, the exponent produced after transformation by $`U`$ is $`in_m(2\xi _{j+1,j}^{\sigma ,\sigma }+A_{j,j}(\sigma ))`$. This term hence appears as a phase additional to $`\gamma _j(\sigma )`$. The parameters $`\xi _{j+1,j}^{\sigma ,\sigma }`$ are (up to an additive constant) given by jumping from $`0`$ to $`2`$ in the exceptional loop in figure 1
$$2(\xi _{j,j1}^{\sigma ,\sigma }\xi _{j+1,j}^{\sigma ,\sigma })=A_{j,j1}(\sigma )+A_{j1,j+1}(\sigma ).$$
(B1)
We will now consider the special case in which $`2\xi _{j+1,j}^{\sigma ,\sigma }+A_{j,j}(\sigma )=0`$ holds. This is exactly what the relation (7) would result in for $`m=j`$. It implies that the subrelevant parameters do not create additional phases in the uncorrelated part of the hopping. Together with (B1) it gives
$$A_{j,j}(\sigma )A_{j,j1}(\sigma )A_{j1,j+1}(\sigma )A_{j1,j1}(\sigma )=0.$$
(B2)
The condition for the relations (7) for $`m=j`$ and $`m=j+1`$ being consistent is $`A_{j1,j}=A_{j1,j1}`$, which can always be fulfilled by properly choosing the irrelevant parameter $`A_{j1,j}`$. Inserting this into Eq. (B2), one finally also exactly obtains relation (7) for $`m=j+1`$. We can therefore conclude that if Eq. (B2) holds, the contributions from subrelevant parameters cancel out. It bridges the void in the recursive relations (7).
Since all the $`\alpha `$’s are relevant, no voids occur in their recursive relations, which is equivalent to all elementary loops being viable.
## C Unitary transformations in the fermion $`u(r)`$ algebra
In this appendix, the transformed number operators $`c_i^{}c_i`$ and hopping operators $`c_i^{}c_j`$ will be studied for the most simple unitary transformation such as
$$U=\mathrm{exp}\left[\mathrm{i}\alpha (c_k^{}c_l+c_l^{}c_k)\right];kl.$$
(C1)
The $`r^2`$ operators $`\{c_k^{}c_l,\mathrm{\hspace{0.33em}1}i,jr\}`$ span the Lie algebra $`u(r)`$:
$$[c_i^{}c_j,c_k^{}c_l]=\delta _{j,k}c_i^{}c_j\delta _{i,l}c_k^{}c_j$$
(C2)
where the Cartan basis is generated by $`H_i=c_i^{}c_i`$, $`\left(i=1\mathrm{}r\right)`$:
$$[H_i,c_j^{}c_k]=(\delta _{i,j}\delta _{i,k})c_j^{}c_k$$
(C3)
The $`u(r)`$ algebraic structure allows us to write
$$\begin{array}{c}\mathrm{e}^{\mathrm{i}\alpha \mathrm{\Phi }_{k,l}^+}n_me^{i\alpha \mathrm{\Phi }_{k,l}}=\hfill \\ \\ n_m+i\mathrm{\Phi }_{k,l}^{}(\delta _{l,m}\delta _{k,m})\hfill \\ \mathrm{sin}^2\alpha (\delta _{l,m}\delta _{k,m})(n_ln_k)+\hfill \\ +i(\mathrm{sin}\alpha \mathrm{cos}\alpha \alpha )(\delta _{l,m}\delta _{k,m})\mathrm{\Phi }_{k,l}^{}\hfill \\ \\ e^{i\alpha \mathrm{\Phi }_{k,l}^+}\mathrm{\Phi }_{m,n}^+e^{i\alpha \mathrm{\Phi }_{k,l}}=\hfill \\ \\ \mathrm{\Phi }_{m,n}^++i\mathrm{sin}\alpha [\mathrm{\Phi }_{k,n}^{}\delta _{l,m}+\mathrm{\Phi }_{l,n}^{}\delta _{k,m}+\hfill \\ +\mathrm{\Phi }_{k,m}^{}\delta _{l,n}+\mathrm{\Phi }_{l,m}^{}\delta _{k,n}]+\hfill \\ \\ +4\mathrm{sin}\alpha [\mathrm{\Phi }_{k,l}^+(\delta _{k,m}\delta _{l,n}+\delta _{l,m}\delta _{k,n})+\hfill \\ +\mathrm{\Phi }_{m,n}^+(\delta _{k,m}+\delta _{l,n}+\delta _{l,m}+\delta _{k,n})]\hfill \end{array}$$
(C4)
where we applied the following equivalence for adjoints
$$[Ad(\mathrm{exp}(A))]B:=e^ABe^A=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}[A,B]_n=:\mathrm{exp}[ad(A)]B,$$
$$[A,B]_0:=B,[A,B]_{n+1}:=[A,[A,B]_n],$$
and we have defined $`\mathrm{\Phi }_{k,l}^+:=c_k^{}c_l+c_l^{}c_k`$, and $`\mathrm{\Phi }_{k,l}^{}:=c_k^{}c_lc_l^{}c_k`$ (the $`\mathrm{\Phi }^+`$ are Hermitean, whereas the $`\mathrm{\Phi }^{}`$ are anti Hermitean).
¿From this it is seen that the transformation of the Hubbard Hamiltonian creates arbitrary-range hopping from and to the sites $`k`$ and $`l`$, as well as pair-hopping created from the interaction term. This can be understood from interpreting $`\mathrm{\Phi }_{k,l}^+`$ as a Hamiltonian itself and $`\alpha `$ as the time. Then, Eq. (C4) gives the number and hopping operators in the Heisenberg picture. ¿From this interpretation it seems reasonable that no linear combination of $`\mathrm{\Phi }_{k,l}^+`$ with $`kl`$ will ever be able to just remove phases in a nearest-neighbor hopping term. However, this is not absolutely true. To point out the exceptional cases, the investigation has to be completed. Since the identity $`[Ad(\mathrm{exp}(A))]B=\mathrm{exp}[ad(A)]B`$ does not considerably simplify calculating the action of a more general unitary transformation, another approach will be taken. But at first, the problem will be reduced as far as possible.
In a general product of creation and annihilation operators, one can at first collect operators occurring in pairs using the exchange algebra. The result is a multilinear form of number operators besides a multi-linear form of creation and annihilation operators, so that no two operators have coinciding indices. So the most general exponent appearing in $`U`$ can be written as follows
$$\begin{array}{c}\xi _{\{p\},\{q\}}_ic_{p_i}^{}c_{q_i}+\xi _{k;\{p\},\{q\}}n_k_ic_{p_i}^{}c_{q_i}+\hfill \\ \\ +\xi _{k,l;\{p\},\{q\}}n_kn_l_ic_{p_i}^{}c_{q_i}+\mathrm{}.\hfill \end{array}$$
(C5)
Here, the different $`\xi `$ are symmetric in the indices before the semicolon and Hermitean in the index sets behind it. They vanish if any two indices coincide. Assuming a pure $`m`$-linear form of the creation/annihilation operators, then a transformed bilinear object contains $`m`$-linear and even higher terms. This leads to multi-particle hopping and interaction terms including more than two number operators . Since the aim is to stay in the class of Hubbard- or XXZ-type Hamiltonians, we consequently can limit ourselves to general multilinear forms of number operators only, as already studied above, or bilinear forms of creation/annihilation operators only. We will discuss the latter in the following.
$$U:=\mathrm{exp}\left[\mathrm{i}\xi _{k,l}c_k^{}c_l\right];\xi _{k,l}=\xi _{l,k}^{},\xi _{k,k}=0.$$
(C6)
We point out that $`Uc_m^{}c_nU^1`$ is bilinear, since $`[c_k^{}c_l,c_m^{}c_n]`$ is bilinear (see C2). So, one can determine the result by projecting on the desired initial and final states.
$`\begin{array}{c}0\left|c_{m_f}\left[\xi _{k,l}c_k^{}c_l\right]^nc_{m_i}^{}\right|0=\hfill \\ =0\left|c_{m_f}\left[\xi _{k,l}c_k^{}c_l\right]c_{m_{n1}}^{}\right|00|c_{m_{n1}}\mathrm{}\hfill \\ \mathrm{}c_{m_1}^{}|00\left|c_{m_1}\left[\xi _{k,l}c_k^{}c_l\right]c_{m_i}^{}\right|0=\hfill \\ =\xi _{m_f,m_{n1}}\xi _{m_{n1},m_{n2}}\mathrm{}\xi _{m_1,m_i}=\left(\xi ^n\right)_{m_f,m_i},\hfill \end{array}`$
all written in sum convention. One directly obtains from this
$`0\left|c_{m_f}Uc_{m_i}^{}\right|0`$ $`=`$ $`\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_f,m_i}`$ (C7)
$`0\left|c_{m_f}Uc_k^{}c_lU^{}c_{m_i}^{}\right|0`$ $`=`$ $`\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_f,k}\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_i,l}^{},`$ (C8)
where hermitecity of the $`\xi `$ was used. The indices $`k,l`$ are not summed over. It can be shown that the transformed interaction term can never include a single-particle hopping term . This already proves that the reverse direction is also impossible. Thus, number operators and nearest-neighbor hopping have to remain “type-invariant” under the transformation, since the type of Hamiltonian should be preserved. This results in restrictions on the matrix $`\xi `$
$`\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_f,k}\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_i,k}^{}`$ $`=`$ $`0m_fm_i`$ (C9)
$`\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_f,j+1}\left(\mathrm{exp}\mathrm{i}\xi \right)_{m_i,j}^{}`$ $`=`$ $`0\text{f}or|m_fm_i|>1.`$ (C10)
The first conditions emerge from transforming the interaction term, whereas the second one comes from the hopping term. Using both, we can deduce the structure of $`\mathrm{exp}\mathrm{i}\xi `$.
$`\begin{array}{c}(\text{C9})k_1m(k)\left(\mathrm{exp}\mathrm{i}\xi \right)_{m(k),k}0\hfill \\ (\text{C10})m(k)=k+const.\hfill \\ \\ \left(\mathrm{exp}\mathrm{i}\xi \right)_{k,l}(\mathrm{\Delta })=:\delta _{l,k+\mathrm{\Delta }}r_k\mathrm{e}^{\mathrm{i}\varphi _k};_kr_l\mathrm{e}^{\mathrm{i}\varphi _l}=1.\hfill \end{array}`$
Unitarity of the matrix implies $`r_l^2=1`$ for all $`l`$. Hence it can be assumed that $`r_l=1`$ for all $`l`$. So, one finally concludes that $`U(\mathrm{\Delta })`$ transforms $`c_{j+1}^{}c_j`$ into $`c_{j+\mathrm{\Delta }+1}^{}c_{j+\mathrm{\Delta }}\mathrm{e}^{\mathrm{i}(\varphi _{j+1}\varphi _j)}`$. Number operators remain unchanged. These phases can be gauged away leaving no boundary phase. It is worth noting that pure $`d`$-range hopping on an $`L`$-site chain can also be obtained iff $`d`$ and $`L`$ are relatively prime. With the analysis used here it can finally be shown that unitary transformations with multinomials of odd degree always produce particle-number violating terms. Hence, they also make us leave the class of models we consider. But it is clear from this, that a huge class of models, which is far from being Hubbard-type can be constructed by unitarily transforming the Hubbard model. They all are solvable and have the same spectrum as the Hubbard model.
## D Proofs
### 1 Proof that transformations such as (C5) exceed the considered class of models
Consider unitary transformations $`U`$ given by
$`\begin{array}{c}U=\mathrm{exp}\mathrm{i}(\mathrm{\Omega }+\mathrm{\Xi })\hfill \\ \\ \mathrm{\Xi }=\xi _{\{p\},\{q\}}_ic_{p_i}^{}c_{q_i}+\xi _{k;\{p\},\{q\}}n_k_ic_{p_i}^{}c_{q_i}+\hfill \\ \\ +\xi _{k,l;\{p\},\{q\}}n_kn_l_ic_{p_i}^{}c_{q_i}+\mathrm{},\hfill \end{array}`$
where $`\mathrm{\Omega }`$ is a multilinear form of number operators only or bilinear in creation/annihilation operators. Thus, the Hermitean $`\mathrm{\Xi }`$ only contains terms higher than bilinear.
Now let us assume
$$\mathrm{exp}(\mathrm{i}\mathrm{\Xi })n_kn_{k+1}\mathrm{exp}(\mathrm{i}\mathrm{\Xi })=\underset{l}{}a_l^kn_ln_{l+1}$$
(D1)
Where $`a_l^k`$ are constants for fixed $`k`$ and $`l`$. This is impossible, because this expression contains terms at least hexa-linear in creation/annihilation operators since $`\mathrm{\Xi }`$ is at least quadri-linear. This still holds including $`\mathrm{\Omega }`$.
Next assume
$$\mathrm{exp}(\mathrm{i}\mathrm{\Xi })n_kn_{k+1}\mathrm{exp}(\mathrm{i}\mathrm{\Xi })\mathrm{exp}(\mathrm{i}f[\{n\}])c_{l+1}^{}c_l$$
(D2)
and consider $`0\left|c_{m_f}\mathrm{exp}(\mathrm{i}\mathrm{\Xi })n_kn_{k+1}\mathrm{exp}(\mathrm{i}\mathrm{\Xi })c_{m_i}^{}\right|0`$. This had to result in $`\delta _{m_f,l+1}\delta _{m_i,l}\mathrm{exp}(\mathrm{i}g[\{n\}])`$. Here, $`f[\{n\}]`$ is an arbitrary functional of the number operators. In general, the functional $`g[\{n\}]`$ differed from $`f[\{n\}]`$. However, $`0\left|c_{m_f}\mathrm{\Xi }c_{m_i}^{}\right|0=0`$, since at least two indices of each coefficient $`\xi `$ in $`\mathrm{\Xi }`$ had to coincide in order to give a non-zero contribution. But then the $`\xi `$ vanish themselves. With the same argument, $`0\left|c_{m_f}\mathrm{\Xi }^nc_{m_i}^{}\right|0=0`$ for all integer $`n`$ (including zero, since $`m_im_f`$ is assumed). Thus, the assumption cannot hold. This completes the proof, since including $`\mathrm{\Omega }`$ in the transformation $`U`$, the only nonvanishing contributions from $`0\left|c_{m_f}\mathrm{exp}\mathrm{i}(\mathrm{\Omega }+\mathrm{\Xi })n_kn_{k+1}\mathrm{exp}\mathrm{i}(\mathrm{\Omega }+\mathrm{\Xi })c_{m_i}^{}\right|0`$ come from terms not including $`\mathrm{\Xi }`$ at all. This means, that the only contributing terms come from $`\mathrm{\Omega }`$ alone.
### 2 Proof that the transform (C6) cannot interchange hopping and interaction
Let us assume, $`Un_jn_{j+1}U^1`$ (in case of XXZ-type models) contained a term $`c_{k+1}^{}c_k`$. Applying Eq.(C8), this means
$$\begin{array}{c}(\mathrm{exp}\mathrm{i}\xi )_{m_f,j}(\mathrm{exp}\mathrm{i}\xi )_{m_i,j}^{}\times \hfill \\ \times (\mathrm{exp}\mathrm{i}\xi )_{\stackrel{~}{m}_f,j+1}(\mathrm{exp}\mathrm{i}\xi )_{\stackrel{~}{m}_i,j+1}^{}c_{m_f}^{}c_{m_i}c_{\stackrel{~}{m}_f}^{}c_{\stackrel{~}{m}_i}\hfill \end{array}$$
(D3)
contained a term $`c_{k+1}^{}c_k`$. Here, $`j`$ is fixed and all the $`m`$’s are summed over. Eq.(D3) implies
1. $`\stackrel{~}{m}_i=k`$, $`m_f=k+1`$ and $`m_i=\stackrel{~}{m}_f=:m`$, or
2. $`m_i=k`$, $`m_f=k+1`$ and $`\stackrel{~}{m}_i=\stackrel{~}{m}_f=:m`$.
Considering 1.: Eq.(D3) in this case becomes
$$\begin{array}{c}(\mathrm{exp}\mathrm{i}\xi )_{k+1,j}(\mathrm{exp}\mathrm{i}\xi )_{m,j}^{}\times \hfill \\ \times (\mathrm{exp}\mathrm{i}\xi )_{m,j+1}(\mathrm{exp}\mathrm{i}\xi )_{k,j+1}^{}c_{k+1}^{}c_mc_m^{}c_k,\hfill \end{array}$$
(D4)
which is proportional to
$`\begin{array}{c}(\mathrm{exp}\mathrm{i}\xi )_{j,m}(\mathrm{exp}\mathrm{i}\xi )_{m,j+1}(\mathrm{exp}\mathrm{i}\xi )_{k+1,j}(\mathrm{exp}\mathrm{i}\xi )_{k,j+1}^{}\times \hfill \\ \times (1n_m)c_{k+1}^{}c_k.\hfill \end{array}`$
Since $`m`$ is summed over, the $`1`$ in the braces vanishes. This is because $`\mathrm{exp}(\mathrm{i}\xi )\mathrm{exp}\mathrm{i}\xi =1\mathrm{l}`$. Additionally, no two-particle hopping must occur. This implies that for arbitrary $`m\overline{m}`$
$$\begin{array}{c}(\mathrm{exp}\mathrm{i}\xi )_{m,j+1}(\mathrm{exp}\mathrm{i}\xi )_{\overline{m},j}^{}\times \hfill \\ \times \underset{E_1}{\underset{}{(\mathrm{exp}\mathrm{i}\xi )_{k+1,j}(\mathrm{exp}\mathrm{i}\xi )_{k,j+1}^{}}}\stackrel{!}{=}0.\hfill \end{array}$$
(D5)
Vanishing $`E_1`$ means no occurrence of nearest neighbor hopping. Hence, $`E_1`$ is assumed to be nonzero. Since (D4)$`0`$ for every fixed $`j`$ one can find an $`m`$ so that $`(\mathrm{exp}\mathrm{i}\xi )_{m,j}0`$ and $`(\mathrm{exp}\mathrm{i}\xi )_{m,j+1}0`$. Now assume the existence of a further $`\overline{m}m`$ so that even $`(\mathrm{exp}\mathrm{i}\xi )_{\overline{m},j}0`$. This implied $`(\mathrm{exp}\mathrm{i}\xi )_{m,j+1}(\mathrm{exp}\mathrm{i}\xi )_{\overline{m},j}0`$ in contradiction to Eq.(D5). Thus, there exists only a single $`m`$ for which $`(\mathrm{exp}\mathrm{i}\xi )_{m,j+1}`$ and $`(\mathrm{exp}\mathrm{i}\xi )_{m,j}`$ don’t vanish. But this would mean $`det\mathrm{exp}\mathrm{i}\xi =0`$, which would contradict the assumptions.
Considering 2.: Eq.(D3) now becomes
$`\begin{array}{c}(\mathrm{exp}\mathrm{i}\xi )_{k+1,j+1}(\mathrm{exp}\mathrm{i}\xi )_{k,j+1}^{}(\mathrm{exp}\mathrm{i}\xi )_{j,m}(\mathrm{exp}\mathrm{i}\xi )_{m,j}^{}\times \hfill \\ \times c_{k+1}^{}c_kc_m^{}c_m.\hfill \end{array}`$
As above, no two-particle hopping must occur. This implies that for arbitrary $`m\overline{m}`$
$$\begin{array}{c}\underset{E_2}{\underset{}{(\mathrm{exp}\mathrm{i}\xi )_{k+1,j+1}(\mathrm{exp}\mathrm{i}\xi )_{k,j+1}^{}}}\times \hfill \\ \times (\mathrm{exp}\mathrm{i}\xi )_{\overline{m},j}(\mathrm{exp}\mathrm{i}\xi )_{m,j}^{}\stackrel{!}{=}0.\hfill \end{array}$$
(D6)
Vanishing $`E_2`$ means no occurrence of nearest neighbor hopping. Hence, $`E_2`$ is assumed to be nonzero. Since $`det\mathrm{exp}\mathrm{i}\xi 0`$, one finds for every fixed $`j`$ an $`m`$ so that $`(\mathrm{exp}\mathrm{i}\xi )_{m,j}0`$. Now assume a further $`\overline{m}m`$ to exist so that $`(\mathrm{exp}\mathrm{i}\xi )_{\overline{m},j}0`$. This implies $`(\mathrm{exp}\mathrm{i}\xi )_{m,j}(\mathrm{exp}\mathrm{i}\xi )_{\overline{m},j}0`$ in contradiction to Eq.(D6). Thus, there exists only a single $`m`$ for which $`(\mathrm{exp}\mathrm{i}\xi )_{m,j}`$ is nonzero. This is in contradiction to $`E_2`$ being nonzero.
This completes the proof. |
warning/0003/quant-ph0003078.html | ar5iv | text | # Transfer of nonclassical features in quantum teleportation via a mixed quantum channel
## I Introduction
Quantum teleportation is one of the important manifestations of quantum mechanics. By quantum teleportation an unknown quantum state is destroyed at a sending station while its replica state appears at a remote receiving station via dual quantum and classical channels. The key to quantum teleportation is the entanglement of the quantum channel. Quantum teleportation has been studied for various systems including two-level systems , $`N`$-dimensional systems , and continuous variables . In particular, quantum teleportation of continuous variable states has been at a focus because of a high detection efficiency and handy manipulation of continuous variable states .
Quantum teleportation of a continuous-variable state was first suggested by Vaidman employing the Einstein-Podolsky-Rosen (EPR) state for the quantum channel in the framework of nonlocal measurements . Braunstein and Kimble made a use of quadrature-phase entanglement in a two-mode squeezed vacuum to teleport the quadrature-phase variables. With the high detection efficiency of the homodyne measurement and highly squeezed light, Ralph and Lam and Furusawa et al. realized quantum teleportation of continuous variable states by experiments. Ralph and Lam produced the required entangled state using two bright squeezed sources. A two-mode squeezed vacuum is entangled with respect not only to quadrature phases but also to photon-number difference and phase sum. Based on this number-phase entanglement, Milburn and Braunstein suggested another protocol to teleport a continuous variable state .
There are a few problems in the quantum teleportation of quadrature-phase variables using the two-mode squeezed vacuum. The perfect quantum teleportation is possible only with a maximally entangled state which means infinite squeezing in the squeezed state. The mean energy of a two-mode squeezed state increases exponentially as the squeezing increases so that the maximally entangled squeezed state is unphysical. As the quantum channel is exposed to the real world, it is influenced by the environment, which turns the pure squeezed state into a mixture and deteriorates the entanglement property. The environmental effect is unavoidable for any type of teleportation and there have been suggestions to purify mixed entangled state into a maximally entangled singlet state for a discrete two-level system . Duan et al. suggested a way to purify a Gaussian continuous variable state . However, their purification protocol may concentrate entanglement only to a finite dimensional Hilbert space. In fact, it is impossible to purify a two-mode squeezed state into a maximally entangled state as it is unphysical. Opartny et al. showed that the problem of not having the maximally entangled squeezed vacuum can be overcome by conditional measurements . Entanglement quantification and purification for continuous-variable states has been studied by Parker et al. . The imperfect detection efficiency and the imperfect realization of unitary transformation at the receiving station can also lower the efficiency of teleportation.
In this paper, we are interested in the efficiency of quantum teleportation in the real world. Nonclassical properties such as sub-Poissonicity and squeezing of the original state can be very useful for communication purposes. As the quantum channel is not maximally entangled, some or all of the nonclassical properties can be lost during the teleportation. Braunstein and Kimble found that when the quantum channel is not squeezed, i.e., when the channel is merely a two-mode vacuum, no quantum features can be observed in the teleported state . This is due to quantum tariffs of vacuum noise, which arises at the sending and receiving stations. The tariff was coined as quduty by Braunstein and Kimble. The pure two-mode squeezed state becomes mixed as the quantum channel is embedded in the environment. Quantum teleportation via the mixed channel can bear a different nature. For example, one may ask “Does the classical correlation play any role to transfer the nonclassical features?” It is not clear so far under which condition any nonclassical features implicit in an original unknown state cannot be transferred by teleportation via a mixed channel. We also consider the fidelity of teleportation to measure how close the teleported state is to the original state when the quantum channel is mixed. Popescu studied quantum teleportation of a discrete two-level system for a mixed quantum channel and found that even when the quantum channel is not maximally entangled, it has the fidelity better than any classical teleportation protocol . In this paper, we restrict ourselves to the situation that the decoherence effect is the same on each mode of the two-mode squeezed vacuum.
The continuous variable state can be easily analyzed using the quasiprobability functions . The description of a quantum-mechanical state in phase space is not unique due to the uncertainty principle; hence there are a family of quasiprobability functions of which the $`P`$, $`Q`$, and Wigner functions are widely used . In particular, it is well-known that the $`P`$ function can be used as a measure of the nonclassicality of a given field because the $`P`$ function is positive well-defined only for a classical state . The nonclassical depth is defined based on how much noise to put into the nonclassical state to have a positive well-defined $`P`$ function.
When teleportation is imperfect, a noise-added replica state is produced at the receiving station. By analyzing the added noise, we find the critical point for the quantum channel not to transfer any nonclassical features which may be imposed in an original unknown state. We examine the coincidence of the critical point with the moment when the quantum channel becomes separable. To do that we find the necessary and sufficient condition of separability for any two-mode Gaussian state , one of which is the mixed two-mode squeezed state. The fidelity, which is defined as the inner product of the original and teleported states, can be represented by the overlap of their Wigner functions . We show that the fidelity is a function of the added noise.
The added noise by teleportation is compared with that by direct transmission of the original state. It is found that the nonclassical nature of the original state can be more easily lost by teleportation than by direct transmission. This is because teleportation relies on the entanglement of the quantum channel, which is very fragile.
## II Quasiprobability functions
Before considering quantum teleportation, we briefly introduce the quasiprobability functions. The family of quasiprobability functions are obtained by the following convolution relation
$$R_\sigma (\alpha )=d^2\beta \left[\frac{2}{\pi (1\sigma )}\mathrm{exp}\left(\frac{2\left|\alpha \beta \right|^2}{1\sigma }\right)\right]P(\beta )$$
(1)
where the $`\sigma `$-parameterized $`R_\sigma (\alpha )`$ function becomes $`Q`$ function for $`\sigma =1`$, Wigner ($`W`$) function for $`\sigma =0`$, and $`P`$ function for $`\sigma =1`$. By the Fourier transform, we find the relation between their characteristic functions
$$C_\sigma ^R(\xi )=\mathrm{exp}\left[\frac{(1\sigma )|\xi |^2}{2}\right]C^P(\xi )$$
(2)
where $`C_\sigma ^R(\xi )`$ and $`C_\sigma ^P(\xi )`$ are the characteristic functions for the $`R`$ and $`P`$ functions, respectively. The family of two-mode quasiprobability functions can be analogously defined as
$`R_\sigma (\alpha ,\beta )=`$ $`{\displaystyle \frac{4}{\pi ^2(1\sigma )^2}}{\displaystyle d^2\varphi d^2\eta }`$ (4)
$`\times \mathrm{exp}\left({\displaystyle \frac{2\left|\alpha \varphi \right|^2}{1\sigma }}{\displaystyle \frac{2\left|\beta \eta \right|^2}{1\sigma }}\right)P(\varphi ,\eta ).`$
## III Teleportation for continuous variables in thermal environments
A continuous variable state $`\widehat{\rho }_o`$ can be teleported with use of a two-mode squeezed vacuum for a quantum channel . Two modes $`b`$ and $`c`$ of the squeezed vacuum are distributed separately to a sending and a receiving stations. The protocol comprises two operations at the sending station and one operation at the receiving station. At the sending station, the original unknown state of mode $`a`$ is mixed with a mode $`b`$ of the quantum channel by a 50/50 beam splitter. Two conjugate quadrature variables are measured respectively for the two output fields of the beam splitter. The measurement results are sent to the receiving station through the classical channel. The other mode $`c`$ of the squeezed vacuum is then displaced at the receiving station according to the measurement results. It is important to displace the photon of mode $`c`$ entangled with the photon measured at the sending station. Braunstein and Kimble considered the teleportation protocol for the pure state of the quantum channel . In this paper we investigate the teleportation via the mixed quantum channel to consider the influence of a thermal environment. We assume that the thermal environment gives the same effect on each mode of the quantum channel and the original state is prepared in a pure state.
The two-mode squeezed vacuum of the quantum channel is entangled and represented by the Wigner function
$`W_{qc}(\alpha _b,\alpha _c)=`$ $`{\displaystyle \frac{4}{\pi ^2}}\mathrm{exp}[2(|\alpha _b|^2+|\alpha _c|^2)\mathrm{cosh}2s_{qc}`$ (6)
$`+2(\alpha _b\alpha _c+\alpha _b^{}\alpha _c^{})\mathrm{sinh}2s_{qc}],`$
where $`s_{qc}`$ is the degree of squeezing and the complex quadrature phase variable $`\alpha _{b,c}=\alpha _{b,c}^r+i\alpha _{b,c}^i`$. When $`s_{qc}\mathrm{}`$ the state (6) manifests the maximum entanglement and becomes an EPR state. However, the mean photon number, which is $`2\mathrm{sinh}^2s_{sq}`$, becomes infinity in this limit.
Before the action of the beam splitter, the total state is a product of the original state and the state of the quantum channel, which is represented by the total Wigner function $`W_t(\alpha _a,\alpha _b,\alpha _c)=W_o(\alpha _a)W_{qc}(\alpha _b,\alpha _c)`$ where $`W_o(\alpha _a)`$ is the Wigner function of the original state $`\widehat{\rho }_o`$. The product state of the original field and quantum channel becomes entangled at the beam splitter. Considering the unitary action of the beam splitter, the quadrature variables $`\alpha _{d,e}`$ of the output fields are related to those of the input fields: $`\alpha _{d,e}=(\alpha _b\pm \alpha _a)/\sqrt{2}`$. The Wigner function $`W_t^B(\alpha _d,\alpha _e,\alpha _c)`$ for the total field after the beam splitter is
$$W_t^B(\alpha _d,\alpha _e,\alpha _c)=W_t(\frac{\alpha _e+\alpha _d}{\sqrt{2}},\frac{\alpha _e\alpha _d}{\sqrt{2}},\alpha _c)$$
(7)
which exhibits entanglement between the modes $`a`$ and $`b`$.
Setting homodyne detectors at the output ports of the beam splitter, the imaginary part of $`\alpha _d`$ and the real part of $`\alpha _e`$ are simultaneously measured by appropriately choosing the phases of reference fields for the homodyne detectors. Each measurement result is transmitted to the receiving station to displace the quadrature variables of the field mode $`c`$. We have to make it sure that the displacement operation is done on the photon of mode $`c`$ entangled with the photon measured at the sending station. After the displacement $`\mathrm{\Delta }(\alpha _d^i,\alpha _e^r)`$ the field of mode $`c`$ becomes to be represented by the Wigner function $`W_r(\alpha _c)`$;
$$W_r(\alpha _c)=d^2\alpha _dd^2\alpha _eW_t^B(\alpha _d,\alpha _e,\alpha _c\mathrm{\Delta }(\alpha _d^i,\alpha _e^r)).$$
(8)
Braunstein and Kimble found that the displacement of $`\mathrm{\Delta }(\alpha _d^i,\alpha _e^r)=\sqrt{2}(\alpha _e^ri\alpha _d^i)`$ maximizes the fidelity when the channel is a pure two-mode squeezed state. The probability $`P(\alpha _d^i,\alpha _e^r)`$ of measuring $`\alpha _d^i`$ and $`\alpha _e^r`$ at the sending station is the same as the marginal Wigner function
$`P(\alpha _d^i,\alpha _e^r)`$ $`=`$ $`{\displaystyle 𝑑\alpha _d^r𝑑\alpha _e^id^2\alpha _cW_t^B(\alpha _d,\alpha _e,\alpha _c)}.`$ (9)
### A two-mode squeezed vacuum in thermal environments
The quantum channel initially in the two-mode squeezed vacuum results in a mixed state by the interaction with the thermal environment. Assuming that two thermal modes are independently coupled with the quantum channel the dynamics of the squeezed field is described by a Fokker-Planck equation in the interaction picture
$`{\displaystyle \frac{W_{qc}(\alpha _b,\alpha _c;t)}{t}}=`$ $`{\displaystyle \frac{\gamma }{2}}{\displaystyle \underset{i=b,c}{}}[{\displaystyle \frac{}{\alpha _i}}\alpha _i+{\displaystyle \frac{}{\alpha _i^{}}}\alpha _i^{}`$ (11)
$`+(1+2\overline{n}){\displaystyle \frac{^2}{\alpha _i\alpha _i^{}}}]W_{qc}(\alpha _b,\alpha _c;t),`$
where $`\overline{n}`$ is the average photon number of the thermal environment. The two thermal modes are assumed to have the same average energy and coupled with the channel in the same strength. This assumption is reasonable as the two modes of the squeezed state are in the same frequency and the temperature of the environment is normally the same. By solving the Fokker-Planck equation (11), the time-dependent Wigner function is obtained as
$`W_{qc}(\alpha _b,\alpha _c;T)=𝒩\mathrm{exp}`$ $`[{\displaystyle \frac{2\mathrm{\Gamma }}{\mathrm{\Gamma }^2\mathrm{\Lambda }^2}}(|\alpha _b|^2+|\alpha _c|^2)`$ (13)
$`+{\displaystyle \frac{2\mathrm{\Lambda }}{\mathrm{\Gamma }^2\mathrm{\Lambda }^2}}(\alpha _b\alpha _c+\alpha _b^{}\alpha _c^{})]`$
where $`𝒩`$ is the normalization factor and two parameters, $`\mathrm{\Gamma }=T(1+2\overline{n})+(1T)\mathrm{cosh}2s_{qc}`$, $`\mathrm{\Lambda }=(1T)\mathrm{sinh}2s_{qc}`$. The renormalized time $`T(t)=1\mathrm{exp}(\gamma t)`$. The relative strength of $`\mathrm{\Lambda }`$ to $`\mathrm{\Gamma }`$ determines how much the mixed channel is entangled. When $`\mathrm{\Lambda }`$ is zero for $`T1`$, the channel loses any correlation so to have neither classical nor quantum correlation. At $`T=0`$ the mixed squeezed state (13) is simply the squeezed vacuum (6).
When the quantum channel is embedded in thermal environments, the teleported state is still represented by the Wigner function (8) with the quantum channel (13). However, a question remains in the unitary operation at the receiving station when the channel is a mixed state. According to the philosophy of the faithful teleportation, the displacement has to be determined to maximize the fidelity of teleportation. The fidelity $``$, which measures how close the teleported state is to the original state, is the projection of the original pure state $`|\mathrm{\Psi }_o`$ onto the teleported state of the density operator $`\widehat{\rho }_r`$: $`=\mathrm{\Psi }_o|\widehat{\rho }_r|\mathrm{\Psi }_o`$. The fidelity is also represented by the overlap between the Wigner functions for the original and teleported states ;
$$=\pi d^2\alpha W_o(\alpha )W_r(\alpha ).$$
(14)
For a maximally entangled quantum channel, the original pure state is reproduced at the receiving station and the fidelity is unity. For an impure or partially entangled channel, the unitary operation at the receiving station may depend on original states to maximize the fidelity, which has been shown for the teleportation of a two-level state . For the infinite dimensional Hilbert space, a formal study is very much complicated. However, we have found that even when the channel is mixed, the displacement of $`\mathrm{\Delta }(\alpha _d^i,\alpha _e^r)=\sqrt{2}(\alpha _e^ri\alpha _d^i)`$ maximizes the fidelity for a coherent projector $`|\mu \nu ^{}|`$, where $`|\mu `$ and $`|\nu ^{}`$ are coherent-state bases. An unknown state can be written as a weighted integral of the coherent projection operators
$$\widehat{\rho }_o=d^2\mu d^2\nu P_o(\mu ,\nu )|\mu \nu ^{}|$$
(15)
where $`P(\mu ,\nu )`$ is proportional to the positive-$`P`$ function . The unitary operation, which maximizes the fidelity, at the receiving station is thus independent of the original state.
### B separability of the quantum channel
A discrete bipartite system of modes $`b`$ and $`c`$ is separable when its density operator is represented by $`\widehat{\rho }=_rP_r\widehat{\rho }_{b,r}\widehat{\rho }_{c,r}`$. Separability and measures of entanglement for continuous variable states has been studied . In particular, Duan et al. found a criterion to determine separability of a two-mode Gaussian state. Here, we have a somewhat different approach to find when a two-mode squeezed vacuum in thermal environments is separable and not quantum-mechanically entangled. Our analysis of separability for the mixed squeezed vacuum is extended and fully described for any two-mode Gaussian state in Appendix.
As shown in Appendix, the mixed two-mode squeezed vacuum in the thermal environment is separable when a positive definite $`P`$ function can be assigned to it. The mixed two-mode squeezed vacuum serving the quantum channel can then be written by a statistical mixture of the direct-product states;
$$\widehat{\rho }_{qc}=d^2\beta 𝒫(\beta )\widehat{\rho }_b(\beta )\widehat{\rho }_c(\beta )$$
(16)
where $`𝒫(\beta )`$ is a probability density function.
With use of Eqs. (4) and (13), we find that the mixed two-mode squeezed vacuum is separable when $`n_\tau =1`$ where $`n_\tau `$ is defined as
$`n_\tau (\overline{n},s_{qc},T)`$ $``$ $`\mathrm{\Gamma }\mathrm{\Lambda }`$ (17)
$`=`$ $`\left(2\overline{n}+1\right)T+(1T)\mathrm{exp}(2s_{qc})`$ (18)
according to the condition (52). This is in agreement with Duan et al.’s separation criterion . The pure two-mode squeezed vacuum for $`T=0`$, is never separable unless $`s_{qc}=0`$. For the zero temperature environment, i.e., $`\overline{n}=0`$, the two-mode squeezed state stays quantum-mechanically entangled at any time. For the reasons given in Sec. IV, we call $`n_\tau `$ as the noise factor.
If $`n_\tau <1`$, the quantum channel state is entangled and the teleportation is performed at the quantum level. When $`n_\tau 1`$, the quantum channel is no longer quantum-mechanically entangled. However the inter-mode correlation is still there as $`\mathrm{\Lambda }0`$ in Eq. (13). Questions then arise: Does this classical correlation influence the teleportation? Can any nonclassical properties imposed in an original state be teleported by the classically-correlated channel? Braunstein and Kimble found that when a pure two-mode squeezed state is separable, i.e., $`s_{sq}=0`$, observation of any nonclassical features in the teleported state is precluded. However, when a pure state is separable there is no classical correlation either.
## IV Transfer of Nonclassical Features
An imperfect replica state is reproduced at the receiving station when the quantum channel is not maximally entangled. It is well known that any linear noise-addition process, for example linear dissipation and amplification, can be described by the convolution relation of the quasiprobability functions. With use of the Wigner functions for an arbitrary original state (15) and for an impure quantum channel (13), we find that the equation (8) leads to the following convolution relation
$$W_r(\alpha )=d^2\beta P_\tau (\alpha \beta )W_o(\beta )$$
(19)
where the function $`P_\tau `$ characterizes the teleportation process;
$$P_\tau (\alpha \beta )=\frac{1}{\pi n_\tau }\mathrm{exp}\left(\frac{1}{n_\tau }\left|\alpha \beta \right|^2\right)$$
(20)
and the noise factor $`n_\tau `$, defined in Eq. (17), is completely determined by the characteristics of the quantum channel. The noise factor increases monotonously as the interaction time $`T`$ increases. The larger the initial squeezing, the less vulnerable the quantum channel is.
The noise factor $`n_\tau `$ is related to the fidelity. With use of Eqs. (14) and (19) the fidelity can be written as
$$=\pi d^2\alpha d^2\beta W_o(\alpha )P_\tau (\alpha \beta )W_o(\beta ).$$
(21)
In the limit of $`n_\tau 0`$, the function $`P_\tau (\alpha \beta )`$ in Eq. (20) becomes a delta function and the fidelity becomes unity. The teleportation loses the original information completely with $`=0`$ in the limit of $`n_\tau \mathrm{}`$.
The properties of the nonclassical states have been calculated and illustrated by quasiprobability functions. The nonclassical features are associated especially with negative values and singularity of the quasiprobability $`P`$ function . Suppose an original state whose $`P`$ function is not positive everywhere in phase space. When this state is teleported, its nonclassical features are certainly transferred to the teleported state if the teleportation is perfect. If the teleportation is poor, the teleported state may have its $`P`$ function positive definite and lose the nonclassical features.
By the Fourier transform of Eq. (19), the convolution relation is represented in terms of the characteristic functions as
$$C_r^W(\xi )=\mathrm{exp}(\overline{n}_\tau |\xi |^2)C_o^W(\xi ).$$
(22)
Using the relation (2) between characteristic functions, Eq.(22) is written as
$$C_r^P(\xi )=\mathrm{exp}[(n_\tau 1)|\xi |^2]C_o^Q(\xi ),$$
(23)
where $`C_o^Q(\xi )`$ is the characteristic function for $`R_{\sigma =1}(\alpha )`$ of the original state. The $`P`$ function is not semi-positive definite if its characteristic function $`C_r^P(\xi )`$ is not inverse-Fourier-transformable. Even when it is inverse-Fourier-transformable, there is a chance for the $`P`$ function to become negative at some points of phase space. Lütkenhaus and Barnett found that only when $`\sigma 1`$ the quasiprobability $`R_\sigma (\alpha )`$ for any state is semi-positive definite. We are sure that, for any original state, the left-hand side of Eq. (23) is inverse-Fourier transformed to a $`P`$ function semi-positive definite only when $`n_\tau 1`$. This condition is the same as the separability condition (17) for the quantum channel. We conclude that when a quantum channel is separable, i.e., not quantum-mechanically entangled, no nonclassical features implicit in an original state is transferred by teleportation. In other words, nonclassical features are not teleported via a classically-correlated channel.
There are two well-known nonclassical properties which a continuous-variable state may have: Sub-Poissonian photon statistics and quadrature squeezing. The two nonclassical properties have been studied for noiseless communication. We analyze the transfer of these properties by teleportation in the following subsections.
### A sub-Poissonian statistics and Fock state
A state is defined to be sub-Poissonian when its photon-number variance $`(\mathrm{\Delta }N)^2`$ is smaller than its mean photon number $`\overline{N}`$. The expectation value of an observable for a state can be obtained by use of the characteristic function $`C^P(\xi )`$ for its $`P`$ function ;
$$(\widehat{a}^{})^m\widehat{a}^n=\frac{^m}{\xi ^m}\frac{^n}{(\xi ^{})^n}C^P(\xi )|_{\xi =\xi ^{}=0}.$$
(24)
Substituting Eq. (23) into Eq. (24), we find that the teleported state is sub-Poissonian when the noise factor,
$$n_\tau <\sqrt{\overline{N}_o^2+\overline{N}_o(\mathrm{\Delta }N_o)^2}\overline{N}_o.$$
(25)
where $`\overline{N}_o`$ and $`(\mathrm{\Delta }N_o)^2`$ are the mean photon number and photon-number variance for the original state. If the original state is Poissonian or super-Poissonian, the right-hand side of the inequality is either negative or imaginary so the teleported state is never sub-Poissonian.
Assuming the largest sub-Poissonicity, $`(\mathrm{\Delta }N_o)^2=0`$, for the original state, it is found that when the noise factor $`n_\tau <\sqrt{\overline{N}_o^2+\overline{N}_o}\overline{N}1/2`$, some sub-Poissonian property is found in the teleported state. Thus, if the noise factor of the quantum channel is larger than or equal to 1/2, the transfer of any sub-Poissonian property is precluded.
A Fock state $`|m`$ has a definite energy and its photon-number variance is zero. When this extreme sub-Poissonian field is teleported, the mean photon number and mean variance are $`\overline{N}_r=m+n_\tau `$ and $`\mathrm{\Delta }N_r^2=(2m+1)n_\tau +n_\tau ^2`$ at the receiving station. The Fock state $`|m`$ is written in the Wigner representation as
$`W_o(\alpha ,m)={\displaystyle \frac{2}{\pi }}(1)^m\mathrm{exp}\left(2|\alpha |^2\right)L_m\left(4|\alpha |^2\right).`$ (26)
where $`L_m`$ is a Laguerre polynomial. From the convolution relation (19), the teleported state is obtained as
$`W_r(\alpha )=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{(2n_\tau 1)^m}{(2n_\tau +1)^{m+1}}}\mathrm{exp}\left({\displaystyle \frac{2|\alpha |^2}{2n_\tau +1}}\right)`$ (28)
$`\times L_m\left({\displaystyle \frac{4|\alpha |^2}{(2n_\tau )^21}}\right).`$
The fidelity for the Fock state is given by Eq. (21);
$`_m={\displaystyle \frac{(1n_\tau )^m}{(1+n_\tau )^{m+1}}}P_m\left({\displaystyle \frac{1+n_\tau ^2}{1n_\tau ^2}}\right)`$ (29)
where $`P_m`$ is a Legendre polynomial. When $`n_\tau =0`$, $`_m=1`$. In the limit of $`n_\tau =1`$, where the teleportation is classical, the fidelity $`_m=(1/4)^m`$ for $`m0`$. The vacuum state has the fidelity $`_0=1/2`$ in the limit.
### B quadrature squeezing and squeezed state
We examine the transfer of quadrature squeezing which an unknown original state may have. The quadrature-phase operator is defined as
$$\widehat{X}(\varphi )=e^{i\varphi }\widehat{a}+e^{i\varphi }\widehat{a}^{}$$
(30)
where $`\widehat{a}`$ ($`\widehat{a}^{}`$) is an annihilation (creation) operator and $`\varphi `$ related to the angle in phase space. A state is said to be squeezed if the quadrature-phase variance $`[\mathrm{\Delta }X(\varphi )]^2<1`$ for an angle $`\varphi `$. Substituting Eq. (23) into Eq. (24), the mean quadrature phase $`\overline{X}(\varphi )`$ and variance $`[\mathrm{\Delta }X(\varphi )]^2`$ can be calculated
$$\overline{X}_r(\varphi )=\overline{X}_o(\varphi );[\mathrm{\Delta }X_r(\varphi )]^2=[\mathrm{\Delta }X_o(\varphi )]^2+2n_\tau ,$$
(31)
where $`\overline{X}_o(\varphi )`$ and $`[\mathrm{\Delta }X_o(\varphi )]^2`$ are the mean quadrature phase and variance for the original state. It is interesting to realize that the mean quadrature phase does not change at all during teleportation. This property holds regardless of the channel entanglement.
The teleported state exhibits quadrature squeezing if
$$n_\tau <\frac{1}{2}\left\{1[\mathrm{\Delta }X_u^2(\varphi )]^2\right\}\frac{1}{2}.$$
(32)
Suppose that the original state has the absolute minimum variance $`[\mathrm{\Delta }X_o(\varphi ^{})]^2=0`$ at $`\varphi =\varphi ^{}`$. Then its teleported state is also squeezed if the quantum channel is entangled enough to be represented by the noise factor $`n_\tau <1/2`$. We note that the condition $`n_\tau <1/2`$ applies to the survival of both quadrature squeezing and sub-Poissonian statistics.
A squeezed vacuum with the degree of squeezing $`s_o`$ is written in the Wigner representation as
$$W_o(\alpha )=\frac{2}{\pi }\mathrm{exp}\left[2\mathrm{exp}(2s_o)\alpha _r^22\mathrm{exp}(2s_o)\alpha _i^2\right]$$
(33)
where $`\alpha _r`$ and $`\alpha _i`$ are real and imaginary parts of $`\alpha `$. Its teleported state is represented by the Wigner function
$$W_r(\alpha )=\frac{2}{\pi \sqrt{A(s_o)A(s_o)}}\mathrm{exp}\left[\frac{2}{A(s_o)}\alpha _r^2\frac{2}{A(s_o)}\alpha _i^2\right]$$
(34)
where the parameter $`A(s_o)=2n_\tau +\mathrm{exp}(2s_o)`$. The fidelity is given by
$$(s_o)=\left(n_\tau ^2+2n_\tau \mathrm{cosh}2s_o+1\right)^{1/2}.$$
(35)
When the teleportation is classical with $`n_\tau =1`$, $`(s_o)=(2+2\mathrm{cosh}2s_o)^{1/2}`$.
## V Remarks
Quantum teleportation can be made more reliable by sophisticated schemes such as purification of the impure or partially entangled quantum channel , detection with perfect efficiency, and well-defined unitary operation. However, in the real world, the influence of noise cannot easily be disregarded. We have been interested in the influence of noise on the transfer of nonclassicalities which may be imposed in an original unknown state. To make the problem simple while honoring the real experimental situation, we assumed that the same amount of noise affects the two modes of the quantum channel. We found that when the quantum channel is separable the transfer of any nonclassicality is impossible: Nonclassical features can not be teleported via a classically-correlated channel. The separability of a two-mode Gaussian state is considered using the possibility of assigning a positive well-defined $`P`$ function to the state after some local unitary operations. We have analyzed the transfer of well-known nonclassical features such as sub-Poissonicity and quadrature squeezing. The teleportation of the two nonclassical features is ruled out under the same noise level. The faithfulness of the teleportation has also been discussed and the fidelities have been found for the initial Fock state and squeezed state. Because one of the important ingredients of teleportation is that the original state is unknown at the sending station. Thus our measure of noise factor $`n_\tau `$, which depends only on the quality of the channel, is important. Of course, to represent the quality of the teleportation by a fidelity we have to know the average fidelity for the teleportation, which is under investigation.
One question still arises: Is the teleportation better than the direct transmission to transfer a nonclassical field? A field may be deteriorated by the thermal environment during the direct transmission. Solving a similar Fokker-Planck equation to Eq. (11) for a single-mode field, we find that, by the direct transmission, the Wigner function at the receiving station can be represented by the same equation as Eq. (19) with the different noise factor $`n_d`$ :
$$n_d=\overline{n}T.$$
(36)
Assuming that the imperfect teleportation is caused only by the impure quantum channel embedded in the thermal environment, we compare the two noise factors $`n_\tau `$ in Eq. (17) and $`n_d`$ in Eq. (36). We have implicitly assumed in this paper that the two-mode squeezed state (quantum channel) generator is located in the middle point between the sending and receiving stations. The squeezed photons in the quantum channel, thus, travel only a half distance between the sending and receiving stations. Bearing it in mind, we find that
$`n_\tau (`$ $`\text{for time}t/2)n_d(\text{for time}t)=`$ (38)
$`\overline{n}[1\sqrt{1T}]^2+1\sqrt{1T}[1\mathrm{exp}(2s_{qc})].`$
The right-hand side is semi-positive so that the noise given by teleportation is more than that by direct transmission. If we consider that this result is obtained for the case when the other operations including detection and unitary transformation in the teleportation protocol is perfect, we conclude that the nonclassical field is more robust in direct transmission than in teleportation. The reason is that the teleportation relies on quantum entanglement of the quantum channel. The quantum entanglement based on inter-mode coherence is much more fragile than the single-mode coherence. However, the quantum teleportation can be made more faithful by purification of the quantum channel while the direct transmission does not have that possibility.
###### Acknowledgements.
This work was supported in part by the Brain Korea 21 grant (D-0055) of the Korean Ministry of Education.
## positivity of $`P`$ function and separability for a Gaussian state
A two-mode Gaussian state $`\widehat{\rho }`$ of mode $`b`$ and $`c`$ is separable when it is represented by a statistical mixture of the direct-product states;
$$\widehat{\rho }=d^2\beta 𝒫(\beta )\widehat{\rho }_b(\beta )\widehat{\rho }_c(\beta )$$
(39)
where $`\widehat{\rho }_{b,c}(\beta )`$ are density matrices, respectively, for modes $`b`$ and $`c`$, and $`𝒫(\beta )`$ is a probability density function with $`𝒫(\beta )0`$. The states of $`\widehat{\rho }_b(\beta )`$ and $`\widehat{\rho }_c(\beta )`$ can be nonclassical and do not have to have their $`P`$ functions positive well-defined. However, because they are Gaussian, it is possible to transform them to assign positive well-defined $`P`$ functions by local unitary transformations . The separable condition, (39), can then be written as
$`\widehat{\rho }^{}=`$ $`{\displaystyle d^2\alpha _bd^2\alpha _cd^2\beta 𝒫(\beta )P(\alpha _b;\beta )P(\alpha _c;\beta )}`$ (41)
$`\times |\alpha _b\alpha _b||\alpha _c\alpha _c|`$
where $`P_b(\alpha _b;\beta )`$ and $`P_c(\alpha _c;\beta )`$ are the $`P`$ functions, respectively, for the fields of modes $`b`$ and $`c`$ after some local unitary operations. $`\widehat{\rho }^{}`$ is for the two-mode Gaussian state after the local operations.
We want to prove in this appendix that if and only if when a two-mode Gaussian state is separable, a positive well-defined $`P`$ function $`P(\alpha _b,\alpha _c)`$ is assigned to it after some local unitary transformations.
consider the sufficient condition. If a two-mode Gaussian state $`\widehat{\rho }`$ is separable, it can be written as Eq. (41) after some local operations. Both $`P_b(\alpha _b;\beta )`$ and $`P_c(\alpha _c;\beta )`$ are positive well-defined and $`𝒫(\beta )`$ is a probability density function so
$$d^2\beta 𝒫(\beta )P(\alpha _b;\beta )P(\alpha _c;\beta )$$
(42)
is a normalized positive function, which we can take as the positive well-defined $`P`$ function $`P(\alpha _b,\alpha _c)`$. We have proved that if a two-mode Gaussian state is separable, it has a positive well-defined $`P`$ function after some local unitary operations.
Now let us prove the necessary condition. If the locally-transformed two-mode Gaussian state is represented by a positive well-defined $`P`$ function $`P(\alpha _b,\alpha _c)`$, the separable condition (41) becomes
$$P(\alpha _b,\alpha _c)=d^2\beta 𝒫(\beta )P_b(\alpha _b;\beta )P_c(\alpha _c;\beta ).$$
(43)
Further by some additional squeezing and rotation it is always possible to have the rotationally-symmetric variance $`[\mathrm{\Delta }\alpha _i(\varphi )]^2`$ for any angle $`\varphi `$. After these transformations, the positive well-defined $`P`$ function $`P(\alpha _b,\alpha _c)`$ can be written as
$`P(\alpha _b,\alpha _c)=`$ $`𝒩\mathrm{exp}[{\displaystyle \underset{i,j=b,c}{}}\alpha _iN_{ij}\alpha _j^{}`$ (45)
$`+{\displaystyle \underset{i=b,c}{}}(\alpha _i\lambda _i^{}+\alpha _i^{}\lambda _i)]`$
where $`𝒩`$ is the normalization constant, $`N_{ij}`$ a Hermitian matrix, and $`\lambda _i`$ a complex number. The linear terms of $`\alpha _i`$ are not considered because they do not affect the proof. In fact, they can always be removed by some local displacement operations. Eq. (45) can be written as
$$P(\alpha _b,\alpha _c)=\frac{\mathrm{Det}N_{ij}}{\pi ^2}\mathrm{exp}\left(\underset{i,j=b,c}{}\alpha _iN_{ij}\alpha _j^{}\right)$$
(46)
where $`\mathrm{Det}N_{ij}`$ is the determinant of the Hermitian matrix $`N_{ij}`$. To find an expression in the form of Eq. (43), let us introduce an auxiliary field ($`\beta `$, $`\beta ^{}`$) enabling the function $`P(\alpha _b,\alpha _c)`$ to be represented by a Gaussian integral;
$`P(\alpha _b,\alpha _c)=`$ $`{\displaystyle \frac{\mathrm{Det}N_{ij}}{\pi ^3}}{\displaystyle }d^2\beta \mathrm{exp}(|\beta |^2E_b(\alpha _b,\beta )`$ (48)
$`E_c(\alpha _c,\beta ))`$
where
$`E_b(\alpha _b,\beta )`$ $`=`$ $`\left(N_{bb}+|N_{bc}|^2\right)|\alpha _b|^2\alpha _bN_{bc}\beta ^{}`$ (50)
$`\alpha _b^{}N_{bc}^{}\beta `$
$`E_c(\alpha _c,\beta )`$ $`=`$ $`\left(N_{cc}+1\right)|\alpha _c|^2+\alpha _c\beta ^{}+\alpha _c^{}\beta `$ (51)
The integrand in Eq. (48) can now be decomposed into three Gaussian functions each of which satisfies the normalization condition because
$$N_{ii}>0\mathrm{and}\mathrm{Det}N_{ij}>0$$
(52)
for positive well-defined $`P(\alpha _b,\alpha _c)`$ in Eq. (45). Taking
$`P_b(\alpha _b;\beta )`$ $`=`$ $`{\displaystyle \frac{M_b}{\pi }}\mathrm{exp}(M_b|\alpha _b|^2+\alpha _bN_{bc}\beta ^{}`$ (54)
$`+\alpha _b^{}N_{bc}^{}\beta {\displaystyle \frac{|N_{bc}|^2}{M_b}}|\beta |^2)`$
$`P_c(\alpha _c;\beta )`$ $`=`$ $`{\displaystyle \frac{M_c}{\pi }}\mathrm{exp}(M_c|\alpha _c|^2\alpha _c\beta ^{}\alpha _c^{}\beta `$ (56)
$`{\displaystyle \frac{1}{M_c}}|\beta |^2)`$
$`𝒫(\beta )`$ $`=`$ $`{\displaystyle \frac{M_s}{\pi }}\mathrm{exp}\left(M_s|\beta |^2\right)`$ (57)
where $`M_b=N_{bb}+|N_{bc}|^2`$, $`M_c=N_{cc}+1`$, and $`M_s=\mathrm{Det}N_{ij}/(M_bM_c)`$, the $`P`$ function is finally obtained in the form of Eq. (43). It is clear that $`𝒫(\beta )`$ is the positive probability density function and the two-mode Gaussian state is separable if it can be transformed to have a positive well-defined $`P`$ function by some local unitary operations. |
warning/0003/hep-th0003090.html | ar5iv | text | # Testing SDLCQ in 2+1 dimensions
## 1 Introduction
The motivations to consider $`𝒩=1`$ supersymmetric Yang-Mills theories in 2+1 dimensions are manifold. For one, there is recent progress in understanding the properties of strongly coupled gauge theories with supersymmetry $`^\mathrm{?}`$, some of which are believed to be interconnected through a web of strong-weak coupling dualities. There is a need to study these issues at a fundamental level. Ideally, we would like to solve for the bound states of these theories directly, and at any coupling. However, solving a field theory from first principles is typically an intractable task. Nevertheless, it has been known for some time that $`1+1`$ dimensional field theories can be solved from first principles via a straightforward application of DLCQ $`^\mathrm{?}`$. Recently, a large class of supersymmetric gauge theories in two dimensions was studied using a supersymmetric form of DLCQ (SDLCQ), which is known to preserve supersymmetry at every stage of the calculation $`^{\mathrm{?},\mathrm{?}}`$. It turns out that this formalism can be applied to higher-dimensional theories $`^\mathrm{?}`$. This is interesting, because in higher dimensions, due to the additional scale, theories have the potential of exhibiting a complex phase structure, which may include a strong-weak coupling duality.
## 2 SDLCQ
We consider a three dimensional SU($`N_c`$) $`𝒩=1`$ super-Yang-Mills theory compactified on the space-time $`𝐑\times S^1\times S^1`$. The calculations are performed in the large $`N_c`$ limit. In particular, we use light-cone coordinates $`x^\pm =\frac{1}{\sqrt{2}}(x^0\pm x^1)`$, compactify $`x^{}`$ on a light-like circle a la DLCQ, and wrap the remaining transverse coordinate $`x^{}`$ on a spatial circle. We are able to solve for bound state wave functions and masses numerically by diagonalizing the discretized light-cone supercharge. This procedure preserves supersymmetry at every step. The action of $`𝒩=1`$ SYM(2+1) is
$$S=d^2x_0^L𝑑x_{}\text{tr}(\frac{1}{4}F^{\mu \nu }F_{\mu \nu }+\mathrm{i}\overline{\mathrm{\Psi }}\gamma ^\mu D_\mu \mathrm{\Psi }).$$
We decompose the spinor $`\mathrm{\Psi }`$ in terms of chiral projections $`\psi ,\chi `$ and choose the light-cone gauge $`A^+=0`$. We solve for the non-dynamical fields $`A^{}`$ and $`\chi `$ and formulate the momentum operators in the physical degrees of freedom ($`\varphi A^2`$)
$`P^+`$ $`=`$ $`{\displaystyle 𝑑x^{}_0^L𝑑x_{}\text{tr}\left[(_{}\varphi )^2+\mathrm{i}\psi _{}\psi \right]},`$
$`P^{}`$ $`=`$ $`{\displaystyle 𝑑x^{}_0^L𝑑x_{}\text{tr}\left[\frac{1}{2}J\frac{1}{_{}^2}J\frac{\mathrm{i}}{2}D_{}\psi \frac{1}{_{}}D_{}\psi \right]}.`$
The canonical commutation relations yield the supersymmetry algebra
$$\{Q^+,Q^+\}=2\sqrt{2}P^+,\{Q^{},Q^{}\}=2\sqrt{2}P^{},\{Q^+,Q^{}\}=4P_{}.$$
We use the standard decomposition of the fields $`\varphi _{ij}(x^{},x^{})`$ and $`\psi _{ij}(x^{},x^{})`$ into momentum modes $`a_{ij}^{}(\underset{¯}{k})`$ and $`b_{ij}^{}(\underset{¯}{k})`$, respectively. We defined $`\underset{¯}{k}(k^+,n^{})`$ for convenience. The (anti-)commutation relations
$`[a_{ij}(\underset{¯}{k}),a_{lk}^{}(\underset{¯}{k}^{})]=\{b_{ij}(\underset{¯}{k}),b_{lk}^{}(\underset{¯}{k}^{})\}=\delta (k^+k^{}_{}{}^{}+)\delta _{n^{},n^{{}_{}{}^{}}}\delta _{il}\delta _{jk}.`$
yield the explicit form of the supercharges, which are listed in Ref. $`^\mathrm{?}`$. For the present discussion it suffices to know that the structure of $`Q^{}`$ is
$`Q^{}`$ $`=`$ $`{\displaystyle \frac{2^{3/4}\pi \mathrm{i}}{L}}{\displaystyle \underset{n^{}𝐙}{}}{\displaystyle _0^{\mathrm{}}}𝑑k^+{\displaystyle \frac{n^{}}{\sqrt{k^+}}}\left[a_{ij}^{}(\underset{¯}{k})b_{ij}(\underset{¯}{k})b_{ij}^{}(\underset{¯}{k})a_{ij}(\underset{¯}{k})\right]+g\stackrel{~}{Q}^{},`$
where $`\stackrel{~}{Q}^{}`$ contains the terms with three operators. We note also that the supercharge $`Q^{}`$ is linear in the coupling $`g`$, and thus the Hamiltonian $`P^{}`$ is quadratic in $`g`$.
We now perform the truncation procedure. The harmonic resolution $`K`$ plays the role of a longitudinal cutoff as usual, and longitudinal momentum fractions take values $`\frac{k_i^+}{P^+}=\frac{n_i}{K},n_i=1,2,\mathrm{},K`$. The transverse cutoff $`T`$ allows for momenta $`k_i^{}=2\pi n_i^{}/L`$ with $`n_i^{}=0,\pm 1,\pm 2,\mathrm{}\pm T`$. This prescription preserves parity symmetry in transverse directions. How does such a truncation affect the supersymmetry properties of the theory? Note first that an operator relation $`[A,B]=C`$ in the full theory is not expected to hold in the truncated formulation. However, if A is quadratic in terms of fields (or in terms of creation and annihilation operators), one can show that the relation $`[A,B]=C`$ implies $`[A_{tr},B_{tr}]=C_{tr}`$ for the truncated operators $`A_{tr}`$,$`B_{tr}`$, and $`C_{tr}`$. In our case, $`Q^+`$ is quadratic, and so the relations $`\{Q_{tr}^+,Q_{tr}^+\}=2\sqrt{2}P_{tr}^+`$ and $`\{Q_{tr}^+,Q_{tr}^{}\}=0`$ are true in the $`P_{}=0`$ sector of the truncated theory. The anticommutator $`\{Q_{tr}^{},Q_{tr}^{}\}`$, however, is not equal to $`2\sqrt{2}P_{tr}^{}`$. So the diagonalization of $`\{Q_{tr}^{},Q_{tr}^{}\}`$ will yield a different bound-state spectrum than the one obtained after diagonalizing $`2\sqrt{2}P_{tr}^{}`$. Of course, the two spectra should agree in the limit $`T\mathrm{}`$. At any finite truncation, however, we have the liberty to diagonalize either of these operators. The choice of $`\{Q_{tr}^{},Q_{tr}^{}\}`$ specifies our regularization scheme. Choosing to diagonalize the light-cone supercharge $`Q_{tr}^{}`$ has an important advantage: the spectrum is exactly supersymmetric for any truncation. In contrast, the spectrum of the Hamiltonian $`P_{tr}^{}`$ becomes supersymmetric only in the infinite resolution limit.
Let us take a look at the discrete symmetries of $`Q^{}`$. The three commuting symmetries $`Z_2`$ are parity in the transverse direction
$$P:a_{ij}(k,n^{})a_{ij}(k,n^{}),b_{ij}(k,n^{})b_{ij}(k,n^{}),$$
which anti-commutes with $`Q^+`$ and $`P_{}`$, and a generalized $`T`$ symmetry
$$S:a_{ij}(k,n^{})a_{ji}(k,n^{}),b_{ij}(k,n^{})b_{ji}(k,n^{}).$$
To close the group, we need a third symmetry, namely $`R=PS`$. An interesting detail of the symmetry considerations is the fact that the $`P`$ symmetry leads to exactly degenerate eigenvalues. This means in turn that all massive eigenvalues are four-fold degenerate. The argument goes as follows. Start with a massive state with positive parity $`|M+`$ which obeys
$$(Q^{})^2|M+=M^2|M+,P|M+=+|M+.$$
Then $`Q^+Q^{}|M+`$ is a state with same mass but opposite parity
$$PQ^+Q^{}|M+=Q^+Q^{}P|M+=Q^+Q^{}|M+.$$
## 3 Numerical Results
With the truncation prescription described above, we can solve the discretized eigenvalue problem
$$2P^+P^{}|\psi =M^2|\psi ,$$
characterized by the cutoffs $`K`$ and $`T`$, on the computer. This is equivalent to constructing the supercharge $`Q^{}`$ in the usual Fock basis, and then diagonalizing it. If the resulting mass (squared) eigenvalues $`M^2`$ are plotted as a function of the dimensionless coupling $`g^{}=g\sqrt{NL/4\pi ^3}`$, several striking features emerge. Namely, as was noted in Ref. $`^\mathrm{?}`$, one finds a very stable strong-coupling spectrum. Secondly, we find states which fall off fast to zero mass with increasing coupling. Since the previous work $`^\mathrm{?}`$ was a calculation of the spectrum at $`T1`$, it is natural to ask whether the well defined large $`g^{}`$ spectrum survives at $`T\mathrm{}`$, and to study the properties of the states with masses decreasing at large $`g^{}`$. A further question is if the number of massless states is independent of the transverse cutoff $`T`$.
Our previous SDLCQ calculations were done using a code written in Mathematica and performed on a PC. This code has now been rewritten in C++ and some of the present work was done on supercomputers. We were able to perform numerical diagonalizations for $`K=2`$ through 7 and for values of $`T`$ up to $`T=9`$ at $`K=4`$ and $`T=1`$ at $`K=7`$.
Massive spectrum: Little is known about the large coupling spectrum of quantum field theories, with the exception of theories in 1+1 dimensions. There, however, the concept of large coupling has no meaning, since the coupling is only a multiplicative constant in the Hamiltonian. In particular, there is no weak/strong duality known in $`𝒩=1`$ SYM(2+1), which could give some clue how the spectrum looks like. We performed therefore a non-perturbative calculation in SDLCQ to directly access the spectrum.
In Fig. 1 we plot the bound state masses squared $`M^2`$ as a function of the transverse resolution $`T`$ for $`K=4`$ and $`K=5`$. We see that the curves are very flat, thus exhibiting fast convergence in transverse cutoff $`T`$. The continuum result can be obtained by extrapolating the curves to $`1/T0`$. Let us look at the bound states as a function of the coupling $`g^{}`$, focusing on the large coupling regime. We see, Fig. 2, that the states are extremely stable in $`g^{}`$, i.e. they are quasi independent of coupling. We find this behavior at every value of $`K`$, and even irrespective of the value of the transverse cutoff $`T`$.
We show the bound state mass as a function of $`1/K`$ in Fig. 3(a). These results are the first calculation of the strong-coupling bound states of $`𝒩=1`$ SYM in $`2+1`$ dimensions. As we increase the resolution we are able to see states that have, as their primary component, more and more partons, and, as we have seen in other supersymmetric theories, many of these states appear at low energies. This accumulation of high-multiplicity low-mass states appears to be a unique property of SUSY theories. In non-SUSY theories the new states appear at increasing energies. In the dimensionally reduced version of this theory we saw that the accumulation point of these low-mass states appeared to be at zero mass $`^{\mathrm{?},\mathrm{?}}`$. Here again we see clear evidence of an accumulation of low mass states, however we don’t have sufficient information to say whether an accumulation point exists.
Massless states: There are two kinds of massless states in the spectrum. Firstly, the states massless for $`g^{}0`$. They are massless because at small coupling, only first term of the supercharge, Eq. (2), gives a contribution. Then all partons with $`n^{}=0`$ (anti-)commute with $`Q^{}`$. Thus all states made out of these partons are massless, and massless states contain just these partons. The set of massless states at $`g^{}=0`$ therefore coincides with a Hilbert space of the theory dimensionally reduced to $`1+1`$. Moreover, the whole infrared spectrum of SYM$`2+1`$ at small coupling is governed by the dimensionally reduced theory (see $`^\mathrm{?}`$ for details) Secondly, we see states that are exactly massless at any coupling. These are $`2(K1)`$ BPS states, fulfilling $`Q^{}|m=0=0`$, $`Q^+|m=00.`$ It is therefore easy to construct the massless states of (2+1) theories, at least at large $`N_c`$.
Unphysical states: Finally, we can unambiguously detect unphysical states due to their special properties. Namely, these states vary strongly with coupling $`g^{}`$ and they appear (predominantly) for $`K`$ odd and decouple for $`T\mathrm{}`$, see Fig. 3. It is thus easy to classify and to remove the unphysical states from the spectrum.
## 4 Conclusions
We have shown that the SDLCQ formalism naturally extends to higher dimensions. Rapid convergence in transverse direction is found. Concerning the specific theory, we obtained the strong coupling spectrum of $`𝒩=1`$ SYM(2+1). The bound states are extremely stable for $`g^{}\mathrm{}`$. The unphysical states in the spectrum can be identified and removed. We found no new massless states at strong coupling compared to previous work $`^\mathrm{?}`$. The massless sector of the theory is completely determined by the dimensional reduced model. The light states turn out to be string-like and might contain physics of dual theories. Also, the theory might be conformal at decompactification limit. We are currently working on analytical and numerical improvements of the approach. We expect to be able to address problems like the very interesting $`𝒩=4`$ SYM in 3+1 and $`𝒩=1`$ SYM(2+1) including a Chern-Simons term conjectured to break supersymmetry, in the near future.
## Acknowledgments
It is a pleasure to thank the workshop organizers for hospitality and support.
## References |
warning/0003/astro-ph0003180.html | ar5iv | text | # The first radius-expansion X-ray burst from GX 3+1
## 1 Introduction
The overall X-ray intensity of the low-mass X-ray binary (LMXB) GX 3+1 varies slowly on time scales of months to years (Makishima et al. 1983; Asai et al. 1993, see also Fig. 1). X-ray bursts in GX 3+1 were discovered by Hakucho, at a time when the persistent X-ray flux was about half of that seen previously (Makishima et al. 1983). During that time roughly one burst per day was observed. The bursts from this source were shown to be thermonuclear flashes on the neutron star surface, i.e. being of type I (Makishima et al. 1983; Asai et al. 1993: Ginga; Molkov et al. 1999: GRANAT/ART-P), but none of them showed evidence for photospheric radius expansion.
GX 3+1 is one of the four brightest so-called “atoll” sources (Hasinger & van der Klis 1989). The sources in this group (including GX 13+1, GX 9+1 and GX 9+9) hardly show any X-ray bursts (if at all), and display properties like those of other atoll sources when these are in their high accretion rate state: their tracks in X-ray colour-colour diagrams are long, diagonal and slightly curved, while their fast timing properties are at all times dominated by a relatively weak (1–4% rms) power-law shaped noise component. Detailed X-ray spectral modeling seems to suggest that they accrete with rates near 10% of the Eddington mass accretion rate, i.e. intermediate between the more frequently bursting atoll sources and that of the so-called “Z” sources (Psaltis & Lamb 1998). At low accretion rates (and therefore probably low intensities) such sources are predicted to display the properties characteristic of the more frequently bursting atoll sources, which in view of the Hakucho result (see above) at least GX 3+1 seems to satisfy.
During one of our series of target of opportunity observations with RXTE aimed at observing GX 3+1 at low intensities, we observed a strong ($``$2.3 Crab \[2–10 keV\] at maximum) and short (15–20 s) X-ray burst. The burst onset occurred on 1999 August 10, 18:35:53.5 UTC. In this paper we discuss its properties.
## 2 Observations and Analysis
Data were acquired with the Proportional Counter Array (PCA; Bradt et al. 1993) in various observation modes. During our observation from 1999 August 10, 17:15 to August 11, 00:00 UTC, only three units were active, i.e. PCU0, PCU2 and PCU3. For the spectral analysis of the persistent emission we used data collected in 16 s intervals with 129 spectral channels. We accumulated data stretches of 96 s just before and after the burst, combining the three PCU’s. In order to study the burst properties we used a mode which provides 64 spectral channels at a time resolution of 16 $`\mu `$s; this mode combines information from all PCU’s. Spectra during the burst were determined every 0.25 sec during the first 10 s, and every 0.5 s for the remainder. All spectra were corrected for background and dead-time using the procedures supplied by the RXTE Guest Observer Facility. A systematic uncertainty of 1% was taken into account. For our spectral fits we confined the energy range to 2.9–20 keV. The hydrogen column density, N<sub>H</sub>, towards GX 3+1 was fixed to that found by the Einstein SSS and MPC measurements (1.7$``$10<sup>22</sup> atoms cm<sup>-2</sup>, Christian & Swank 1997).
Large amplitude, high coherence brightness oscillations have been observed during various X-ray bursts in other LMXBs (see Strohmayer 1998, 2000). In our search for possible burst oscillations we made fast Fourier transforms (FFTs) of data segments ranging from 0.25 s to 2 s long during the burst, with time steps of 0.125 s (so-called sliding FFTs). We used the 64 spectral channel and 16 $`\mu `$sec data set, and limited our search to the 50 to 2048 Hz frequency range. We performed the search in the whole PCA energy range (2–60 keV), as well as in a relatively high energy range (8–60 keV).
## 3 Results
### 3.1 Temporal behaviour
The light curve of the burst at low energies is single-peaked, whereas at high energies it is double-peaked (Fig. 2a,b). The corresponding hardness curve (Fig. 2c) shows that the burst first hardens, softens, hardens again, and then gradually softens again. Our search for burst oscillations yielded negative results. Previous burst oscillations were found mainly during the rise to burst maximum and after the radius-expansion phase (see Strohmayer 1998, 2000). Using the 0.25 s long FFTs in the full energy range, we derive upper limits on the modulation amplitude of burst oscillations of $``$45% at the start of the rise, $``$15% at the maximum of the burst, and $``$20% just after the radius expansion phase. The upper limits in the 8–60 keV energy band are $``$60%, $``$20%, and $``$25%, respectively. The 2 s long FFTs give more stringent upper limits of $``$16%, $``$6%, and $``$9%, respectively, for the full energy range, whereas we derive $``$24%, $``$7%, and $``$10%, respectively, in the 8–60 keV energy band.
### 3.2 Spectral behaviour
The net burst spectra (i.e. total burst spectrum minus persistent spectrum) were satisfactorily ($`\chi _{\mathrm{red}}^2`$ less than $``$2) modeled by black-body emission. The results are shown in Fig. 3. A dip in the black-body temperature, T<sub>bb</sub>, $``$2 s after the burst onset is apparent, simultaneous with an increase of a factor of $``$2 in the black-body radius, R<sub>bb</sub>. The total increase/decrease phase of R<sub>bb</sub> lasts only $``$1.5 sec.
We note that the X-ray spectral analysis during bursts can be significantly hampered when the persistent emission contains a black-body contribution from the same surface of the neutron star that emits the burst emission (van Paradijs & Lewin 1985). In that case our spectral fits to the net burst spectra may contain systematic errors in the black-body temperature and radius, especially near the end of the burst. We therefore repeated our spectral analysis to the total burst spectra, fixing the non black-body component in the persistent emission (see Table 1). The black-body component, which now includes all emission from the neutron star surface, is left free. This procedure only slightly alters our estimated black-body flux. The absence of significant differences between the two methods is mainly due the fact that the burst is sufficiently stronger than the persistent emission (which is reflected by the burst parameter $`\gamma `$, see below), as also noted by Asai et al. (1993).
The persistent emission just before and after the burst can be satisfactorily modeled by a black-body plus a cut-off power-law component (Table 1). Using the X-ray spectral fits we can determine the peak flux (i.e. including persistent emission), F<sub>peak</sub>, and the total burst fluence (i.e. the integrated net burst flux), E<sub>b</sub>, and hence the burst parameters $`\gamma `$ (=F<sub>pers</sub>/\[F<sub>peak</sub>-F<sub>pers</sub>\]) and $`\tau `$ (=E<sub>b</sub>/F<sub>peak</sub>). For the burst parameter $`\alpha `$ (=F<sub>pers</sub>/(E<sub>b</sub>/$`\mathrm{\Delta }`$t)), where $`\mathrm{\Delta }`$t is the time since the previous burst) we can only give a lower limit, since the source is not observed during South Atlantic Anomaly passages and earth occultations. However, for a crude estimate we also used the mean burst rate of $``$2 per day, as observed during the 1999 August to October BeppoSAX Wide Field Camera campaign (Muller et al. 2000, in preparation). All burst parameters are also shown in Table 1.
## 4 Discussion
The light curve and X-ray spectral behaviour of the X-ray burst in GX 3+1 observed with RXTE show clear evidence for radius expansion of the neutron star photosphere due to near-Eddington luminosities during a themonuclear runaway on the neutron star surface (for a review see e.g. Lewin et al. 1993). The total time for the expansion and contraction phase is only $``$1.5 s, during which the radius expanded only by a factor of $``$2. Such short bursts with small expansion phases have been seen in other bright X-ray sources, such as Cyg X-2 (Smale 1998). The gradual softening at the end of the burst is attributed to cooling of the neutron star surface, which is characteristic for type-I bursts (Hoffman et al. 1978).
During the burst our derived black-body temperatures are smaller, whereas our inferred black-body radii (all at the same assumed distance) are larger, than reported for previous GX 3+1 bursts (Makishima et al. 1983, see also Inoue et al. 1981; Asai et al. 1993; Molkov et al. 1999). The burst parameter values for $`\gamma `$ ($``$0.10–0.20) and $`\tau `$ (4–8 s) quoted by Asai et al. (1993), and inferred from the observations presented by Makishima et al. (1983) and Molkov et al. (1999), are similar to our findings. We note (see also Asai et al. 1993) that $`\gamma `$, $`\tau `$ and our estimate of $`\alpha `$ fall on the extreme end of relations between $`\tau `$ vs. $`\gamma `$ and $`\alpha `$ vs. $`\gamma `$ as presented by Van Paradijs et al. (1988) for typical type I bursters. This shows that if bright sources burst, the burst duration tends to be short (order of 10 s; note however, that some bursts in the bright “Z” source GX 17+2 have a duration of the order of minutes, see e.g. Kuulkers et al. 1997 and references therein). It is interesting to note that our X-ray burst from GX 3+1 is very similar to the radius expansion burst seen in Cyg X-2 with RXTE in most of its facets, except notably for the $`\gamma `$ being a factor 4.3 larger for Cyg X-2 (Smale 1998). Note also that during the burst from Cyg X-2 no evidence for pulsations was reported, similar to what we conclude for GX 3+1, both with upper limits on the modulation strength which are significantly lower than for bursts during which oscillations were seen (see Strohmayer 1998, 2000).
A convenient way to display the burst properties as they vary in time, is a flux-temperature diagram, see Fig. 4. In such a diagram the phases of expansion/contraction and subsequent cooling of the neutron star photosphere are distinguished by two separate tracks (see e.g. Lewin et al. 1993). GX 3+1 moves from the middle bottom to top left (rising phase), top middle (expansion phase), back towards top left (contraction phase), and finally to the lower right part of the diagram (cooling phase). We can adequately fit ($`\chi _{\mathrm{red}}^2`$/dof = 0.9/34) $`\mathrm{log}\mathrm{F}_{\mathrm{bol}}`$ versus $`\mathrm{log}\mathrm{T}_{\mathrm{bb}}`$ during the cooling phase of the burst by a straight line with a slope of 3.97$`\pm `$0.15 (dotted line in Fig. 4). This means that F<sub>bol</sub> is consistent with being proportional to T$`{}_{}{}^{4}{}_{\mathrm{bb}}{}^{}`$, which indicates that the neutron star photosphere radiates as a black-body during the cooling phase, at a constant radius R<sub>bb</sub>. We note that burst spectra are generally not described by pure black-body radiation, especially near the Eddington limit (see Lewin et al. 1993, and references therein). Instead the black-body radiation is modified mainly at energies below $``$3 keV and above $``$10 keV. Since our burst spectra are analysed in the 2.9–20 keV energy range, we are, therefore, not greatly affected by modified black-body radiation. Note that we then probably underestimate our bolometric fluxes.
At the start of the expansion phase the black-body bolometric flux and temperature values do not match those at the end of the contraction phase. Our estimated emission areas are the same at these instants; the above then means that the photosphere is cooler at the end of the contraction phase with respect to the start of the expansion phase. From Fig. 4 we see that F<sub>bol</sub> drops below the constant peak flux before the end of the contraction phase. We infer that the cooling phase therefore already started before the end of the contraction phase.
Using the fact that during the expansion and contraction phase of the neutron star photosphere the burst luminosity equals the Eddington luminosity one can get an estimate of the distance (see e.g. Lewin et al. 1993). Assuming standard burst paramaters (isotropy, cosmic abundances and a canonical neutron star mass of 1.4 M) and taking into account gravitational redshift effects we find $`d=4.5\pm 0.1`$ kpc. If we assume a neutron star mass of 2.0 M we instead find $`d=5.1\pm 0.1`$ kpc. For bright sources like GX 3+1 most of the hydrogen content is being burned persistently, so during the expansion/contraction phase the neutron star atmosphere is likely to lack hydrogen. Using the Eddington luminosity appropriate for hydrogen-poor matter then leads to $`d=6.1\pm 0.1`$ kpc. Dropping only our assumption that the burst radiates isotropically, and assuming anisotropy values of $`0.5<\xi <2`$ (e.g. van Paradijs & Lewin 1987), we derive distances between 3–7 kpc. On the other hand, if the peak luminosities during radius expansion bursts are standard candles we can use the mean peak luminosity for such bursts seen in globular cluster sources for which the distances are known, i.e. $`3.0\times 10^{38}`$ erg s<sup>-1</sup> (Lewin et al. 1993). In this case we derive $`d5.6`$ kpc. These distance estimates show that in principle one can get an idea of the distance to the source, but the exact value still remains rather uncertain by $``$30%.
The persistent flux during our observations is the same within a factor of $``$2 with respect to the previous reports when GX 3+1 was bursting, i.e. low ($``$0.2 Crab). Using the fact that during the peak of the burst the observed (net-burst) luminosity is at near Eddington values we can now for the first time estimate the persistent flux in terms of the Eddington luminosity for the bright atoll sources like GX 3+1 (i.e. GX 13+1, GX 9+1 and GX 9+9). For GX 3+1 we find L$`{}_{\mathrm{pers}}{}^{}0.17`$ L<sub>edd</sub> (assuming the burst and persistent emission is radiated in the same directions). This is consistent with that inferred from models of X-ray spectra, i.e. $``$0.1 L<sub>Edd</sub> (Psaltis & Lamb 1998).
GX 13+1 has been seen to burst sporadically (Matsuba et al. 1995), whereas no bursts have been reported for GX 9+1 and GX 9+9. This may mean that GX 3+1 and GX 13+1 are accreting near to the critical mass accretion rate at which bursts cease to occur, whereas GX 9+9 and GX 9+1 accrete above this limit. However, this does not explain the fact that some sources that are accreting at even higher rates (near Eddington), i.e. Cyg X-2 and GX 17+2, also irreglularly show bursts.
###### Acknowledgements.
We thank Lars Bildsten, Mariano Méndez and Dimitrios Psaltis for discussions. We acknowledge the use of daily averaged quick-look results provided by the ASM/RXTE team. This work was supported in part by the Netherlands Organization for Scientific Research (NWO) grant 614-51-002. |
warning/0003/hep-th0003139.html | ar5iv | text | # Chaos in Superstring Cosmology
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## Abstract
It is shown that the general solution near a spacelike singularity of the Einstein-dilaton-$`p`$-form field equations relevant to superstring theories and $`M`$-theory exhibits an oscillatory behaviour of the Belinskii-Khalatnikov-Lifshitz type. String dualities play a significant role in the analysis.
\]
An outstanding result in theoretical cosmology has been the discovery by Belinskii, Khalatnikov and Lifshitz (BKL) that the generic solution of the four-dimensional Einstein’s vacuum equations near a cosmological singularity exhibits a never ending oscillatory behaviour . The oscillatory approach toward the singularity has the character of a random process, whose chaotic nature has been intensively studied . However, two results cast a doubt on the physical applicability, to our universe, of the BKL picture. First, it was surprisingly found that the chaotic BKL oscillatory behaviour disappears from the generic solution of the vacuum Einstein equations in spacetime dimension $`D11`$ and is replaced by a monotonic Kasner-like power-law behaviour . Second, it was proved that the generic solution of the four-dimensional Einstein-scalar equations also exhibits a non-oscillatory, power-law behaviour , .
Superstring theory suggests that the massless (bosonic) degrees of freedom which can be generically excited near a cosmological singularity correspond to a high-dimension ($`D=10`$ or $`11`$) Kaluza-Klein-type model containing, in addition to Einstein’s $`D`$-dimensional gravity, several other fields (scalars, vectors and/or forms). In view of the results quoted above, it is a priori unclear whether the full field content of superstring theory will imply, as generic cosmological solution, a chaotic BKL-like behaviour, or a monotonic Kasner-like one. In this letter we report the result that the massless bosonic content of all superstring models ($`D=10`$ IIA, IIB, I, $`\mathrm{het}_\mathrm{E}`$, $`\mathrm{het}_{\mathrm{SO}}`$), as well as of $`M`$-theory ($`D=11`$ supergravity), generically implies a chaotic BKL-like oscillatory behaviour near a cosmological singularity. \[Our analysis applies at scales large enough to excite all Kaluza-Klein-type modes, but small enough to be able to neglect the stringy and non-perturbative massive states.\] It is the presence of various form fields (e.g. the three form in $`\mathrm{SUGRA}_{11}`$) which provides the crucial source of this generic oscillatory behaviour.
Let us consider a model of the general form
$`S`$ $`=`$ $`{\displaystyle }d^Dx\sqrt{g}[R(g)_\mu \phi ^\mu \phi `$ (2)
$`{\displaystyle \underset{p}{}}{\displaystyle \frac{1}{(p+1)!}}e^{\lambda _p\phi }(dA_p)^2].`$
Here, the spacetime dimension $`D`$ is left unspecified. We work (as a convenient common formulation) in the Einstein conformal frame, and we normalize the kinetic term of the “dilaton” $`\phi `$ with a weight 1 with respect to the Ricci scalar. The integer $`p0`$ labels the various $`p`$-forms $`A_pA_{\mu _1\mathrm{}\mu _p}`$ present in the theory, with field strengths $`F_{p+1}dA_p`$, i.e. $`F_{\mu _0\mu _1\mathrm{}\mu _p}=_{\mu _0}A_{\mu _1\mathrm{}\mu _p}\pm p`$ permutations. The real parameter $`\lambda _p`$ plays the crucial role of measuring the strength of the coupling of the dilaton to the $`p`$-form $`A_p`$ (in the Einstein frame). When $`p=0`$, we assume that $`\lambda _00`$ (this is the case in type IIB where there is a second scalar). The Einstein metric $`g_{\mu \nu }`$ is used to lower or raise all indices in Eq. (2) ($`gdetg_{\mu \nu }`$). The model (2) is, as it reads, not quite general enough to represent in detail all the superstring actions. Indeed, it lacks additional terms involving possible couplings between the form fields (e.g. Yang-Mills couplings for $`p=1`$ multiplets, Chern-Simons terms, $`(dC_2C_0dB_2)^2`$-type terms in type IIB). However, we have verified in all relevant cases that these additional terms do not qualitatively modify the BKL behaviour to be discussed below. On the other hand, in the case of $`M`$-theory (i.e. $`D=11`$ SUGRA) the dilaton $`\phi `$ is absent, and one must cancell its contributions to the dynamics.
The leading Kasner-like approximation to the solution of the field equations for $`g_{\mu \nu }`$ and $`\phi `$ derived from (2) is, as usual ,
$`g_{\mu \nu }dx^\mu dx^\nu `$ $``$ $`dt^2+{\displaystyle \underset{i=1}{\overset{d}{}}}t^{2p_i(x)}(\omega ^i)^2,`$ (4)
$`\phi `$ $``$ $`p_\phi (x)\mathrm{ln}t+\psi (x),`$ (6)
where $`dD1`$ denotes the spatial dimension and where $`\omega ^i(x)=e_j^i(x)dx^j`$ is a time-independent $`d`$-bein. The spatially dependent Kasner exponents $`p_i(x)`$, $`p_\phi (x)`$ must satisfy the famous Kasner constraints (modified by the presence of the dilaton):
$`Q(p)p_\phi ^2+{\displaystyle \underset{i=1}{\overset{d}{}}}p_i^2\left({\displaystyle \underset{i=1}{\overset{d}{}}}p_i\right)^2`$ $`=`$ $`0,`$ (8)
$`{\displaystyle \underset{i=1}{\overset{d}{}}}p_i`$ $`=`$ $`1.`$ (10)
The set of parameters satisfying Eqs. (Chaos in Superstring Cosmology) is topologically a $`(d1)`$-dimensional sphere: the “Kasner sphere”. When the dilaton is absent, one must set $`p_\phi `$ to zero in Eq.(8). In that case the dimension of the Kasner sphere is $`d2=D3`$.
The approximate solution Eqs. (Chaos in Superstring Cosmology) is obtained by neglecting in the field equations for $`g_{\mu \nu }`$ and $`\phi `$: (i) the effect of the spatial derivatives of $`g_{\mu \nu }`$ and $`\phi `$, and (ii) the contributions of the various $`p`$-form fields $`A_p`$. The condition for the “stability” of the solution (Chaos in Superstring Cosmology), i.e. for the absence of BKL oscillations at $`t0`$, is that the inclusion in the field equations of the discarded contributions (i) and (ii) (computed within the assumption (Chaos in Superstring Cosmology)) be fractionally negligible as $`t0`$. As usual, the fractional effect of the spatial derivatives of $`\phi `$ is found to be negligible, while the fractional effect (with respect to the leading terms, which are $`t^2`$) of the spatial derivatives of the metric, i.e. the quantities $`t^2\overline{R}_j^i`$ (where $`\overline{R}_j^i`$ denotes the $`d`$-dimensional Ricci tensor) contains, as only “dangerous terms” when $`t0`$ a sum of terms $`t^{2g_{ijk}}`$, where the gravitational exponents $`g_{ijk}`$ ($`ij`$, $`ik`$, $`jk`$) are the following combinations of the Kasner exponents
$$g_{ijk}(p)=2p_i+\underset{\mathrm{}i,j,k}{}p_{\mathrm{}}=1+p_ip_jp_k.$$
(11)
The “gravitational” stability condition is that all the exponents $`g_{ijk}(p)`$ be positive. In the presence of form fields $`A_p`$ there are additional stability conditions related to the contributions of the form fields to the Einstein-dilaton equations. They are obtained by solving, à la BKL, the $`p`$-form field equations in the background (Chaos in Superstring Cosmology) and then estimating the corresponding “dangerous” terms in the Einstein field equations. When neglecting the spatial derivatives in the Maxwell equations in first-order form $`dF=0`$ and $`\delta (e^{\lambda _p\phi }F)=0`$, where $`\delta d`$ is the (Hodge) dual of the Cartan differential $`d`$ and $`F_{p+1}=dA_p`$, one finds that both the “electric” components $`\sqrt{g}e^{\lambda _p\phi }F^{0i_1\mathrm{}i_p}`$, and the “magnetic” components $`F_{j_1\mathrm{}j_{p+1}}`$, are constant in time. Combining this information with the approximate results (Chaos in Superstring Cosmology) one can estimate the fractional effect of the $`p`$-form contributions in the right-hand-side of the $`g_{\mu \nu }`$\- and $`\phi `$-field equations, i.e. the quantities $`t^2T_{(A)0}^0`$ and $`t^2T_{(A)j}^i`$ where $`T_{(A)\nu }^\mu `$ denotes the stress-energy tensor of the $`p`$-form. \[As usual the mixed terms $`T_{(A)i}^0`$, which enter the momentum constraints play a rather different role and do not need to be explicitly considered.\] Finally, one gets as “dangerous” terms when $`t0`$ a sum of “electric” contributions $`t^{2e_{i_1\mathrm{}i_p}^{(p)}}`$ and of “magnetic” ones $`t^{2b_{j_1\mathrm{}j_{dp1}}^{(p)}}`$. Here, the electric exponents $`e_{i_1\mathrm{}i_p}^{(p)}`$ (where all the indices $`i_n`$ are different) are defined as
$$e_{i_1\mathrm{}i_p}^{(p)}(p)=p_{i_1}+p_{i_2}+\mathrm{}+p_{i_p}\frac{1}{2}\lambda _pp_\phi ,$$
(12)
while the magnetic exponents $`b_{j_1\mathrm{}j_{dp1}}^{(p)}`$ (where all the indices $`j_n`$ are different) are
$$b_{j_1\mathrm{}j_{dp1}}^{(p)}(p)=p_{j_1}+p_{j_2}+\mathrm{}+p_{j_{dp1}}+\frac{1}{2}\lambda _pp_\phi .$$
(13)
To each $`p`$-form is associated a (duality-invariant) double family of “stability” exponents $`e^{(p)}`$, $`b^{(p)}`$. The “electric” (respectively “magnetic”) stability condition is that all the exponents $`e^{(p)}`$ (respectively, $`b^{(p)}`$) be positive. This result generalizes the results of on the effect of vector fields in $`D=4`$.
The main result reported here is that, for all superstring models, there exists no open region of the Kasner sphere where all the stability exponents $`g(p)`$, $`e(p)`$, $`b(p)`$ are strictly positive. To define the set of stability conditions for the various superstring models, let us review their field content and give the values of the crucial dilaton couplings $`\lambda _p`$. The simplest case is the massless bosonic sector of $`M`$-theory, i.e of SUGRA in $`D=11`$. In that case, there is a 3-form and no dilaton. The parameters $`p_\alpha ^M`$, $`\alpha =1,\mathrm{},10`$, run over the 8-dimensional sphere $`S_M^8`$ defined by $`_\alpha (p_\alpha ^M)^2=1=_\alpha p_\alpha ^M`$. The presence of a 3-form $`A_3`$ uncoupled to any dilaton implies that the electric and magnetic stability exponents are respectively given by (12) and (13) with $`p=3`$, $`\lambda _p=0`$ and $`d=10`$, i.e., $`e_{\alpha _1\alpha _2\alpha _3}^{M(3)}=p_{\alpha _1}^M+p_{\alpha _2}^M+p_{\alpha _3}^M`$ and $`b_{\alpha _1\mathrm{}\alpha _6}^{M(3)}=p_{\alpha _1}^M+\mathrm{}+p_{\alpha _6}^M`$.
The $`D=10`$ type IIA string theory involves, besides $`g_{\mu \nu }`$ and a dilaton $`\phi =\mathrm{\Phi }/\sqrt{2}`$ (with $`g_s=e^\mathrm{\Phi }`$ being the string coupling) a 1-form, a 2-form and a 3-form. The (Einstein-frame) dilaton coupling parameters of the forms are $`\lambda _1^A=3\sqrt{2}/2`$, $`\lambda _2^A=\sqrt{2}`$ and $`\lambda _3^A=\sqrt{2}/2`$, respectively. Besides the dilaton Kasner exponent $`p_\phi ^A`$, there are nine metric exponents $`p_i^A`$, $`i=1,\mathrm{},9`$. They run over $`S_A^8`$ defined by Eqs. (Chaos in Superstring Cosmology).
The $`D=10`$ type IIB string theory involves (besides $`g_{\mu \nu }`$): two scalars: the dilaton $`\phi =\mathrm{\Phi }/\sqrt{2}`$ and the $`RR`$ 0-form $`C_0`$, two 2-forms $`B_2(NSNS)`$ and $`C_2(RR)`$, and one “self-dual” $`RR`$ 4-form $`C_4`$. The dilaton coupling strengths of the forms are: $`\lambda _{C_0}^B=2\sqrt{2}`$, $`\lambda _{B_2}^B=\sqrt{2}`$, $`\lambda _{C_2}^B=+\sqrt{2}`$ and $`\lambda _{C_4}^B=0`$. \[$`\lambda _{C_2}^B`$ refers to the more complicated mixed coupling $`e^\mathrm{\Phi }(dC_2C_0dB_2)^2`$\]. The Kasner exponents $`p_\phi ^B`$, $`p_i^B`$ ($`i=1,\mathrm{},9`$) run over $`S_B^8`$ defined by Eqs. (Chaos in Superstring Cosmology).
The $`D=10`$ type I string theory involves (besides $`g_{\mu \nu }`$ and $`\phi `$): an $`\mathrm{SO}(32)`$ vector potential, and a 2-form. The dilaton couplings are $`\lambda _1^I=\sqrt{2}/2`$ and $`\lambda _2^I=+\sqrt{2}`$<sup>*</sup><sup>*</sup>*Note the misprint in Eq. (12.1.34b) of , corrected on the author’s web page.. The Kasner exponents $`p_\phi ^I`$, $`p_i^I`$ ($`i=1,\mathrm{},9`$) run, as for IIA and IIB, on the $`S_I^8`$ defined by Eqs. (Chaos in Superstring Cosmology).
Finally, the $`D=10`$ heterotic string theories involve (besides $`g_{\mu \nu }`$ and $`\phi `$): an $`\mathrm{SO}(32)`$ or $`E_8\times E_8`$ vector potential, and a 2-form. Their respective (Einstein-frame) dilaton couplings are: $`\lambda _1^h=\sqrt{2}/2`$, $`\lambda _2^h=\sqrt{2}`$. The Kasner sphere $`S_h^8`$ for $`p_\phi ^h`$, $`p_i^h`$ ($`i=1,\mathrm{},9`$) is the same as for IIA, IIB or I.
Let us denote for each given theory “th” (where $`\mathrm{th}=M,A,B,I,h`$ labels the theory) the full (finite) sequence of stability exponents as $`w_J^{\mathrm{th}}(p)`$, where $`J`$ labels all the possible exponents within each theory. E.g., when $`\mathrm{th}=M`$ the label $`J`$ takes 690 values corresponding to the set $`\{w_J^M\}=\{g_{\alpha \beta \gamma }^M,e_{\alpha _1\alpha _2\alpha _3}^{M(3)},b_{\beta _1\mathrm{}\beta _6}^{M(3)}\}`$. The condition of “Kasner stability” of each theory is that there exist an open region of the corresponding Kasner sphere $`S_{\mathrm{th}}^8`$ where $`w_J^{\mathrm{th}}(p)>0`$ for all the labels $`J`$. However, we have proven that, for all theories, $`\mathrm{inf}_Jw_J^{\mathrm{th}}(p)`$ is strictly negative for all values of $`pS_{\mathrm{th}}^8`$, except at a finite number of isolated points where it vanishes.
Let us first consider $`M`$-theory. We have proven a stronger result, namely, that the electric stability conditions alone are never fulfilled. If, at any point on $`S_M^8`$, we order the Kasner exponents as $`p_1^Mp_2^M\mathrm{}p_{10}^M`$, the most stringent electric stability criterion involves $`f_0(p)p_1^M+p_2^M+p_3^M`$. To show that this function is non-positive on the cell $`p_1^M\mathrm{}p_{10}^M`$ of the Kasner sphere, we maximize it subject to the constraints $`(p_\alpha ^M)^2=1`$, $`p_\alpha ^M=1`$. These constraints can be taken into account by introducing two Lagrange multipliers. After a straightforward (but rather long) exhaustive analysis, we have found that $`f_0^{\mathrm{max}}=0`$, this maximum being reached only at $`p_1=\mathrm{}=p_9=0`$, $`p_{10}=1`$.
To deal with the type IIA theory, we use the fact that IIA is the Kaluza-Klein (KK) reduction of $`M`$ on a circle. This fact dictates the link between the field variables of the two models. If we label by the letter $`y`$ the compactified dimension this link implies the following relation between the (Einstein-frame) Kasner exponents of the two theories ($`i=1,\mathrm{},9`$)
$$p_\phi ^A=\frac{6\sqrt{2}p_y^M}{8+p_y^M},p_i^A=\frac{8p_i^M+p_y^M}{8+p_y^M}.$$
(14)
Forgetting about this Kaluza-Klein motivationNote that the fact that the IIA field variables depend on one less variable than the $`M`$-ones is unimportant. What is important is the map (14) and the fact that we have taken into account in the stability criteria all possible dangerous terms in a generic solution., we can consider that Eqs. (14) define a one-to-one map $`\pi _{AM}`$ from $`S_M^8`$ to $`S_A^8`$: $`p_\alpha ^A=\pi _{AM}(p_\beta ^M)`$. Using this map, we have then shown that the complete set of IIA stability conditions is logically equivalent to the complete set of $`M`$ stability conditions. The instability of the Kasner behaviour of $`M`$-theory proven above then implies that the Kasner behaviour of IIA is also unstable.
To deal with the type IIB theory, we use the fact that IIA and IIB are related by $`T`$-duality. The link between the field variables of the two models dictated by $`T`$-duality enables one to derive a certain fractionally linear map $`\pi _{BA}`$ between their (Einstein-frame) Kasner exponents, which can be used, as above, to prove the Kasner-stability equivalence of the types IIA and IIB theories. Since type IIA is unstable, type IIB is also unstable.
At this stage, we know that $`M`$, IIA and IIB are equivalent with respect to Kasner stability, and are all unstable. It remains to tackle the type I and heterotic theories, which are equivalent because their stability conditions are algebraically mapped onto each other by the $`S`$-duality transformation $`p_\phi ^I=p_\phi ^h,p_i^I=p_i^h`$. To study the Kasner-stability of the heterotic theory, we found very convenient to replace the Einstein-frame Kasner exponents $`(p_\phi ^h,p_i^h)`$ by their string-frame counterparts $`(\alpha _i^h)`$. The link between the two is (in $`d+1`$ spacetime dimensions, see, e.g. )
$$p_\phi =\frac{\sqrt{d1}\sigma }{d1\sigma },p_i=\frac{(d1)\alpha _i\sigma }{d1\sigma },$$
(15)
with $`\sigma \left(_i\alpha _i\right)1`$ and $`i=1,\mathrm{},d`$. In terms of the $`\alpha `$’s the Kasner sphere $`S^{d1}`$ is simply the usual unit sphere, $`_i(\alpha _i)^2=1`$. In our case, $`d=9`$ and one should add a label “$`h`$” to both the $`p`$’s and the $`\alpha `$’s. In terms of the string-frame exponents it is found that the $`h`$-stability conditions are equivalent to the simpler inequalities $`\alpha _i^h>0`$ and $`\alpha _i^h+\alpha _j^h+\alpha _k^h<1`$ (where $`i,j,k`$ are all different) subjected to the constraints $`_i(\alpha _i^h)^2=1`$. It is easy to verify that these inequalities can never hold when the space dimension is $`d=9`$. In that case, the closest one comes to satisfying the inequalities is the isotropic point $`\alpha _i=1/3`$ for which the second inequality is saturated. This concludes our proof that the heterotic model (and therefore also the type I one) is Kasner unstable. Finally the two blocks of theories $`(M,A,B)`$ and $`(I,h)`$ are both Kasner unstable, though for different algebraic reasons.
Our results so far show that the generic solution of the low-energy string models can never reach a monotonic Kasner-like behaviour. Following the BKL approach one can go further and study the evolution near a cosmological singularity as a sequence of Kasner-like “free flights” interrupted by “collisions” against the “potential walls” corresponding to the various stability-violating exponents $`g`$, $`e`$ or $`b`$. We have studied this problem and found the following universal “collision law” giving the Kasner exponents $`\overline{p}^\mu `$ of the Kasner epoch following a collision in terms of the old ones:
$$\overline{p}^\mu =\left(12\frac{(w\overline{p})(wu)}{(ww)}\right)^1\left[\overline{p}^\mu 2\frac{(w\overline{p})w^\mu }{(ww)}\right].$$
(16)
Here, $`\overline{p}^\mu `$ stands for $`\overline{p}^0p_\phi `$ and $`\overline{p}^ip_i`$. The scalar products are computed with the metric $`G_{\mu \nu }`$ occurring in the quadratic form (8), namely, $`G_{00}=1`$, $`G_{0i}=0`$, $`G_{ij}=\delta _{ij}1`$, while the vector $`u`$ (entering (10)) has “covariant” components $`u_0=0`$, $`u_i=1`$. Finally, the “contravariant” vector $`w^\mu `$ characterizes the “wall” responsible for the collision and is defined in such a way that the corresponding exponent ($`g(p)`$, $`e(p)`$ or $`b(p)`$) reads $`w(p)=w_\mu \overline{p}^\mu G_{\mu \nu }w^\mu \overline{p}^\nu `$. \[E.g., for the wall associated with the electric exponent $`e_{123}(p)p_1+p_2+p_3`$, $`w_\mu `$ reads $`w_0=0`$, $`w_i=1`$ for $`i=1,2,3`$ and $`w_i=0`$ for $`i>3`$.\] The result (16) (which is a rescaled geometrical reflection in the hyperplane $`w_\mu \overline{p}^\mu =0`$) applies uniformly to all possible “walls”: gravitational, electric or magnetic. It generalizes particular results derived by many authors .
Summarizing: In all string models, the generic solution near a cosmological singularity for the massless bosonic degrees of freedom exhibits BKL-type oscillations, i.e. a (formally infinite) alternation of Kasner-epochs. The primary sources of this BKL behaviour are (i) the presence of $`p`$-forms in the field spectrum of the theories and, (ii) the strength of their dilaton couplings. In the absence of $`p`$-forms, or if the $`\lambda _p`$’s were somewhat smaller, the monotonic Kasner behaviour would be stable and generic. The general rule defining the change of Kasner exponents from one epoch to the next is given by Eq. (16), where $`w`$ is the “wall” (among the various gravitational, electric or magnetic ones) for which $`w(p)=w_\mu \overline{p}^\mu `$ is most negative. We anticipate that the discrete dynamics (16) will define (in all string models) a chaotic motion on the Kasner sphere. At this stage, the physical consequences of such a chaotic motion are unclear. It might constitute a problem for the pre-big-bang scenario which strongly relies on the existence, near a (future) cosmological singularity, of relatively large, quasi-uniform patches of space following a monotonic, dilaton-driven Kasner behaviour. By contrast our findings suggest that the spatial inhomogeneity continuously increases toward a singularity, as all quasi-uniform patches of space get broken up into smaller and smaller ones by the chaotic oscillatory evolution. In other words, the spacetime structure tends to develop a kind of “turbulence” .
We are aware of the limitations of our result (tree-level bosonic massless modes only) but we think that our finding suggests that the full quantum, string-theory behaviour might be at least as complicated, near a cosmological singularity, as our simplified analysis shows.
We thank Volodia Belinskii, Isaak Khalatnikov and Ilan Vardi for useful exchanges of ideas. T.D. is grateful to David Gross, Gary Horowitz and Joe Polchinski for informative discussions. M. H. thanks the Institut des Hautes Etudes Scientifiques for its kind hospitality. |
warning/0003/hep-th0003165.html | ar5iv | text | # 1 Introduction
## 1 Introduction
D-branes play important roles to describe the solitonic modes in string theory and could make clear dynamics in strong coupling regions. The physical observables of D-brane’s effective theories have dependences on moduli of compactified strings or wrapped D-branes. We expect that properties of compactified internal spaces essentially control non-perturbative effects in low energy theories. In this paper, we focus on the type II superstring compactified on Calabi-Yau manifold and study its topological sector from the point of view of topological sigma models (A- and B-models) to investigate properties of moduli spaces. Because both the A- and B-models are topological theories, they are characterized by their two-point and three point functions that play important roles as the constituent blocks in these models-. Topological metrics are two-point functions and receive no quantum corrections. On the other hand, the three-point functions of the A-model have information about the fusion structure of observables. The remaining fundamental blocks are a Kähler potential $`K`$ and associated hermitian two-point functions. They are hermitian and describe correlations of topological ant anti-topological sectors,,,. Also the $`K`$ has information about intersections of homology cycles with even dimensions in the A-model case.
The aim of this paper is to develop a concrete method to construct Kähler potentials applicable in the large radius regions of Calabi-Yau $`d`$-folds and to investigate their properties in order to understand structures of the moduli spaces. We present formulae of Kähler potentials for $`d`$ dimensional Calabi-Yau manifolds explicitly.
The paper is organized as follows. In section 2, we explain a mirror manifold paired with a Calabi-Yau $`d`$-fold embedded in $`CP^{d+1}`$. We also explain the results in about a Kähler potential $`K`$ in the small complex structure region of the B-model in order to fix notations. In section 3, we introduce a set of periods valid in the large complex structure region. By relating two sets of periods in the large and small complex structure regions, we construct a formula of the $`K`$ applicable in the large complex structure region of the B-model. In sections 4 and 5, we construct a mirror map and a Kähler potential. The scalar curvature of the Kähler moduli is investigated. The set of correlation functions associated with the Kähler moduli are calculated in the large radius region of the A-model. A concrete application of our result is explained in the quintic case in section 6. Also the result there is generalized to propose a formula of the Kähler potential of a Calabi-Yau $`d`$-fold in the complete intersection type in section 7. Section 8 is devoted to conclusions and comments. In appendix A, we summarize several examples of the expansion coefficients of a function $`\widehat{K}`$ in lower dimensional cases.
## 2 Small Complex Structure Region
We take a one-parameter family of Calabi-Yau $`d`$-fold $`M`$ realized as a zero locus of a hypersurface embedded in a $`CP^{d+1}`$
$`M;p=X_1^N+X_2^N+\mathrm{}+X_N^NN\psi X_1X_2\mathrm{}X_N=0.`$
The $`N`$ is related with the complex dimension $`d`$ of M, $`N=d+2`$. A mirror manifold $`W`$ paired with this $`M`$ is constructed as a orbifold divided by some maximally discrete group $`G=𝐙_N^{(N1)}`$
$`W;\widehat{\{p=0\}/G}.`$
When one thinks about Hodge structure of the $`G`$-invariant parts of the cohomology group $`H^d(W)`$, related Hodge numbers are written as
$`h^{d,0}=h^{d1,1}=\mathrm{}=h^{1,d1}=h^{0,d}=1.`$
In our previous paper, we study the formula of the Kähler potential of the Calabi-Yau $`d`$-fold $`W`$ with a moduli parameter $`\psi `$ of the complex structure.
The $`K`$ is constructed by combining a set of periods $`\stackrel{~}{\varpi }_k`$ quadratically
$`e^K={\displaystyle \underset{k=1}{\overset{N1}{}}}I_k\stackrel{~}{\varpi }_k^{}\stackrel{~}{\varpi }_k,`$ (1)
$`I_k={\displaystyle \frac{1}{\pi ^NN^{N+2}}}(1)^{k1}\left(\mathrm{sin}{\displaystyle \frac{\pi k}{N}}\right)^N,`$
$`\stackrel{~}{\varpi }_k(\psi )=\left[\mathrm{\Gamma }\left({\displaystyle \frac{k}{N}}\right)\right]^N{\displaystyle \frac{(N\psi )^k}{\mathrm{\Gamma }(k)}}`$
$`\times \left[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{k}{N}+n\right)}{\mathrm{\Gamma }\left(\frac{k}{N}\right)}}\right]^N{\displaystyle \frac{\mathrm{\Gamma }(k)}{\mathrm{\Gamma }(Nn+k)}}(N\psi )^{Nn}\right].`$
We determined the coefficients $`I_k`$ in the by requiring consistency conditions with the results of the CFT at the Gepner point. This formula is valid in the small $`\psi `$ region because of the convergence of the series expansion.
In the following section, we investigate the large complex structure region of the W. Its formula is important for the large radius analyses of the manifold $`M`$.
## 3 Large Complex Structure Region
We try to rewrite the $`K`$ in a formula that is valid in the large $`\psi `$ by an analytic continuation. First we choose a set of periods $`\{\mathrm{\Omega }_m\}`$ ($`m=0,1,\mathrm{},N2`$) appropriate to describe the large complex structure region of the mirror W. A generating function of the $`\mathrm{\Omega }_m`$ is defined by using a formal parameter $`\rho `$ with $`\rho ^{N1}=0`$
$`{\displaystyle \underset{m=0}{\overset{N2}{}}}\mathrm{\Omega }_m\rho ^m=\sqrt{\widehat{K}(\rho )}\varpi ({\displaystyle \frac{\rho }{2\pi i}};z),`$
$`\varpi (v)=z^v{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a(n+v)}{a(v)}}z^n,z=(N\psi )^N,`$ (2)
$`a(v)={\displaystyle \frac{\mathrm{\Gamma }(Nv+1)}{[\mathrm{\Gamma }(v+1)]^N}}.`$
Here we introduce a function $`\widehat{K}(\rho )`$.
$`\widehat{K}(\rho ):={\displaystyle \frac{a\left(+{\displaystyle \frac{\rho }{2\pi i}}\right)}{a\left({\displaystyle \frac{\rho }{2\pi i}}\right)}}=\mathrm{exp}\left[2{\displaystyle \underset{m=1}{}}{\displaystyle \frac{NN^{2m+1}}{2m+1}}\zeta (2m+1)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^{2m+1}\right]`$
$`=1+2\zeta (3){\displaystyle \frac{c_3}{N}}\left({\displaystyle \frac{\rho }{2\pi i}}\right)^3+𝒪(\rho ^5).`$
The leading term in the $`\widehat{K}`$ of the variety W is a 3rd Chern class of the M. Generally the coefficients of $`\widehat{K}`$ contain topological information of M. In fact, Chern classes of $`c_{\mathrm{}}`$ $`(\mathrm{}=1,2,\mathrm{},N2)`$ of the manifold $`M`$ are generated by a function $`c(\rho )`$
$`c(\rho )={\displaystyle \frac{(1+\rho )^N}{1+N\rho }}=1+{\displaystyle \underset{\mathrm{}1}{}}\rho ^{\mathrm{}}{\displaystyle \frac{c_{\mathrm{}}}{N}}.`$
Typical coefficients $`X_{\mathrm{}}=NN^{\mathrm{}}`$ in the $`\widehat{K}`$ are some combinations of Chern classes $`c_{\mathrm{}}`$
$`c(\rho )=1+{\displaystyle \underset{\mathrm{}1}{}}\rho ^{\mathrm{}}{\displaystyle \frac{c_{\mathrm{}}}{N}}=\mathrm{exp}\left({\displaystyle \underset{\mathrm{}1}{}}(1)^\mathrm{}1\rho ^{\mathrm{}}{\displaystyle \frac{X_{\mathrm{}}}{\mathrm{}}}\right).`$
For examples, we list several $`X_{\mathrm{}}`$ for the Calabi-Yau case
$`X_1=0,X_2={\displaystyle \frac{2}{N}}c_2,X_3={\displaystyle \frac{3}{N}}c_3,`$
$`X_4={\displaystyle \frac{4}{N}}c_4+{\displaystyle \frac{2}{N^2}}c_2^2,X_5={\displaystyle \frac{5}{N}}c_5{\displaystyle \frac{5}{N^2}}c_3c_2.`$
The series expansion Eq.(2) converges around $`z0`$, that is, large complex structure point of the $`W`$. For the purpose of an analytic continuation into the large complex structure region, we find that the two sets of the periods are related by a transformation matrix $`\stackrel{~}{M}`$ with components $`\stackrel{~}{M}_k\mathrm{}`$
$`\stackrel{~}{\varpi }_k={\displaystyle \underset{\mathrm{}=0}{\overset{N2}{}}}\stackrel{~}{M}_k\mathrm{}\mathrm{\Omega }_{\mathrm{}}(k=1,2,\mathrm{},N1),`$
$`\stackrel{~}{M}_k\mathrm{}=(N)(2\pi i)^{N1}\times \left[\sqrt{\widehat{A}(\rho )}{\displaystyle \frac{\alpha ^k}{e^\rho \alpha ^k}}(\rho )^{\mathrm{}}\right]|_{\rho ^{N2}}`$
$`=(N)(2\pi i)^{N1}{\displaystyle \underset{m=0}{\overset{N2}{}}}G_{k,m}V_{m,\mathrm{}},`$
$`G_{k,m}={\displaystyle \frac{\alpha ^k}{(\alpha ^k1)^{m+1}}}(1kN1,\mathrm{\hspace{0.17em}\hspace{0.17em}0}mN2),`$
$`\alpha =e^{\frac{2\pi i}{N}},`$
$`V_{m,\mathrm{}}=[\sqrt{\widehat{A}(\rho )}(e^\rho 1)^m(\rho )^{\mathrm{}}]|_{\rho ^{N2}}(0mN2,\mathrm{\hspace{0.17em}\hspace{0.17em}0}\mathrm{}N2).`$
Here the transformation matrix $`V`$ contains a square root of a topological invariant “A-roof” of the Calabi-Yau space $`M`$
$`\widehat{A}(\rho )=\left({\displaystyle \frac{{\displaystyle \frac{\rho }{2}}}{\mathrm{sinh}{\displaystyle \frac{\rho }{2}}}}\right)^N\left({\displaystyle \frac{\mathrm{sinh}{\displaystyle \frac{N\rho }{2}}}{{\displaystyle \frac{N\rho }{2}}}}\right)={\displaystyle \frac{1}{a({\displaystyle \frac{\rho }{2\pi i}})a(+{\displaystyle \frac{\rho }{2\pi i}})}}`$
$`=\mathrm{exp}\left[+{\displaystyle \underset{m=1}{}}{\displaystyle \frac{(1)^mB_m}{(2m)!}}{\displaystyle \frac{NN^{2m}}{2m}}\rho ^{2m}\right]`$
$`=1+{\displaystyle \frac{1}{12}}{\displaystyle \frac{c_2}{N}}\rho ^2+𝒪(\rho ^4).`$
The $`B_m`$s are Bernoulli numbers and are defined in our convention as
$`{\displaystyle \frac{x}{e^x1}}=1{\displaystyle \frac{x}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^nB_n}{(2n)!}}x^{2n},`$
$`B_1={\displaystyle \frac{1}{6}},B_2={\displaystyle \frac{1}{30}},B_3={\displaystyle \frac{1}{42}},B_4={\displaystyle \frac{1}{30}},\mathrm{}.`$
Now we return to the Kähler potential $`K`$. By performing the analytic continuation of the Eq.(1), we can obtain a formula of the $`K`$ applicable in the large $`\psi `$ region.
$`e^K=(1)^N\left({\displaystyle \frac{2\pi i}{N}}\right)^{N2}{\displaystyle \frac{1}{N^2}}{\displaystyle \underset{\mathrm{},\mathrm{}^{}=0}{\overset{N2}{}}}\left(V^{}V\right)_\mathrm{},\mathrm{}^{}(1)^\mathrm{}+\mathrm{}^{}\overline{\mathrm{\Omega }}_{\mathrm{}}\mathrm{\Omega }_{\mathrm{}^{}},`$
Here the $``$ is a triangular matrix that combines $`\overline{\mathrm{\Omega }}`$ and $`\mathrm{\Omega }`$ quadratically
$`_{m,m^{}}=2^Ni^N{\displaystyle \underset{k=1}{\overset{N1}{}}}{\displaystyle \frac{(1)^k\left(\mathrm{sin}\frac{\pi k}{N}\right)^N}{(\alpha ^k1)^{m+1}(\alpha ^k1)^{m^{}+1}}}=(1)^m\delta _{m+m^{},N2}+\mathrm{}.`$
The $`_{m,m^{}}`$ has non-vanishing components only at $`m+m^{}N2`$ ($`m,m^{}=0,1,\mathrm{},N2`$) and we find an expression for the matrix $`V^{}V`$
$`\left(V^{}V\right)_\mathrm{},\mathrm{}^{}=(1)^{\mathrm{}^{}}\delta _{\mathrm{}+\mathrm{}^{},N2}.`$ (3)
We check validity of this equation Eq.(3) concretely for $`N27`$ and propose this formula for arbitrary $`N(3)`$ cases as a conjecture.
Finally we obtain a formula of the $`K`$ of $`W`$ with this equation
$`e^K=(1)^d\left({\displaystyle \frac{2\pi i}{N}}\right)^d{\displaystyle \frac{1}{N^2}}\left(\mathrm{\Omega }^{}\mathrm{\Sigma }\mathrm{\Omega }\right),`$
$`\mathrm{\Sigma }_\mathrm{},\mathrm{}^{}=(1)^{\mathrm{}}\delta _{\mathrm{}+\mathrm{}^{},N2}.`$
The $`e^K`$ is constructed by combining a holomorphic $`d`$ form and an anti-holomorphic one quadratically. Both parts are decomposed by a dual basis of (real) homology cycles and their coefficients are realized as periods. Then we can understand that a matrix which combines the periods and their complex conjugates is an intersection matrix of the cycles. In our case, the $`\mathrm{\Sigma }`$ is an intersection matrix of homology cycles associated with the set of periods $`\mathrm{\Omega }_m`$ of $`W`$. The result means that the cycles we used here are combined into a symplectic $`\text{USp}(d+1)`$ or an $`\text{SO}(\frac{d}{2}+1,\frac{d}{2})`$ invariant bases for respectively $`d=`$odd or $`d=`$even cases. But the basis is not an integral one and we have to perform an appropriate linear transformation with fractional rational numbers to construct a canonical basis of a central charge. More details will appear in our next paper.
Let us study behaviors of a metric, a curvature in the large $`\psi `$ region of $`W`$. Powers of logarithm of $`\psi `$ appear in the Kähler potential $`K`$ in the leading expansion
$`e^K={\displaystyle \frac{1}{d!N^2}}\{2\mathrm{log}(N|\psi |)\}^d\times \left[1+{\displaystyle \underset{n=1}{\overset{d}{}}}{\displaystyle \frac{d!}{(dn)!}}{\displaystyle \frac{(1)^n\widehat{K}_n}{\{2N\mathrm{log}(N|\psi |)\}^n}}\right]+\mathrm{}`$
$`={\displaystyle \frac{1}{d!N^2}}\{2\mathrm{log}(N|\psi |)\}^d\times [1+4\left(\begin{array}{c}d\\ 3\end{array}\right)(N^3N)\zeta (3)\{2N\mathrm{log}(N|\psi |)\}^3`$
$`+48\left(\begin{array}{c}d\\ 5\end{array}\right)(N^5N)\zeta (5)\{2N\mathrm{log}(N|\psi |)\}^5+\mathrm{}].`$
Here the $`\widehat{K}_n`$s are coefficients of the series expansion of the $`\widehat{K}`$
$`\widehat{K}(\rho )=1+{\displaystyle \underset{n=3}{}}\widehat{K}_n\left({\displaystyle \frac{\rho }{2\pi i}}\right)^n,`$
$`\widehat{K}_3={\displaystyle \frac{2}{3}}(N^3N)\zeta (3)={\displaystyle \frac{2}{N}}c_3\zeta (3),`$
$`\widehat{K}_5={\displaystyle \frac{2}{5}}(N^5N)\zeta (5)={\displaystyle \frac{2}{N}}(c_5{\displaystyle \frac{1}{N}}c_3c_2)\zeta (5),`$
They are related to the Chern classes of the $`d`$-fold $`M`$. Also there appear $`\zeta (2m+1)`$s ($`m=1,2,\mathrm{}`$ with $`2m+1d`$) , which might be transcendental numbers, in this formula. We summarize several concrete examples of the $`\widehat{K}_n`$ in the appendix. The exponent of power of the logarithm is at most $`d`$ for the $`d`$-fold. In addition, there are parts of infinite series with respect to the $`\psi `$ in the $`e^K`$. They are omitted as an abbreviated symbol “$`\mathrm{}`$”.
Next we calculate the Kähler metric of this B-model associated with $`W`$
$`g_{\psi \overline{\psi }}={\displaystyle \frac{d}{\{2|\psi |\mathrm{log}(N|\psi |)\}^2}}\times [116\left(\begin{array}{c}d1\\ 2\end{array}\right)(N^3N)\zeta (3)\{2N\mathrm{log}(N|\psi |)\}^3`$
$`288\left(\begin{array}{c}d1\\ 4\end{array}\right)(N^5N)\zeta (5)\{2N\mathrm{log}(N|\psi |)\}^5+\mathrm{}].`$
Also we obtain the scalar curvature in this large $`\psi `$ region
$`R={\displaystyle \frac{4}{d}}\times [180\left(\begin{array}{c}d1\\ 2\end{array}\right)(N^3N)\zeta (3)\{2N\mathrm{log}(N|\psi |)\}^3`$
$`4032\left(\begin{array}{c}d1\\ 4\end{array}\right)(N^5N)\zeta (5)\{2N\mathrm{log}(N|\psi |)\}^5+\mathrm{}].`$
In the $`|\psi |=\mathrm{}`$ limit, the curvature is a negative constant and its absolute value is inversely proportional to the dimension $`d`$. The $`|R|`$ decreases apart from the point $`|\psi |=\mathrm{}`$ for $`N5`$ cases. The leading term of the corrections in the brackets is inversely proportional to the $`\{\mathrm{log}(N\psi )\}^3`$ and contains $`\zeta (3)`$ as its coefficient. For the $`N=3,4`$ cases, the associated curvatures $`R`$s are constants except for points in the $`\psi `$-plane with $`\psi ^N=1`$.
Next we discuss an invariant coupling. The Kähler potential is not a function but a section of a line bundle and there is an arbitrariness of multiplications by holomorphic and anti-holomorphic functions. Also a $`d`$-point correlation function in the B-model is a section with a weight $`2`$
$`\kappa _{\underset{d\text{times}}{\underset{}{\psi \psi \mathrm{}\psi }}}={\displaystyle \frac{1}{N^d}}{\displaystyle \frac{N\psi ^2}{1\psi ^N}}.`$
But there is an invariant $`d`$-point function $`\kappa `$ that is constructed by combining the metric, the $`K`$ and $`\kappa _{\psi \mathrm{}\psi }`$
$`\kappa =(g_{\psi \overline{\psi }})^{d/2}e^K|\kappa _{\psi \psi \mathrm{}\psi }|.`$
In our normalization, it is expressed in the large $`|\psi |`$ limit as
$`\kappa ={\displaystyle \frac{d!}{d^{d/2}N^{d3}}}\times [1+20\left(\begin{array}{c}d\\ 3\end{array}\right)(N^3N)\zeta (3)\{2N\mathrm{log}(N|\psi |)\}^3`$
$`+672\left(\begin{array}{c}d\\ 5\end{array}\right)(N^5N)\zeta (5)\{2N\mathrm{log}(N|\psi |)\}^5+\mathrm{}].`$
Up to coefficients, the corrections have the same structures for the metric, the $`R`$ and the $`\kappa `$. In the small complex structure limit (at $`\psi =0`$), this invariant coupling is evaluated as
$`\kappa =\left[{\displaystyle \frac{\mathrm{\Gamma }\left({\displaystyle \frac{1}{N}}\right)}{\mathrm{\Gamma }\left(1{\displaystyle \frac{1}{N}}\right)}}\right]^{\frac{1}{2}N(d2)}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(1{\displaystyle \frac{2}{N}}\right)}{\mathrm{\Gamma }\left({\displaystyle \frac{2}{N}}\right)}}\right]^{\frac{1}{2}Nd}\times {\displaystyle \frac{1}{N^{d3}}}.`$
For an example, we write down the result in the $`d=3`$ case
$`\kappa ={\displaystyle \frac{2}{\sqrt{3}}}\left[112c_3\zeta (3)\{10\mathrm{log}(5|\psi |)\}^3+\mathrm{}\right],c_3=200.`$
This shows that the coupling increases apart from the $`|\psi |=\mathrm{}`$.
Also we compare the leading value of this $`\kappa `$ in the small $`|\psi |`$ with that in the large $`|\psi |`$ limits
$`\text{small }\psi ;\kappa =\left({\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{5})}{\mathrm{\Gamma }(\frac{4}{5})}}\right)^{5/2}\left({\displaystyle \frac{\mathrm{\Gamma }(\frac{3}{5})}{\mathrm{\Gamma }(\frac{2}{5})}}\right)^{15/2}`$
$`=1.55531898996323897249994495854237822\mathrm{}`$
$`\text{large }\psi ;\kappa ={\displaystyle \frac{2}{\sqrt{3}}}`$
$`=1.154700538379251529018297561003914911\mathrm{}.`$
The $`\kappa `$ in the small $`\psi `$ case can be analyzed by using our previous result in . The coupling in the small $`\psi `$ region is stronger than that in the large $`\psi `$ region in an amount of 34.7%.
## 4 Large Radius Region
In the previous section, we study properties of physical quantities in the large complex structure point of the $`d`$-fold $`W`$. It is known that the large complex structure point is related to the large radius point of the partner $`M`$. It is possible to translate the results of $`W`$ to those of the $`M`$. For the purpose of this program, we introduce a mirror map $`t`$
$`2\pi it=\mathrm{log}z+{\displaystyle \frac{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(Nn)!}{(n!)^N}}\left({\displaystyle \underset{m=n+1}{\overset{Nn}{}}}{\displaystyle \frac{N}{m}}\right)z^n}{{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(Nn)!}{(n!)^N}}z^n}},q=e^{2\pi it},`$
$`x_n={\displaystyle \frac{1}{(2\pi i)^n}}{\displaystyle \frac{1}{n!}}_\rho ^n\mathrm{log}\left[{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(N(m+\rho )+1)}{\mathrm{\Gamma }(N\rho +1)}}\left({\displaystyle \frac{\mathrm{\Gamma }(\rho +1)}{\mathrm{\Gamma }(m+\rho +1)}}\right)^Nz^m\right]|_{\rho =0}`$
$`={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}a_{n,m}q^m(n2),`$
$`a_{n,1}=N!S_n(\beta _1(1),\mathrm{},\beta _n(1)),`$
$`a_{n,2}={\displaystyle \frac{(2N)!}{2^N}}S_n(\beta _1(2),\mathrm{},\beta _n(2))`$
$`(N!)^2\left({\displaystyle \underset{m=2}{\overset{N}{}}}{\displaystyle \frac{N}{m}}\right)S_n(\beta _1(1),\mathrm{},\beta _n(1))`$
$`{\displaystyle \frac{1}{2}}(N!)^2S_n(2\beta _1(1),\mathrm{},2\beta _n(1)),`$
$`\beta _m(n)=N^m{\displaystyle \underset{k=1}{\overset{Nn}{}}}{\displaystyle \frac{(1)^{k1}(k1)!}{k^m}}N{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{(1)^{k1}(k1)!}{k^m}}.`$
This $`t`$ is a coordinate of the Kähler moduli space or the coefficient of the complexified Kähler form
$`B+iJ=t[D],[D]\text{H}^2(M).`$
The $`[D]`$ is a Poincaré dual of a divisor “$`D`$” of the $`M`$. The $`S_n(x_1,x_2,\mathrm{},x_n)`$s are Schur polynomials and are defined as
$`\mathrm{exp}\left({\displaystyle \underset{m=1}{}}x_mu^m\right)={\displaystyle \underset{n=0}{}}S_n(x_1,x_2,\mathrm{},x_n)u^n.`$
Next we rewrite the $`K`$ in this coordinate
$`e^K=(1)^d\left({\displaystyle \frac{2\pi i}{N}}\right)^d{\displaystyle \frac{1}{N^2}}(\mathrm{\Omega }^{}\mathrm{\Sigma }\mathrm{\Omega }),`$
$`\left(\mathrm{\Omega }^{}\mathrm{\Sigma }\mathrm{\Omega }\right)={\displaystyle \underset{\mathrm{}=0}{\overset{N2}{}}}\overline{\mathrm{\Omega }}_{\mathrm{}}(1)^{\mathrm{}}\mathrm{\Omega }_{N2\mathrm{}}=\left[\widehat{K}(\rho )\overline{\varpi ({\displaystyle \frac{\rho }{2\pi i}};z)}\varpi (+{\displaystyle \frac{\rho }{2\pi i}};z)\right]|_{\rho ^{N2}}`$
$`=|\varpi _0|^2\times \left[\widehat{K}(\rho )e^{\rho (t\overline{t})}\mathrm{exp}\left({\displaystyle \underset{n2}{}}\rho ^n(x_n+(1)^n\overline{x}_n)\right)\right]|_{\rho ^{N2}}.`$ (4)
Here the first term in the brackets contains characteristic classes of the $`M`$. Its coefficients in the series expansion are represented as some combinations of Chern classes of the $`M`$. Also there appear Riemann’s zeta functions evaluated at positive odd integers. These are irrational numbers and might be transcendental numbers. It implies some arithmetic properties of this model. Next the second term is associated to the imaginary part of the $`t`$. When we translate the formal parameter $`\rho `$ into a divisor “$`D`$” of the hyperplane, the $`\rho (t\overline{t})`$ is identified with the Kähler form $`J`$ of $`M`$
$`\rho (t\overline{t})=2i\rho \text{Im}(t)2i\text{Im}(t)[D]=2iJ.`$
This second term contains only imaginary part of $`t`$ and is invariant under an arbitrary shift of $`\text{Re}(t)\text{Re}(t)+𝐚`$ for $`{}_{}{}^{}𝐚𝐑`$. This symmetry is a classical one and is broken at the quantum level. In fact, the third term in Eq.(4) in not invariant under this arbitrary shift because the $`x_n`$s are expressed as series expansions of the variable $`q=e^{2\pi it}`$. This 3rd term is invariant under only integral shifts of $`tt+n`$ with $`n𝐙`$. This term contains information about non-perturbative effects of the worldsheet instantons. They break the Peccei-Quinn symmetry into a symmetry under the integral shift of the $`\text{Re}(t)`$. The real part of the $`t`$ is related to the 2nd rank antisymmetric field $`B`$ in the NS-NS sector
$`\text{Re}(t)[D]=B,`$
and the shift of the $`\text{Re}(t)\text{Re}(t)+n`$ is equivalent to a shift of the $`B`$ field, $`BB+[D]`$.
Finally we will write down our result for the $`K`$
$`e^K=(1)^d\left({\displaystyle \frac{2\pi i}{N}}\right)^d{\displaystyle \frac{1}{N^2}}|\varpi _0|^2S_d(\stackrel{~}{y}_1,\stackrel{~}{y}_2,\mathrm{},\stackrel{~}{y}_d)`$
$`=(1)^d\left({\displaystyle \frac{2\pi i}{N}}\right)^d{\displaystyle \frac{1}{N^2}}|\varpi _0|^2{\displaystyle \frac{(t\overline{t})^d}{d!}}\times \left[1+{\displaystyle \underset{\mathrm{}=2}{\overset{d}{}}}\left(\begin{array}{c}d\\ \mathrm{}\end{array}\right){\displaystyle \frac{\mathrm{}!S_{\mathrm{}}(0,\stackrel{~}{y}_2,\mathrm{},\stackrel{~}{y}_{\mathrm{}})}{(t\overline{t})^{\mathrm{}}}}\right],`$
$`\widehat{k}_{2n+1}:=2{\displaystyle \frac{NN^{2n+1}}{2n+1}}{\displaystyle \frac{\zeta (2n+1)}{(2\pi i)^{2n+1}}}(n1),`$
$`\widehat{K}(\rho )=\mathrm{exp}\left({\displaystyle \underset{m1}{}}\widehat{k}_{2m+1}\rho ^{2m+1}\right)=1+{\displaystyle \underset{n1}{}}\widehat{K}_n\left({\displaystyle \frac{\rho }{2\pi i}}\right)^n,`$
$`y_n:=x_n+(1)^n\overline{x}_n(n2),`$
$`\stackrel{~}{y}_n=\{\begin{array}{cc}t\overline{t}& (n=1)\\ y_{2m}& (n=2m;m1)\\ y_{2m+1}+\widehat{k}_{2m+1}& (n=2m+1;m1)\end{array}.`$
The Schur functions contain “loop”<sup>1</sup><sup>1</sup>1The $`\frac{c_3\zeta (3)}{(2\pi i)^3}`$ is interpreted as an effect of loop corrections at the 4-loop perturbative calculation of the sigma model with 3 dimensional Calabi-Yau target spaces. We do not know that the other terms $`\widehat{K}_n`$s in the $`\widehat{K}`$ can be interpreted directly as perturbative loop corrections of the sigma model at higher loop calculations. and non-perturbative effects of the model
$`S_2=2\text{Re}(x_2),S_3=2i\text{Im}(x_3)+\widehat{k}_3,`$
$`S_4=2[\text{Re}(x_4)+(\text{Re}(x_2))^2],`$
$`S_5=2i\text{Im}(y_5)+4i\text{Re}(x_2)\text{Im}(x_3)+\widehat{k}_5+2\text{Re}(x_2)\widehat{k}_3.`$
Especially we obtain an asymptotic formula of the $`y_n`$s in the range of the large radius volume
$`\{\begin{array}{cc}\stackrel{~}{y}_{2m}0& (m1)\\ \stackrel{~}{y}_{2m+1}\widehat{k}_{2m+1}& (m1)\end{array}`$
$`S_{\mathrm{}}(0,\stackrel{~}{y}_2,\mathrm{},\stackrel{~}{y}_{\mathrm{}})\widehat{K}_{\mathrm{}}(\mathrm{}2).`$
In this limit, we can neglect non-perturbative corrections and write down the $`K`$, a metric and a scalar curvature of the moduli space of $`M`$ as power series of the $`(t\overline{t})`$
$`e^K=(1)^d\left({\displaystyle \frac{2\pi i}{N}}\right)^d{\displaystyle \frac{1}{N^2}}{\displaystyle \frac{1}{d!}}(t\overline{t})^d\left[1+{\displaystyle \underset{\mathrm{}=3}{\overset{d}{}}}\left(\begin{array}{c}d\\ \mathrm{}\end{array}\right){\displaystyle \frac{\mathrm{}!\widehat{K}_{\mathrm{}}}{(2\pi i)^{\mathrm{}}(t\overline{t})^{\mathrm{}}}}\right],`$
$`g_{t\overline{t}}={\displaystyle \frac{d}{(t\overline{t})^2}}\left[1{\displaystyle \frac{72\left(\begin{array}{c}d\\ 3\end{array}\right)\widehat{K}_3}{(2\pi i)^3d\left(t\overline{t}\right)^3}}{\displaystyle \frac{3600\left(\begin{array}{c}d\\ 5\end{array}\right)\widehat{K}_5}{(2\pi i)^5d\left(t\overline{t}\right)^5}}+\mathrm{}\right],`$
$`R_{t\overline{t}}=_t_{\overline{t}}\mathrm{log}g_{t\overline{t}}={\displaystyle \frac{2}{(t\overline{t})^2}}\left[1{\displaystyle \frac{432\left(\begin{array}{c}d\\ 3\end{array}\right)\widehat{K}_3}{(2\pi i)^3d\left(t\overline{t}\right)^3}}{\displaystyle \frac{54000\left(\begin{array}{c}d\\ 5\end{array}\right)\widehat{K}_5}{(2\pi i)^5d\left(t\overline{t}\right)^5}}+\mathrm{}\right],`$
$`R=2g^{t\overline{t}}R_{t\overline{t}}={\displaystyle \frac{4}{d}}\left[1{\displaystyle \frac{360\left(\begin{array}{c}d\\ 3\end{array}\right)\widehat{K}_3}{(2\pi i)^3d\left(t\overline{t}\right)^3}}{\displaystyle \frac{50400\left(\begin{array}{c}d\\ 5\end{array}\right)\widehat{K}_5}{(2\pi i)^5d\left(t\overline{t}\right)^5}}+\mathrm{}\right].`$
The line element of the A-model moduli space is given as $`ds^2=g_{t\overline{t}}dtd\overline{t}`$. The large $`\text{Im}(t)`$ limit, the metric and the Ricci tensor are inversely proportional to the $`(t\overline{t})`$. But apart from the point $`\text{Im}(t)=\mathrm{}`$, there appear terms $`\text{Im}(t)^n`$ ($`n>2`$) and also non-perturbative corrections with $`e^{2\pi int}`$ or $`e^{2\pi in\overline{t}}`$ ($`n1`$) for $`d3`$ cases. The large $`\text{Im}(t)`$ point corresponds to the large $`\psi `$ point. On the other hand, a related value of the “$`t`$” at the $`\psi =0`$ point does not vanish, but it is finite
$`t(\psi =0)={\displaystyle \frac{1}{2}}+{\displaystyle \frac{i}{2\mathrm{tan}\frac{\pi }{N}}}.`$
At the $`\psi =0`$ point, the $`R`$, $`R_{t\overline{t}}`$ and $`g_{t\overline{t}}`$ do not vanish and are evaluated as
$`g_{t\overline{t}}=\left(2\mathrm{sin}{\displaystyle \frac{2\pi }{N}}\right)^2\left(2\mathrm{cos}{\displaystyle \frac{\pi }{N}}\right)^{N+2},`$
$`R=4+2\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{N}\right)\mathrm{\Gamma }\left(\frac{3}{N}\right)}{\mathrm{\Gamma }\left(1\frac{1}{N}\right)\mathrm{\Gamma }\left(1\frac{3}{N}\right)}}\right]^N\left[{\displaystyle \frac{\mathrm{\Gamma }\left(1\frac{2}{N}\right)}{\mathrm{\Gamma }\left(\frac{2}{N}\right)}}\right]^{2N},`$
$`R_{t\overline{t}}={\displaystyle \frac{1}{2}}g_{t\overline{t}}R.`$
The scalar curvature is positive around the Gepner point for $`d3`$. When we increase the $`\psi `$ from zero to one, the $`R`$ increases monotonically with the $`\psi `$ for each $`N`$. The point $`\psi =1`$ is a singular point where associated scalar curvatures blow up. When one passes through the $`\psi =1`$ and increases the value of the $`\psi `$, the associated $`R`$ decreases monotonically and vanishes at some point $`\psi =\psi _0`$ for each $`N`$. In the range $`\psi >\psi _0`$, the $`R`$ is negative and its asymptotic value is $`4/d`$. That is to say, it is negative in the large $`\psi `$ region for $`d3`$. The behaviors of the curvatures for $`N=5,6,7`$ are shown in Fig.1
For the torus case, the curvature is $`R=4`$ both in the small $`\psi `$ and in the large $`\psi `$ regions. Similarly, for the K3 case, the $`R`$ is $`(2)`$ in the two limits. Associated Ricci tensors at the $`\psi =0`$ are obtained for $`N=3,4`$ respectively
$`R_{\psi \overline{\psi }}=18\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{2}{3}\right)}{\mathrm{\Gamma }\left(\frac{1}{3}\right)}}\right]^6=0.30021677774546778674\mathrm{},(N=3\text{ case}),`$
$`R_{\psi \overline{\psi }}=16\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{3}{4}\right)}{\mathrm{\Gamma }\left(\frac{1}{4}\right)}}\right]^4=0.20880017792822456419\mathrm{},(N=4\text{ case}).`$
When we increase the $`\psi `$ from zero to one, the $`R_{\psi \overline{\psi }}`$ decreases monotonically with the $`\psi `$ for each $`N`$. The point $`\psi =1`$ is a singular point where the $`R_{\psi \overline{\psi }}`$s blow up. When one passes through the $`\psi =1`$ and increases the value of the $`\psi `$, the associated $`R_{\psi \overline{\psi }}`$ decreases monotonically for each $`N`$ as shown in Fig.2. But the scalar curvatures are constants except for a point $`\psi =1`$.
## 5 Two-Point Functions
The other constituent blocks of the topological model are three-point couplings $`\{\kappa _{\mathrm{}}\}`$ and two-point functions. Each A-model operator $`𝒪^{(\mathrm{})}`$ is associated with a cohomology element $`e_{\mathrm{}}\text{H}^2\mathrm{}`$(M). The $`\kappa _{\mathrm{}}`$ is a fusion coupling of $`𝒪^{(1)}`$ and $`𝒪^{(\mathrm{})}`$
$`𝒪^{(1)}𝒪^{(\mathrm{})}=\kappa _{\mathrm{}}𝒪^{(\mathrm{}+1)}(\mathrm{}=0,1,\mathrm{},d1).`$
$`𝒪^{(1)}𝒪^{(d)}=0,`$
$`\kappa _0=1,`$
$`\kappa _{\mathrm{}}={\displaystyle \frac{1}{\kappa _\mathrm{}1}}{\displaystyle \frac{1}{\kappa _\mathrm{}2}}\mathrm{}{\displaystyle \frac{1}{\kappa _1}}{\displaystyle \frac{1}{\kappa _0}}S_{\mathrm{}+1}(t,x_2,x_3,\mathrm{},x_{\mathrm{}+1})`$
$`=1+𝒪(q),(1\mathrm{}d1).`$
Also the topological metric $`\eta _{\mathrm{},m}`$ is given as
$`\eta _{\mathrm{},m}=𝒪^{(\mathrm{})}𝒪^{(m)}=N\delta _{\mathrm{}+m,d}.`$
On the other hand, the hermitian two-point functions $`\overline{𝒪}^{(\mathrm{})}|𝒪^{(m)}`$ are calculated by using a method of the $`tt^{}`$-fusion. In our case, the correlators have diagonal forms and are given as
$`\overline{𝒪}^{(\mathrm{})}|𝒪^{(m)}=e^q_{\mathrm{}}\delta _{\mathrm{},m}(0\mathrm{}d,\mathrm{\hspace{0.17em}0}mm),`$
$`\stackrel{~}{q}_0=q_0,`$
$`\stackrel{~}{q}_{\mathrm{}}=q_{\mathrm{}}+{\displaystyle \underset{n=0}{\overset{\mathrm{}1}{}}}\mathrm{log}|\kappa _n|^2(\mathrm{}1),`$
$`\overline{}\stackrel{~}{q}_0+e^{\stackrel{~}{q}_1\stackrel{~}{q}_0}=0,`$
$`\overline{}\stackrel{~}{q}_{\mathrm{}}+e^{\stackrel{~}{q}_{\mathrm{}+1}\stackrel{~}{q}_{\mathrm{}}}e^{\stackrel{~}{q}_{\mathrm{}}\stackrel{~}{q}_\mathrm{}1}=0(1\mathrm{}d1),`$
$`\overline{}\stackrel{~}{q}_de^{\stackrel{~}{q}_d\stackrel{~}{q}_{d1}}=0.`$
In the large radius limit, the $`\stackrel{~}{q}_{\mathrm{}}`$s behave as
$`e^{\stackrel{~}{q}_0}=e^K=(1)^d\left({\displaystyle \frac{2\pi i}{N}}\right)^d{\displaystyle \frac{1}{N^2}}{\displaystyle \frac{(t\overline{t})^d}{d!}}+\mathrm{},`$
$`e^{\stackrel{~}{q}_{\mathrm{}+1}\stackrel{~}{q}_{\mathrm{}}}={\displaystyle \frac{(\mathrm{}+1)(d\mathrm{})}{(t\overline{t})^2}}+\mathrm{}(0\mathrm{}d1).`$
In this case, any corrections are suppressed and normalized two-point functions are inversely proportional to the $`\text{Im}(t)^2\mathrm{}`$ ($`\mathrm{}=0,1,2,\mathrm{},d`$)
$`{\displaystyle \frac{\overline{𝒪}^{(\mathrm{})}|𝒪^{(\mathrm{})}}{\overline{𝒪}^{(0)}|𝒪^{(0)}}}=e^{\stackrel{~}{q}_{\mathrm{}}\stackrel{~}{q}_0}=(1)^{\mathrm{}}(\mathrm{}!)^2\left(\begin{array}{c}d\\ \mathrm{}\end{array}\right){\displaystyle \frac{1}{(t\overline{t})^2\mathrm{}}}+\mathrm{}(0\mathrm{}d).`$
In the finite $`\text{Im}(t)`$ case, there appear corrections in power series and non-perturbative corrections of exponential types. Those are abbreviated as a symbol “$`\mathrm{}`$” in the above formula. Even in that generic case, these $`q_{\mathrm{}}`$s are described by combining the curvature $`R`$, the metric $`g_{t\overline{t}}`$ and the $`K`$
$`e^{\stackrel{~}{q}_0}=e^K,e^{\stackrel{~}{q}_1\stackrel{~}{q}_0}=g_{t\overline{t}},e^{\stackrel{~}{q}_2\stackrel{~}{q}_1}=g_{t\overline{t}}\left({\displaystyle \frac{R}{2}}+2\right),`$
$`e^{\stackrel{~}{q}_3\stackrel{~}{q}_2}=g_{t\overline{t}}\left[3\left({\displaystyle \frac{R}{2}}+1\right)g^{t\overline{t}}_t\overline{}_{\overline{t}}\mathrm{log}\left({\displaystyle \frac{R}{2}}+2\right)\right],`$
$`e^{\stackrel{~}{q}_4\stackrel{~}{q}_3}=g_{t\overline{t}}[4+3R2g^{t\overline{t}}_t\overline{}_{\overline{t}}\mathrm{log}({\displaystyle \frac{R}{2}}+2)`$
$`g^{t\overline{t}}_t\overline{}_{\overline{t}}\mathrm{log}[3({\displaystyle \frac{R}{2}}+1)g^{t\overline{t}}_t\overline{}_{\overline{t}}\mathrm{log}({\displaystyle \frac{R}{2}}+2)]],`$
$`\mathrm{}.`$
We know the formula of the $`K`$, $`R`$, and $`g_{t\overline{t}}`$ and can evaluate moduli dependences of these correlators.
## 6 Quintic
In this section, we investigate the $`K`$ for a quintic $`M`$ case and compare our results with those by Candelas et al. The set of periods $`\{\mathrm{\Omega }_{\mathrm{}}\}`$ of an associated mirror $`W`$ is expressed by using the $`\varpi `$
$`\varpi \left({\displaystyle \frac{\rho }{2\pi i}}\right)\sqrt{\widehat{K}(\rho )}=:{\displaystyle \underset{\mathrm{}0}{}}\rho ^{\mathrm{}}\mathrm{\Omega }_{\mathrm{}},`$
$`\sqrt{\widehat{A}(\rho )}=1+{\displaystyle \frac{5}{12}}\rho ^2,\sqrt{\widehat{K}(\rho )}=1{\displaystyle \frac{40}{(2\pi i)^3}}\zeta (3)\rho ^3,`$
$`\widehat{c}={\displaystyle \frac{40}{(2\pi i)^3}}\zeta (3),`$
$`\mathrm{\Omega }=\left(\begin{array}{c}\mathrm{\Omega }_0\\ \mathrm{\Omega }_1\\ \mathrm{\Omega }_2\\ \mathrm{\Omega }_3\end{array}\right)=\left(\begin{array}{c}1\\ t\\ \frac{1}{2}t^2+S_2(0,x_2)\\ \frac{1}{6}t^3\widehat{c}+tS_2(0,x_2)+S_3(0,x_2,x_3)\end{array}\right).`$ (5)
The prepotential $`F`$ of $`M`$ is expressed as a sum of a polynomial part of $`t`$ and a non-perturbative part $`f`$
$`F={\displaystyle \frac{\kappa }{6}}t^3+{\displaystyle \frac{1}{2}}at^2+bt+{\displaystyle \frac{1}{2}}c+f,`$
$`a={\displaystyle \frac{11}{2}},b={\displaystyle \frac{25}{12}},c={\displaystyle \frac{c_3\zeta (3)}{(2\pi i)^3}}=5\widehat{c},c_3=200.`$
The effects of the instantons are encoded in this function $`f`$. The $`F`$ leads to a a canonical set of basis $`\{\mathrm{\Pi }_{\mathrm{}}\}`$. The $`\mathrm{\Pi }_{\mathrm{}}`$s are represented as some linear combinations of periods
$`\mathrm{\Pi }=\left(\begin{array}{c}\mathrm{\Pi }_0\\ \mathrm{\Pi }_1\\ \mathrm{\Pi }_2\\ \mathrm{\Pi }_4\end{array}\right)=\left(\begin{array}{c}1\\ t\\ _tF\\ t_tF2F\end{array}\right)=𝒩\mathrm{\Omega },`$ (6)
$`𝒩=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ b& a& \kappa & 0\\ 0& b& 0& \kappa \end{array}\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ \frac{25}{12}& \frac{11}{2}& 5& 0\\ 0& \frac{25}{12}& 0& 5\end{array}\right),`$
$`\kappa =5,a=11/2,b=25/12.`$
By comparing these two approaches Eqs.(5),(6), we obtain relations for the $`c`$ and the $`f`$
$`c=\kappa \widehat{c}={\displaystyle \frac{c_3\zeta (3)}{(2\pi i)^3}},`$
$`f={\displaystyle \frac{\kappa }{2}}S_3(0,x_2,x_3),`$
$`_tf=\kappa S_2(0,x_2).`$
Then the Kähler potential $`K`$ in the A-model is evaluated by using the $`F`$ up to a overall normalization factor with a suitable choice of a section of an associated line bundle
$`e^K=(t\overline{t})(F+\overline{}\overline{F})2(F\overline{F})`$
$`={\displaystyle \frac{\kappa }{6}}(t\overline{t})^3+(\overline{a}a)t\overline{t}(b\overline{b})(t\overline{t})(c\overline{c})+(t\overline{t})(f+\overline{}\overline{f})2(f\overline{f})`$
$`={\displaystyle \frac{\kappa }{6}}(t\overline{t})^32c+(t\overline{t})(f+\overline{}\overline{f})2(f\overline{f}),`$
$`=\kappa \left[{\displaystyle \frac{1}{6}}(t\overline{t})^32\widehat{c}+(t\overline{t})(S_2+\overline{S}_2)+(S_3\overline{S}_3)\right]`$
$`=(1)^3{\displaystyle _{\text{CY}_3}}\left[\widehat{K}(\rho )\overline{\varpi ({\displaystyle \frac{\rho }{2\pi i}})}\varpi (+{\displaystyle \frac{\rho }{2\pi i}})\right]|_{\rho =[D]},`$
$`{\displaystyle _{\text{CY}_3}}[D][D][D]=\kappa .`$
Here we used the facts the $`a`$ and $`b`$ are real numbers and the $`c`$ is a pure imaginary number. This formula coincides with our formula Eq.(4). Now we make a remark here: The basis $`\mathrm{\Omega }`$ does not coincide with the $`\mathrm{\Pi }`$, but it is related to the $`\mathrm{\Pi }`$ by a kind of a symplectic transformation, that is represented as a matrix $`𝒩`$
$`\mathrm{\Pi }=𝒩\mathrm{\Omega },𝒩^t\mathrm{\Sigma }𝒩=(5)\mathrm{\Sigma }.`$
It leads to the same $`K`$ up to a multiplicative factor because the following relation is satisfied
$`e^K\mathrm{\Pi }^{}\mathrm{\Sigma }\mathrm{\Pi }=\mathrm{\Omega }^{}𝒩^t\mathrm{\Sigma }𝒩\mathrm{\Omega }=(5)\mathrm{\Omega }^{}\mathrm{\Sigma }\mathrm{\Omega }.`$
## 7 Generalization
In this section, we propose a formula of the $`K`$ by generalizing the previous result of the Fermat type. We consider complete intersections $`M`$ of $`\mathrm{}`$ hypersurfaces $`\{p_j=0\}`$ in products of $`k`$ projective spaces $``$
$`M:=\left(\begin{array}{c}𝐏^{n_1}(w_1^{(1)},\mathrm{},w_{n_1+1}^{(1)})\\ \mathrm{}\\ 𝐏^{n_k}(w_1^{(k)},\mathrm{},w_{n_1+1}^{(k)})\end{array}|\right|\begin{array}{c}d_1^{(1)}\mathrm{}d_{\mathrm{}}^{(1)}\\ \mathrm{}\\ d_1^{(k)}\mathrm{}d_{\mathrm{}}^{(k)}\end{array}).`$
The $`d_j^{(i)}`$ are degrees of the coordinates of $`𝐏^{n_i}(w_1^{(i)},\mathrm{},w_{n_1+1}^{(i)})`$ in the $`j`$-th polynomial $`p_j`$ ($`i=1,2,\mathrm{},k;j=1,2,\mathrm{},\mathrm{}`$). We propose a formula of the Kähler potential of the A-model up to some overall normalization factors
$`e^K=(1)^d{\displaystyle _M}\left[\widehat{K}(\lambda )\overline{\varpi ({\displaystyle \frac{\lambda [D_i]}{2\pi i}};z_k)}\varpi (+{\displaystyle \frac{\lambda [D_i]}{2\pi i}};z_k)\right]_{\lambda ^d}\times |\varpi _0|^2`$
$`=(1)^d{\displaystyle _{}}\left[\widehat{K}(\lambda )\overline{\varpi ({\displaystyle \frac{\lambda [D_i]}{2\pi i}};z_k)}\varpi (+{\displaystyle \frac{\lambda [D_i]}{2\pi i}};z_k)\right]_{\lambda ^d}[H]\times |\varpi _0|^2,`$
$`\varpi (v;z):={\displaystyle \underset{n}{}}{\displaystyle \frac{a(n+v)}{a(v)}}z^{n+v},z^{n+v}:={\displaystyle \underset{k=1}{\overset{p}{}}}z_k^{n_k+v_k},`$
$`a(n+v):={\displaystyle \frac{{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }\left(1+{\displaystyle \underset{i=1}{\overset{k}{}}}d_j^{(i)}(n_i+v_i)\right)}{{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j^{}=1}{\overset{n_i+1}{}}}\mathrm{\Gamma }\left(1+w_j^{}^{(i)}(n_i+v_i)\right)}},`$
$`\widehat{K}(\lambda ):={\displaystyle \frac{a\left(+{\displaystyle \frac{\lambda [D]}{2\pi i}}\right)}{a\left({\displaystyle \frac{\lambda [D]}{2\pi i}}\right)}}=\mathrm{exp}\left[+2{\displaystyle \underset{m=1}{}}{\displaystyle \frac{\zeta (2m+1)}{2m+1}}\left({\displaystyle \frac{\lambda }{2\pi i}}\right)^{2m+1}X_{2m+1}\right],`$
$`X_n:={\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j^{}=1}{\overset{n_i+1}{}}}(w_j^{}^{(i)}[D_i])^n{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \underset{i=1}{\overset{k}{}}}d_j^{(i)}[D_i]\right)^n,[H]={\displaystyle \frac{{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \underset{i=1}{\overset{k}{}}}d_j^{(i)}[D_i]\right)}{{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j^{}=1}{\overset{n_i+1}{}}}w_j^{}^{(i)}}},`$
$`d=\mathrm{}+{\displaystyle \underset{i=1}{\overset{k}{}}}n_i,(\text{dimension}).`$
The $`[D_i]`$s ($`i=1,2,\mathrm{},k`$) are Poincaré duals of divisors “$`D_i`$” of the model. We will rewrite these in order to interpret them in the A-model language. First mirror maps are given as
$`2\pi it^i=\mathrm{log}z^i+_{v_i}\mathrm{log}\widehat{\varpi },`$
$`\widehat{\varpi }={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a(n+v)}{a(v)}}z^n.`$
Then the normalized $`\varpi `$ is expressed as
$`\varpi _0^1\varpi \left({\displaystyle \frac{\lambda [D_i]}{2\pi i}}\right)=\mathrm{exp}\left(\lambda [D_i]t^i+{\displaystyle \underset{\mathrm{}=2}{}}\lambda ^{\mathrm{}}x_{\mathrm{}}\right),`$
$`x_{\mathrm{}}={\displaystyle \frac{1}{(2\pi i)^{\mathrm{}}}}{\displaystyle \frac{1}{\mathrm{}!}}([D])^{\mathrm{}}\mathrm{log}\widehat{\varpi },`$
$`[D]:={\displaystyle \underset{i}{}}[D_i]{\displaystyle \frac{}{v_i}},`$
$`\widehat{K}(\lambda )=\mathrm{exp}\left({\displaystyle \underset{m=1}{}}\lambda ^{2m+1}\widehat{k}_{2m+1}\right),`$
$`\widehat{k}_{2m+1}=2{\displaystyle \frac{X_{2m+1}}{2m+1}}{\displaystyle \frac{\zeta (2m+1)}{(2\pi i)^{2m+1}}}.`$
By using the above relations, we can express the Kähler potential in a compact form
$`e^K=(1)^d{\displaystyle _{}}S_d(\stackrel{~}{y}_1,\stackrel{~}{y}_2,\mathrm{},\stackrel{~}{y}_d)[H],`$ (7)
$`\stackrel{~}{y}_1={\displaystyle \underset{i}{}}[D_i](t^i\overline{t}^i),`$
$`\stackrel{~}{y}_{2n}=x_{2n}+\overline{x}_{2n},(n1),`$
$`\stackrel{~}{y}_{2n+1}=x_{2n+1}\overline{x}_{2n+1}+\widehat{k}_{2n+1},(n1).`$
To confirm the validity of this formula, we restrict ourselves to the $`3`$-fold case and consider the $`K`$. First let us consider a general type of Calabi-Yau $`3`$-fold $`M`$ with Kähler parameters $`t^i`$ ($`i=1,2,\mathrm{},h^{1,1}`$). Its prepotential $`F`$ is described as
$`F={\displaystyle \frac{1}{6}}\kappa _{ijk}t^it^jt^k+{\displaystyle \frac{1}{2}}a_{ij}t^it^j+b_it^i+{\displaystyle \frac{1}{2}}c+f,`$
$`e^K=(t^i\overline{t}^i)(_iF+\overline{}_i\overline{F})2(F\overline{F})`$
$`={\displaystyle \frac{1}{6}}\kappa _{ijk}(t^i\overline{t}^i)(t^j\overline{t}^j)(t^k\overline{t}^k)(a_{ij}\overline{a}_{ji})t^i\overline{t}^j`$
$`(b_i\overline{b}_i)(t^i\overline{t}^i)(c\overline{c})+(t^i\overline{t}^i)(_if+\overline{}_i\overline{f})2(f\overline{f}).`$
When we impose a condition that the $`e^K`$ is invariant under a constant shift of $`t^it^i+1`$ for an arbitrary $`i`$, the term $`(a_{ij}\overline{a}_{ji})t^i\overline{t}^j`$ must vanish. It means that the $`a_{ij}`$s must be hermitian as components of a matrix
$`\overline{a}_{ij}=a_{ji}.`$
Also the $`b_i`$s are related to a second Chern class $`c_2`$ of the $`M`$
$`b_i={\displaystyle \frac{1}{24}}{\displaystyle _M}c_2J_i={\displaystyle _M}\sqrt{\widehat{A}}J_i,`$
and they are real numbers. On the other hand, the $`c`$ is a pure imaginary number and is associated with a 3rd Chern class
$`c={\displaystyle \frac{c_3\zeta (3)}{(2\pi i)^3}}.`$
Collecting all these facts, we can rewrite the $`K`$
$`e^K={\displaystyle \frac{1}{6}}\kappa _{ijk}(t^i\overline{t}^i)(t^j\overline{t}^j)(t^k\overline{t}^k)2c`$
$`+(t^i\overline{t}^i)(_if+\overline{}_i\overline{f})2(f\overline{f}).`$ (8)
The Eq.(8) is obtained by using a prepotential $`F`$ for a $`3`$-fold. In contrast, we can obtain our result for the formula $`K`$ without requiring an existence of the prepotential. Let us write down our result Eq.(7) in the 3 dimensional case. All we have to do is to evaluate the $`S_3`$ in this case
$`S_3={\displaystyle \frac{1}{6}}[D_i][D_j][D_k](t\overline{t})^i(t\overline{t})^j(t\overline{t})^k`$
$`+[D_i](t\overline{t})^i(x_2+\overline{x}_2)+(x_3\overline{x}_3)+\widehat{k}_3.`$
Then the $`K`$ is obtained
$`e^K={\displaystyle \frac{1}{6}}\kappa _{ijk}(t\overline{t})^i(t\overline{t})^j(t\overline{t})^k`$
$`+(t\overline{t})^i(_{t^i}f+_{\overline{t}^i}\overline{f})2(f\overline{f})2{\displaystyle \frac{c_3\zeta (3)}{(2\pi i)^3}},`$ (9)
$`\kappa _{ijk}={\displaystyle _{}}[D_i][D_j][D_k][H],`$
$`_{t^i}f={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(2\pi i)^2}}\kappa _{ijk}{\displaystyle \frac{}{v_j}}{\displaystyle \frac{}{v_k}}\mathrm{log}\widehat{\varpi },`$
$`f={\displaystyle \frac{1}{12}}{\displaystyle \frac{1}{(2\pi i)^3}}\kappa _{ijk}{\displaystyle \frac{}{v_i}}{\displaystyle \frac{}{v_j}}{\displaystyle \frac{}{v_k}}\mathrm{log}\widehat{\varpi },`$
$`c_3={\displaystyle \frac{1}{3}}\left[{\displaystyle \underset{i}{}}{\displaystyle \underset{j^{}}{}}(w_j^{}^{(i)})^3\kappa _{iii}{\displaystyle \underset{i,j,k}{}}{\displaystyle \underset{j^{}}{}}d_j^{}^{(i)}d_j^{}^{(j)}d_j^{}^{(k)}\kappa _{ijk}\right].`$
This formula Eq.(9) coincides with that of the $`3`$-fold case, Eq.(8). But Eq.(7) is not restricted to this 3 dimensional case. In fact, Eq.(7) is applicable to smooth $`d`$-folds cases because the periods $`\sqrt{\widehat{K}}\varpi `$ have appropriate intersection forms
$`\widehat{K}(\lambda )\overline{\varpi \left({\displaystyle \frac{\lambda [D_i]}{(2\pi i)}}\right)}\varpi \left(+{\displaystyle \frac{\lambda [D_i]}{(2\pi i)}}\right)|_{\lambda ^d}={\displaystyle \underset{n,n^{}=0}{\overset{d}{}}}\overline{\mathrm{\Omega }}_n\mathrm{\Sigma }_{n,n^{}}\mathrm{\Omega }_n^{},`$
$`\mathrm{\Omega }_m={\displaystyle \frac{1}{m!}}{\displaystyle \frac{1}{(2\pi i)^m}}([D])^m\left[\sqrt{\widehat{K}}\varpi \right],\mathrm{\Sigma }_{n,n^{}}=(1)^n\delta _{n+n^{},d}.`$
For singular $`d`$-fold case, we might have to modify some parts which have information about intersection numbers associated with the divisors.
## 8 Conclusions and Discussions
In this article, we develop a method to calculate the Kähler potential in the topological A-model. Generally the Kähler potential $`K`$ is represented as $`(t\overline{t})^d`$ when no quantum corrections exist. But it is known that there are loop corrections in the two dimensional $`N=2`$ non-linear sigma models with Calabi-Yau target spaces. A term $`(t\overline{t})^{d3}\times {\displaystyle \frac{c_3\zeta (3)}{(2\pi i)^3}}`$ in the $`K`$ reflects a perturbative correction at the four loop calculation. The $`c_3`$ is the 3rd Chern class of the CY. In our case, the $`\widehat{K}`$ seems to describe loop corrections to the sigma model. In general, there might be more corrections and they could be controlled by one function $`\widehat{K}`$. It is interesting to give a field theoretical interpretation to these $`\widehat{K}_n`$s from the point of view of direct higher loop calculations.
The basis we pick here is a symplectic (or an SO-invariant) basis with intersection matrix $`\mathrm{\Sigma }`$. But it is not integrable but rational basis. In order to obtain a set of canonical basis, we have to do some linear transformation on the $`\mathrm{\Omega }`$. It is needed to discuss the D-branes charges or central charges in BPS mass formulae. We will study these topics in the next paper.
## Acknowledgment
This work is supported by the Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture 10740117.
## Appendix A Examples of $`\widehat{K}`$
We write down several concrete examples of the $`\widehat{K}_n`$ for the $`d`$-fold $`M`$. First the generating function $`\widehat{K}(\rho )`$ is defined by using Riemann’s zeta functions
$`\widehat{K}(\rho )=\mathrm{exp}\left(2{\displaystyle \underset{m=1}{}}{\displaystyle \frac{NN^{2m+1}}{2m+1}}\zeta (2m+1)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^{2m+1}\right)`$
$`=1+{\displaystyle \underset{n=1}{\overset{d}{}}}\widehat{K}_n\left({\displaystyle \frac{\rho }{2\pi i}}\right)^n,N=d+2.`$
The coefficients in the series expansion are represented as some combinations of Chern classes of the $`M`$
$`\widehat{K}_1=\widehat{K}_2=0,\widehat{K}_3={\displaystyle \frac{2}{N}}c_3\zeta (3),\widehat{K}_4=0,`$
$`\widehat{K}_5=\left({\displaystyle \frac{2}{N}}c_5{\displaystyle \frac{2}{N^2}}c_3c_2\right)\zeta (5),\widehat{K}_6={\displaystyle \frac{2}{N^2}}c_3^2\zeta (3)^2,`$
$`\widehat{K}_7=\left({\displaystyle \frac{2}{N}}c_7{\displaystyle \frac{2}{N^2}}c_5c_2{\displaystyle \frac{2}{N^2}}c_4c_3+{\displaystyle \frac{2}{N^3}}c_3c_2^2\right)\zeta (7),`$
$`\widehat{K}_8=\left({\displaystyle \frac{4}{N^2}}c_5c_3{\displaystyle \frac{4}{N^3}}c_3^2c_2\right)\zeta (3)\zeta (5).`$
A finite number of $`\widehat{K}_n`$s ($`nd`$) appear for the $`d`$-fold case. The $`\widehat{K}_1`$ and $`\widehat{K}_2`$ always vanish and the $`\widehat{K}(\rho )`$ is identity for the torus and K3 cases ($`N=3,4`$ respectively). We will summarize several examples for lower dimensional cases
$`\widehat{K}(N=3)=1,\widehat{K}(N=4)=1,`$
$`\widehat{K}(N=5)=180\zeta (3)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^3,\widehat{K}(N=6)=1140\zeta (3)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^3,`$
$`\widehat{K}(N=7)=1224\zeta (3)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^36720\zeta (5)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^5,`$
$`\widehat{K}(N=8)=1336\zeta (3)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^313104\zeta (5)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^5+56448\zeta (3)^2\left({\displaystyle \frac{\rho }{2\pi i}}\right)^6,`$
$`\widehat{K}(N=9)=1480\zeta (3)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^323616\zeta (5)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^5`$
$`+115200\zeta (3)^2\left({\displaystyle \frac{\rho }{2\pi i}}\right)^61366560\zeta (7)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^7,`$
$`\widehat{K}(N=10)=1660\zeta (3)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^339996\zeta (5)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^5`$
$`+217800\zeta (3)^2\left({\displaystyle \frac{\rho }{2\pi i}}\right)^62857140\zeta (7)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^7+26397360\zeta (3)\zeta (5)\left({\displaystyle \frac{\rho }{2\pi i}}\right)^8.`$ |
warning/0003/gr-qc0003050.html | ar5iv | text | # Boost-rotation symmetric type D radiative metrics in Bondi coordinates.
## 1 Introduction
Boost-rotation symmetric spacetimes describe “uniformly accelerated particles” that approach the speed of light asymptotically. The smoothness of the solution requires the spacetimes to be reflection symmetric; therefore at least two particles with opposite acceleration are present, and thus future null infinity contains at least two singular points. Furthermore, these solutions are found to be time symmetric, thus the puncturing of null infinity is present at both $`I^+`$ and $`I^{}`$. The null infinity of a boost-rotation symmetric spacetime can be global, in the sense that it admits spherical cuts, but the generators are not complete .
This family of spacetimes possess an Abelian $`G_2`$ of isometries, in which one of the Killing vector fields has closed orbits, and the other symmetry is the curved spacetime generalization of the boost rotations of Minkowski spacetime that leave invariant the null cone with vertex at the origin . The boost-rotation symmetry is of special interest because it is the only other continuous isometry an axially symmetric spacetime can have that does not exclude the possibility of radiation . Bičák, Hoenselaers & Schmidt have shown how to systematically construct an infinite number of these spacetimes by prescribing the multipolar structure of the particles undergoing the hyperbolic motion .
One of the most well known representatives of the family of boost-rotation symmetric spacetimes is the C-metric, a type D spacetime with hypersurface orthogonal Killing vectors which has been interpreted as describing a pair of black holes receding from one another, and joined by a singular axis . The C-metric can be generalized to include Maxwell field with principal null directions aligned with the pairs of repeated principal null directions of the Weyl tensor. Early works on boost-rotation symmetric spacetimes demanded the Killing vector fields to be hypersurface orthogonal, but as discussed by Bičák and Pravdová , this hypothesis can be set aside. Thus, another possible generalization is the twisting C-metric . In the charged version it is found that the two accelerated black holes have opposite charge so that the overall spacetime is electrically neutral . Ashtekar & Dray have shown that the C-metric admits a conformal completion such that the cuts of $`I`$ are isomorphic to the 2-sphere, and therefore it can be regarded as asymptotically flat at null infinity. Subsequent work by Dray showed that the charged C-metric is also asymptotically flat at spatial infinity . He also showed that the symmetries of the spacetime give rise to a vanishing ADM mass.
The boost-rotation symmetric spacetimes with hypersurface orthogonal Killing vectors are usually given in a coordinate system that clearly exhibits their symmetries . In these symmetry-adapted coordinates, the metric of the spacetime contains only two functions. The different particular examples can be obtained from boost-rotation symmetric solutions of the flat spacetime wave equation with sources. Once one of these is prescribed, the metric functions can be obtained by quadratures. If one wants to study the radiative properties of these spacetimes (radiation patterns, mass loss, Newman-Penrose constants), then one has to rewrite the spacetime in terms of Bondi coordinates like the ones discussed in chapter 3. The transformation between the symmetry adapted coordinate system and the Bondi coordinates used in the asymptotic expansions of the gravitational field has to be given in terms of series . These expansions are extremely messy, and usually only the leading terms can be calculated explicitly . To add to the problem, the coefficients in terms of which some quantities of interest like the Newman-Penrose constants are defined are found deep into the series expansions. One is bound to look for better methods to calculate the quantities of physical interest.
Bičák and Pravdová have obtained a series of constraint equations for the news function and the mass aspect of electrovacuum boost-rotation symmetric spacetimes. The constraint equations can be solved so that the news and the mass aspect depend on arbitrary functions of the ratio $`w=\mathrm{sin}\theta /u`$. These functions have, of course, to satisfy the field equations. Therefore, if one makes some extra assumptions about the spacetime one can obtain a closed system of ordinary differential equations for the mass aspect, the gravitational and electromagnetic news function.
Taking the C-metric as paradigm, our attention will be restricted to type D spacetimes. For electrovacuum spacetimes it will be necessary to make a further assumption about the principal null directions of the electromagnetic field; it will be taken to be algebraically general, but with null principal directions parallel to the pairs of null principal directions of the Weyl tensor. In order to be able to make some discussion on the polarization of the gravitational and electromagnetic radiation fields, the Killing vector fields will not be assumed hypersurface orthogonal. The family of boost-rotation spacetimes that will come out of these assumptions will clearly contain the C-metric as a particular case: hypersurface orthogonal Killing vectors, and no electromagnetic field.
Some analysis of the radiative properties of the spinning, charged C-metric of Plebanski & Demianski was done by Farhoosh & Zimmerman . They wrote the spacetime in terms of Bondi coordinates and evaluated some radiative properties of the spacetime (the news function and the mass aspect) in the linear regime. However, they failed to notice the global character of the solution, and thus for example they did not find that the overall electromagnetic charge of the spacetime vanishes. More recently, there has been some work on the spinning C-metric by Bičák & Pravda in which they discuss briefly the radiative character of this spacetime, and show how it indeed describes two uniformly accelerated, spinning black holes connected by a conical singularity, or with conical singularities extending from each horizon to infinity.
At this point it is worth making a note about the radiative character of the C-metric in particular, and the boost-rotation symmetric spacetimes as a whole. A boost-rotation symmetric spacetime can be divided in three regions depending on whether the Killing vector is timelike, null or spacelike. The region where the Killing field is null is known as the *roof*, and it intersects $`I`$ at the cuts where the particles in hyperbolic motion puncture null infinity. Below the roof, the Killing vector is timelike, and thus the spacetime is stationary in this region. Above the roof, the Killing vector field is timelike, so that the spacetime can be radiative. In this region, the boost-rotation symmetric spacetimes can be shown to be locally isometric to cylindrical waves (Einstein-Rosen waves) . It is also important to note that these solutions have a time reflexion symmetry. Thus, what is true for $`I^+`$ is also true for $`I^{}`$. Hence, the particles come from past null infinity, leaving a puncture on it. They approach to each other, and then they recede again. Finally, they again puncture null infinity (but now at $`I^+`$).
Previous attempts at addressing the radiative properties of boost-rotation symmetric spacetime proved to be unfruitful, mainly because they relied on finding a transformation between the symmetry adapted coordinates and the Bondi coordinates. Unfortunately, it is not possible in general to express the coordinate transformation in a closed form. This difficulties have motivated us to adopt a rather different approach here, as we construct the solutions by using series expansions. Here we take advantage of the series of constraint equations for the news function and the mass aspect of electrovacuum boost-rotation symmetric spacetimes found elsewhere . These equations are the tools which will help us gain some insight in the structure and properties of the Newman-Penrose constants of the mentioned class of radiative spacetimes. Only those solutions to the constraint equations can be solved such that the news and the mass aspect depend on an arbitrary function of the ratio $`w=\mathrm{sin}\theta /u`$ are considered here. Needless to say that these functions will have to satisfy the field equations. Under the assumptions made, we obtain a closed system of ordinary differential equations for the mass aspect, the gravitational and electromagnetic news function. With these at hand we are able to describe and interpreting the asymptotic and radiative properties of the spacetimes under consideration, in what their late time limit is regarded. We give expressions for the total charge, Bondi and NUT masses and the Newman-Penrose constants and with these at hand we succeed in attaching a physical meaning to the main parameters defining the solutions. We also give some additional arguments which provide an independent backup to our interpretation. Then, we carry out a counting of the degrees of freedom of both the gravitational and electromagnetic field. Finally, we outline our main conclusions.
## 2 Constraints
### 2.1 Asymptotic flatness and Bondi coordinates.
To begin with, consider a foliation of an asymptotically flat spacetime by —say— future oriented null hypersurfaces. Let this foliation be parametrized by a retarded time coordinate $`u`$. On each null hypersurface one takes a generic geodesic generator as a null curve parametrized by an affine parameter $`r`$. At a particular cut of $`I`$ (the intersection of the null hypersurfaces with null infinity), coordinates $`x^\alpha `$ are introduced, where $`\alpha =2,3`$. These will be taken to be the usual angular coordinates $`(\theta ,\phi )`$, although a complex stereographic coordinate $`\zeta `$ could be used as well. In addition, these coordinates can be propagated along $`I`$ and into the interior of the spacetime by requiring them to remain constant along the generators of $`I`$ and the outgoing null geodesics respectively. This coordinate construction will be referred as to *Bondi coordinates*.
In order to construct a null tetrad, one takes as its first vector $`l^a`$ the vector tangent the null hypersurfaces giving the foliation. The scaling in the affine parameter $`r`$ is then chosen so that the null vector coincides with $`_r`$. Now, by looking at the 2-surfaces $`u=const.`$, $`r=const.`$ ($`S_{u,r}`$) it can be seen that, at each point on these, there is another null vector $`n^a`$ which is future pointing and orthogonal to the 2-surface. Then, $`n^a`$ will be chosen as the second vector of the tetrad. Finally, the vectors $`m^a`$ and $`\overline{m}^a`$ are chosen so that they span the tangent space of $`S_{u,r}`$. From the previous construction it can be seen that:
$`l^a=_r,`$ (1)
$`n^a=_u+Q_r+C^\alpha _\alpha ,`$ (2)
$`m^a=\xi ^a_\alpha ,`$ (3)
with $`\alpha =\theta `$, $`\phi `$. The remaining freedom in the tetrad construction consists on a boost ($`lAl`$, $`nA^1n`$), which can be used to rescale $`r`$; and a spin ($`me^{i\theta }m`$), which in turn can be used in to make the spin coefficient $`ϵ`$ vanish. It can be also shown that $`\kappa =0`$, $`\tau =\overline{\pi }=\overline{\alpha }+\beta `$, and $`\rho `$ and $`\mu `$ are real.
The spacetimes under consideration will be assumed to be asymptotically flat with smooth sections of $`I`$, which is equivalent to say that the Weyl tensor of the spacetime peels off. Similarly, we will assume that the electromagnetic field also peels off. Thus,
$`\mathrm{\Psi }_n=𝒪(r^{n5}),`$ (4)
$`\varphi _m=𝒪(r^{m3}),`$ (5)
with $`n=0\mathrm{}4`$, $`m=0\mathrm{}2`$. The asymptotic expansions in powers of $`1/r`$ for peeling-off asymptotically flat Einstein-Maxwell systems are well known, and can be found in several places in the literature . Here Stewart’s version will be used <sup>1</sup><sup>1</sup>1It has to be pointed out that these expansions differ slightly from those appearing in Penrose & Rindler , as the tetrads used there are also different.. Explicitly, from the peeling-off theorem one has:
$`\mathrm{\Psi }_n=\mathrm{\Psi }_n^{n5}r^{n6}+\mathrm{\Psi }_n^{n6}r^{n6}+𝒪(r^{n7}),`$ (6)
$`\varphi _m=\varphi _m^{m3}r^{m3}+\varphi _m^{m4}r^{m4}+𝒪(r^{m5})`$ (7)
with $`n=0\mathrm{}4`$, $`m=0\mathrm{}2`$. It also follows $`\sigma =\sigma _2r^2+𝒪(r^3)`$. In general, the coefficients in the expansions (6)-(7) will depend on $`(u,\theta ,\phi )`$. The coefficients in $`\mathrm{\Psi }_0`$ ($`\mathrm{\Psi }_0^5`$, $`\mathrm{\Psi }_0^6,\mathrm{}`$) are the initial data quantities over a null hypersurface $`𝒩_0`$ in a characteristic initial value problem, whereas $`\mathrm{\Psi }_1^4`$ and $`\text{Re}\mathrm{\Psi }_2^3`$ (the Coulomb part of the gravitational field) are data that have to be supplied at $`𝒩_0I^+`$. A similar thing happens with the Maxwell field: $`\varphi _0`$ ($`\varphi _0^3`$, $`\varphi _0^4,\mathrm{}`$) being data at $`𝒩_0`$ and $`\varphi _1^2`$ (the Coulomb part of the Maxwell field) being data at $`𝒩_0I^+`$. The remaining data is set by supplying $`\varphi _2^1`$ and the leading term of $`\sigma `$. The coefficient $`\sigma _2`$ determines the radiative part of the gravitational field via:
$`\mathrm{\Psi }_4^1=\ddot{\overline{\sigma }}_2,`$ (8)
$`\mathrm{\Psi }_3^2=ð\dot{\overline{\sigma }}_2.`$ (9)
The imaginary part of $`\mathrm{\Psi }_2^3`$ is related to $`\sigma _2`$ through:
$$2\text{Im}(\mathrm{\Psi }_2^3)=\overline{ð}^2\sigma _2ð^2\overline{\sigma }_2+\overline{\sigma }_2\dot{\sigma }_2\dot{\overline{\sigma }}_2\sigma _2.$$
(10)
In addition, the following relations will be required later:
$`\mathrm{\Psi }_1^5=6\overline{\varphi }_1^2\varphi _0^3\overline{ð}\mathrm{\Psi }_0^5,`$ (11)
$`\mathrm{\Psi }_2^4=4\varphi _1^2\overline{\varphi }_1^2\overline{ð}\mathrm{\Psi }_1^4,`$ (12)
$`\mathrm{\Psi }_3^3=2\overline{\varphi }_1^2\varphi _2^1\overline{ð}\mathrm{\Psi }_2^3,`$ (13)
$`\mathrm{\Psi }_4^2=\overline{ð}\mathrm{\Psi }_3^2,`$ (14)
together with,
$`\varphi _1^3=\overline{ð}\varphi _0^3,`$ (15)
$`\varphi _2^2=\overline{ð}\varphi _1^2.`$ (16)
These relation are obtained from the expansions of the $`D`$-Bianchi identities and Maxwell equations. Similarly, from the $`\mathrm{\Delta }`$-Bianchi identities one obtains the evolution equations (i.e. the equations with derivatives with respect to $`u`$):
$`\dot{\mathrm{\Psi }}_0^5=ð\mathrm{\Psi }_1^4+3\sigma _2\mathrm{\Psi }_2^3+6\varphi _0^3\overline{\varphi }_2^1,`$ (17)
$`\dot{\mathrm{\Psi }}_0^6=4\overline{ð}(\sigma _2\mathrm{\Psi }_1^4)\overline{ð}ð\mathrm{\Psi }_0^5+8\overline{\varphi }_1^2ð\varphi _0^316\sigma _2\varphi _1^2\overline{\varphi }_1^2+8\varphi _0^4\overline{\varphi }_2^2,`$ (18)
$`\dot{\mathrm{\Psi }}_1^4=ð\mathrm{\Psi }_2^3+2\sigma _2\mathrm{\Psi }_3^2+4\varphi _1^2\overline{\varphi }_2^1,`$ (19)
$`\dot{\mathrm{\Psi }}_2^3=ð\mathrm{\Psi }_3^2+\sigma _2\mathrm{\Psi }_4^1+2\varphi _2^1\overline{\varphi }_2^1,`$ (20)
and
$`\dot{\varphi }_0^3=ð\varphi _1^2+\sigma _2\varphi _2^1,`$ (21)
$`\dot{\varphi }_0^4=\overline{ð}ð\varphi _0^32\overline{ð}(\sigma _2\varphi _1^2),`$ (22)
$`\dot{\varphi }_1^2=ð\varphi _2^1.`$ (23)
### 2.2 The algebraic type of the gravitational and electromagnetic fields
As remarked in the introduction, we make extra assumptions are required in order to obtain a closed system of ordinary differential equations. Our analysis will be restricted to type D spacetimes (i.e. spacetimes with two pairs of repeated null principal directions). A type D spacetime is obtained by demanding the quartic
$$\mathrm{\Psi }_04c_w\mathrm{\Psi }_1+6c_w^2\mathrm{\Psi }_24c_w^3\mathrm{\Psi }_3+c_w^4\mathrm{\Psi }_4=0,$$
(24)
to have two different pairs of repeated roots. This is the case if:
$`\mathrm{\Psi }_0=(\mathrm{\Psi }_4)^3\left[\mathrm{\hspace{0.17em}3}\mathrm{\Psi }_2\mathrm{\Psi }_42(\mathrm{\Psi }_3)^2\right]^2,`$ (25)
$`\mathrm{\Psi }_1=\mathrm{\Psi }_3(\mathrm{\Psi }_4)^2\left[3\mathrm{\Psi }_2\mathrm{\Psi }_42(\mathrm{\Psi }_3)^2\right],`$ (26)
provided that $`(\mathrm{\Psi }_3)^2\mathrm{\Psi }_2\mathrm{\Psi }_4`$, otherwise one gets a type N spacetime. From these two relations one obtains expressions for the coefficients in the expansions of $`\mathrm{\Psi }_0`$ and $`\mathrm{\Psi }_1`$ (i.e. $`\mathrm{\Psi }_0^5`$, $`\mathrm{\Psi }_0^6`$, $`\mathrm{\Psi }_1^4`$). The relations are:
$`\mathrm{\Psi }_0^5=`$ $`(\mathrm{\Psi }_4^1)^3\left[\mathrm{\hspace{0.17em}3}\mathrm{\Psi }_2^3\mathrm{\Psi }_4^12(\mathrm{\Psi }_3^2)^2\right]^2,`$ (27)
$`\mathrm{\Psi }_0^6=`$ $`(\mathrm{\Psi }_4^1)^4[2(\mathrm{\Psi }_3^2)^23\mathrm{\Psi }_4^1\mathrm{\Psi }_2^3)(3\mathrm{\Psi }_4^2(\mathrm{\Psi }_4^1\mathrm{\Psi }_2^32(\mathrm{\Psi }_3^2)^2)6\mathrm{\Psi }_2^4(\mathrm{\Psi }_4^1)^2+8\mathrm{\Psi }_4^1\mathrm{\Psi }_3^2\mathrm{\Psi }_3^3],`$ (28)
$`\mathrm{\Psi }_1^4=`$ $`(\mathrm{\Psi }_4^1)^2\mathrm{\Psi }_3^2\left[3\mathrm{\Psi }_2^3\mathrm{\Psi }_4^12(\mathrm{\Psi }_3^2)^2\right](\mathrm{\Psi }_4^1)^2.`$ (29)
The class of electromagnetic fields to be considered are such that their principal null directions are aligned with the two pairs of repeated principal null directions of the Weyl tensor. This requirement can be implemented as follows. The condition for a null rotation (about $`n^a`$) to give $`l^a`$ pointing in a principal null direction is given by equation (24). If the spacetime is of type D (i.e. equations (25) and (26) hold), then the double repeated solutions of this quartic equation are:
$$c_w=\frac{\mathrm{\Psi }_3\pm \sqrt{3\left((\mathrm{\Psi }_3)^2\mathrm{\Psi }_2\mathrm{\Psi }_4\right)}}{\mathrm{\Psi }_4}.$$
(30)
The analogous condition for the Maxwell field is given in terms of a quadratic equation,
$$\varphi _02c_m\varphi _1+c_m^2\varphi _2=0,$$
(31)
the roots of which are given by
$$c_m=\frac{\varphi _1\pm \sqrt{(\varphi _1)^2\varphi _0\varphi _2}}{\varphi _2}.$$
(32)
Demanding the roots of the quartic and quadratic equations to be equal by pairs, so that the principal null directions of the Maxwell field are aligned with the pairs of repeated principal null directions of the Weyl tensor, one arrives at the condition:
$`\varphi _0=\varphi _2\left(3\mathrm{\Psi }_2\mathrm{\Psi }_{4}^{}{}_{}{}^{1}2(\mathrm{\Psi }_3)^2(\mathrm{\Psi }_4)^2\right),`$ (33)
from which it can be seen that:
$`\varphi _0^3=\varphi _2^1\left(3\mathrm{\Psi }_2^3\mathrm{\Psi }_{4}^{1}{}_{}{}^{1}2(\mathrm{\Psi }_3^2)^2(\mathrm{\Psi }_4^1)^2\right),`$ (34)
$`\varphi _0^4=2(\mathrm{\Psi }_4^1)^3\varphi _3^2\left(\varphi _3^2\left(2\varphi _4^2\mathrm{\Psi }_2^1\varphi _4^1\mathrm{\Psi }_2^2\right)2\varphi _3^3\varphi _4^1\mathrm{\Psi }_2^1\right)+3\varphi _4^1\left(\varphi _2^4\varphi _4^1\mathrm{\Psi }_2^1\varphi _2^3\left(\left(\varphi _4^2\mathrm{\Psi }_2^1\right)\varphi _4^1\mathrm{\Psi }_2^2\right)\right),`$ (35)
### 2.3 The Killing vector fields
The spacetimes under consideration will be assumed to have two Killing vector fields ($`\xi _1^a`$ and $`\xi _2^a`$). The Killing vector $`\xi _1^a`$ will be taken to be an axial vector field (closed orbits). The Bondi coordinates can be chosen so that $`\xi _1^a=_\phi `$. Because of the closed orbits of $`\xi _1^a`$, the $`G_2`$ will be necessarily Abelian .
It was shown elsewhere that the asymptotic form of the most general Killing vector field compatible with asymptotic flatness is given by:
$$\xi ^a=(a_1u+𝒪(1/r),a_1r+𝒪(1),\frac{1}{\sqrt{2}}(\overline{c}_1+c_1)+𝒪(1/r),\frac{\text{i}\mathrm{csc}\theta }{\sqrt{2}}(\overline{c}_1c_1)+𝒪(1/r)),$$
(36)
where
$`a_1={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=1}{\overset{1}{}}}\{\overline{A}_m+(1)^mA_m\}(_0Y_{1,m}),`$ (37)
$`c_1={\displaystyle \underset{m=1}{\overset{1}{}}}(1)^{m+1}A_m\left({}_{1}{}^{}Y_{1,m}^{}\right),`$ (38)
$`\overline{c}_1={\displaystyle \underset{m=1}{\overset{1}{}}}\overline{A}_m\left({}_{1}{}^{}Y_{1,m}^{}\right),`$ (39)
and the $`A_m`$ are arbitrary complex numbers. In the case of axial symmetry, one has
$`a_1=0,`$ (40)
$`c_1={\displaystyle \frac{\text{i}\mathrm{sin}\theta }{\sqrt{2}}}.`$ (41)
The other Killing vector compatible with a radiative gravitational field is well known to be the *boost-rotation Killing vector* . Asymptotically, it reads:
$$\xi _2^\mu =(u\mathrm{cos}\theta ,r\mathrm{cos}\theta ,\mathrm{sin}\theta ,0).$$
(42)
Hence, in this case:
$`c_1=\mathrm{sin}\theta ,`$ (43)
$`a_1=\sqrt{2}\mathrm{cos}\theta .`$ (44)
Using these two Killing vectors, constraints on the diverse quantities of physical interest can be deduced . The relevant results for us are the following:
$`\sigma _2=\mathrm{csc}(\theta )\left(K_1(w)+iK_2(w)\right),`$ (45)
$`\text{Re}\mathrm{\Psi }_2^3={\displaystyle \frac{L(w)}{u^3}}{\displaystyle \frac{1}{2u}}\left(wK_1^{\prime \prime }(w)+2K_1^{}(w)\right),`$ (46)
$`\varphi _2^1={\displaystyle \frac{1}{u^2}}\left(H_1(w)+iH_2(w)\right),`$ (47)
$`\varphi _1^2={\displaystyle \frac{1}{u^2}}\left(G_1(w)+iG_2(w)\right){\displaystyle \frac{1}{\sqrt{2}u}}\mathrm{cot}\theta \left(H_1(w)+iH_2(w)\right),`$ (48)
where $`w=\mathrm{sin}\theta /u`$, and $`K_1`$, $`K_2`$, $`L`$, $`H_1`$, $`H_2`$, $`G_1`$, $`G_2`$ are arbitrary functions of the argument. Regularity at the poles requires $`G_1=G_2=0`$ .
## 3 The equations for $`K_1`$, $`K_2`$, $`L`$, $`H_1`$ and $`H_2`$
As mentioned previously, we will only consider peeling gravitational and electromagnetic fields, as polyhomogeneous boost-rotation symmetric fields happen to be singular at the “North” and “South” poles . Combining the evolution equations for the different leading coefficients of the Maxwell and Weyl fields together with the conditions of Petrov type D, the axial symmetry and the boost-rotation symmetry one arrives to a system of 5 *ordinary* differential equations for the functions $`K_1`$, $`K_2`$, $`L`$, $`H_1`$ and $`H_2`$. These differential equations are highly coupled and non-linear, and far too lengthy to be displayed here. In order to get around the difficulty to extract some relevant physical information out of them, we adopt the approach of solving the equations using formal expansions in $`w=\mathrm{sin}\theta /u`$. One might expect these expansions to hold for $`w`$ close to $`0`$, i.e. for very late times ($`u>>1`$).
### 3.1 The simplest case: the C-metric
How far one can go if one tries when trying an exact approach to the equations for $`K_1`$, $`K_2`$, $`L`$, $`H_1`$, $`H_2`$? In order to answer this, let us consider the easiest possible case in our analysis, that of a type D, boost-rotation symmetric spacetime with hypersurface orthogonal Killing vectors and no electromagnetic field. Under these assumptions one obtains the following system of ordinary differential equations for $`K=K_1`$ and $`L`$ which describe the news and mass aspect of the C-metric:
$$3K^{\prime \prime }+4KK^{}+2wKK^{\prime \prime }+wK^{\prime \prime \prime }2w^2L^{}6wL=0$$
(49)
and
$`32(K^{})^38w^2(K^{})^2K^{\prime \prime \prime }+16w^2(K^{})^2L+16w^4KK^{}K^{\prime \prime }L72w^4K^{}L^22w^4K^{\prime \prime }K^{\prime \prime \prime }+`$
$`24wK^2K^{\prime \prime }+4w^4K^{\prime \prime }L+10w^5K^{\prime \prime }L^{}48w^5K^{}LL^{}+4w^5KK^{\prime \prime }L+40w^3(K^{})^2L^{}+`$
$`40w^4K^{}K^{\prime \prime }L^{}+16w^3K^{}K^{\prime \prime }L8w^3K^{}K^{\prime \prime }K^{\prime \prime \prime }+16w^3K(K^{})^2L2w^3(K^{\prime \prime })^3`$
$`24w^6K^{\prime \prime }LL^{}+12w^6K^{\prime \prime \prime }L=0.`$ (50)
These equations can be decoupled yielding a quadrature for $`L`$ and the following fourth order differential equation for $`K`$:
$`10w^5K^{}K^{\prime \prime }K^{\prime \prime \prime }10w^5KK^{\prime \prime }K^{(iv)}+14w^5K(K^{\prime \prime \prime })^2+60w^4K^{}(K^{\prime \prime })^2+48w^4KK^{\prime \prime }K^{\prime \prime \prime }+`$
$`5w^4K^{\prime \prime }K^{(iv)}7w^4(K^{\prime \prime \prime })^220w^4KK^{}K^{(iv)}+20w^4(K^{})^2K^{\prime \prime \prime }+144w^3K(K^{\prime \prime })^219w^3K^{\prime \prime }K^{\prime \prime \prime }+`$
$`180w^3(K^{})^2K^{\prime \prime }+10w^3K^{}K^{(iv)}72w^3KK^{}K^{\prime \prime \prime }42w^2(K^{\prime \prime })^2+48w^2KK^{}K^{\prime \prime }+120w^2(K^{})^3+`$
$`46w^2K^{}K^{\prime \prime \prime }+66wK^{}K^{\prime \prime }+24wK(K^{})^2+48(K^{})^2=0,`$ (51)
which can be shown to have 3 three Lie point symmetries , and therefore can be reduced to the following Abel ordinary differential equation:
$$y^{}=\frac{2}{5}(224x^4+160x^3+22x^2x)y^3+\frac{1}{5}(4x+7)y^2\frac{60x+8}{20x^2+5x}y.$$
(52)
To the best of our knowledge this equation does not fall in any of the known solvable categories of Abelian equations . It is indeed somehow frustrating having gone so far, and not being able to perform the last step in the program! One may require a qualitative study of equation (51). This as well, is not an easy task due to the highly nonlinear character of the fourth order differential equation, which makes it very difficult to analyze the behavior of the phase space close to critical points.
### 3.2 The series solutions
As discussed previously, we will attempt to extract some physical information of the spacetimes using series expansions for $`|w|<<1`$, which describe the behavior of the system for late retarded times. It is always possible to choose the size of this region so that it will be completely contained in the radiative region of the spacetime (the region above the “roof”). The use of series carries some natural several drawbacks. In particular, it will not be possible to make any statement regarding the long term behavior and existence of the spacetime. The exceedingly complicated structure of the differential equations makes things worse, not allowing us to make estimates of the convergence radius or obtain the form of the general term in the expansion. Moreover, the coordinates $`(u,\theta )`$ are not well behaved at $`u=0`$ (see the discussion in ).
If the regularity of the the solutions at the poles is demanded then the leading behavior of the series is given by:
$`K_1=a_0+𝒪(w),`$ (53)
$`K_2=b_0+𝒪(w),`$ (54)
$`H_1=c_0+𝒪(w),`$ (55)
$`H_2=d_0+𝒪(w),`$ (56)
$`L=e_0+𝒪(w),`$ (57)
where $`a_i`$, $`b_i`$, $`c_i`$, $`d_i`$ and $`e_i`$ ($`i0`$) are constants. The constant $`a_0`$ can be removed using a super-translation, so that in fact $`K_1=𝒪(w)`$. The asymptotic shear transforms under super-translations according to :
$$\stackrel{~}{\sigma }_{2,0}(\stackrel{~}{u},x^\alpha )=\sigma _{2,0}(\stackrel{~}{u}\alpha ,x^\alpha )+ð^2\alpha (\theta ,\varphi ),$$
(58)
and this being so if initially one has $`\sigma _2=a_0\mathrm{csc}\theta +𝒪(w)`$, then in order to get $`\stackrel{~}{\sigma }_2=𝒪(w)`$ it is necessary to find a function $`\alpha (\theta ,\phi )`$ such that $`ð^2\alpha =a_0\mathrm{csc}\theta `$. This is given by,
$$\alpha (\theta ,\phi )=2a_0\mathrm{sin}\theta ,$$
(59)
that despite its apparent simplicity defines a true super-translation<sup>2</sup><sup>2</sup>2This transformation is not just a translation because the spherical harmonics expansion of the sine function consists of an infinite number of terms of the form $`Y_{2n,0}`$, $`n=0,1,\mathrm{}.`$. This choice of cuts of $`I^+`$ leads to the following series solutions:
$`K_1(w)`$ $`=a_3w^3+a_5w^5+𝒪(w^7),`$ (60)
$`K_2(w)`$ $`=b_3w^3+b_5w^5+𝒪(w^7),`$ (61)
$`H_1(w)`$ $`=c_1w+{\displaystyle \frac{3}{10\left(a_{3}^{}{}_{}{}^{2}+b_{3}^{}{}_{}{}^{2}\right)}}(5(b_5(b_3c_1+a_3d_1)+a_5(a_3c_1b_3d_1))`$ (62)
$`(a_3c_1b_3d_1)(c_{1}^{}{}_{}{}^{2}+d_{1}^{}{}_{}{}^{2}))w^3+𝒪(w^4),`$
$`H_2(w)`$ $`=wd_1+{\displaystyle \frac{3}{10\left(a_{3}^{}{}_{}{}^{2}+b_{3}^{}{}_{}{}^{2}\right)}}(5(a_5(b_3c_1+a_3d_1)b_5(a_3c_1b_3d_1))`$ (63)
$`(b_3c_1+a_3d_1)(c_{1}^{}{}_{}{}^{2}+d_{1}^{}{}_{}{}^{2}))w^3+𝒪(w^4),`$
$`L(w)`$ $`=4a_3+{\displaystyle \frac{2}{5}}\left(30a_5(c_{1}^{}{}_{}{}^{2}+d_{1}^{}{}_{}{}^{2})\right)w^2+𝒪(w^4).`$ (64)
The coefficients $`a_3`$, $`b_3`$, $`a_5`$, $`b_5`$ are free parameters of the gravitational field, whereas $`c_1`$, $`d_1`$ play the same role for the electromagnetic field . In order to shorten the expressions of the components of the Weyl and Maxwell spinors, the following complex parameters will be introduced:
$`q_1=c_1+id_1,`$ (65)
$`p_3=a_3+ib_3,`$ (66)
$`p_5=a_5+ib_5.`$ (67)
With this new notation, the leading terms of the components of the Weyl spinor are found to be
$`\mathrm{\Psi }_4^1=`$ $`6\mathrm{sin}^2\theta \left(2\overline{p}_3u^5+5\overline{p}_5\mathrm{sin}^2\theta u^7\right)+𝒪(u^9),`$ (68)
$`\mathrm{\Psi }_3^2=`$ $`3\sqrt{2}\mathrm{cos}\theta \mathrm{sin}\theta \left(2\overline{p}_3u^4+5\overline{p}_5\mathrm{sin}^2\theta u^6\right)+𝒪(u^8),`$ (69)
$`\mathrm{\Psi }_2^3=`$ $`2\overline{p}_3\left(23\mathrm{sin}^2\theta \right)u^3+\mathrm{sin}^2\theta \left(3\overline{p}_5\left(45\mathrm{sin}^2\theta \right){\displaystyle \frac{2q_1\overline{q}_1}{5}}\right)u^5+𝒪(u^7),`$ (70)
$`\mathrm{\Psi }_1^4=`$ $`3\sqrt{2}\mathrm{cos}\theta \mathrm{sin}\theta \left(\overline{p}_3u^2+{\displaystyle \frac{1}{20}}\left(5\overline{p}_5\left(15\mathrm{cos}2\theta \right)+4q_1\overline{q}_1\right)u^4\right)+𝒪(u^6),`$ (71)
$`\mathrm{\Psi }_0^5=`$ $`3\mathrm{sin}^2\theta \left(\overline{p}_3u^1{\displaystyle \frac{1}{20}}\left(5\overline{p}_5\left(3+5\mathrm{cos}2\theta \right)8q_1\overline{q}_1\right)u^3\right)+\text{O}(u^5),`$ (72)
$`\mathrm{\Psi }_0^6=`$ $`\mathrm{sin}^2\theta \left({\displaystyle \frac{15}{2}}\overline{p}_3{\displaystyle \frac{3}{8}}\left(5\overline{p}_5\left(5+7\mathrm{cos}2\theta \right)8q_1\overline{q}_1\right)u^2\right)+\text{O}(u^4);`$ (73)
whereas the relevant terms of the components of the Maxwell spinor are
$`\varphi _2^1=`$ $`q_1\mathrm{sin}\theta \left(u^3+{\displaystyle \frac{3}{10\overline{p}_3}}\mathrm{sin}^2\theta \left(5p_5q_1\overline{q}_1\right)u^5\right)+𝒪(u^7),`$ (74)
$`\varphi _1^2=`$ $`{\displaystyle \frac{q_1\mathrm{cos}\theta }{\sqrt{2}}}\left(u^2{\displaystyle \frac{3}{10\overline{p}_3}}\mathrm{sin}^2\theta \left(5\overline{p}_5q_1\overline{q}_1\right)u^4\right)+𝒪(u^6),`$ (75)
$`\varphi _0^3=`$ $`{\displaystyle \frac{q_1\mathrm{sin}\theta }{2}}\left(u^1{\displaystyle \frac{1}{10\overline{p}_3}}\left(23\mathrm{sin}^2\theta \right)\left(5\overline{p}_5q_1\overline{q}_1\right)u^3\right)+𝒪(u^5),`$ (76)
$`\varphi _0^4=`$ $`{\displaystyle \frac{3q_1\mathrm{sin}\theta }{4}}\left(1{\displaystyle \frac{1}{10\overline{p}_3}}\left(45\mathrm{sin}^2\theta \right)\left(5\overline{p}_5q_1\overline{q}_1\right)u^2\right)+𝒪(u^4).`$ (77)
It is worth noting the existence of a peeling off property of the $`u`$ dependence of these coefficients:
$`\mathrm{\Psi }_n^{5n}=\text{O}(u^{1n}),`$ (78)
$`\varphi _m^{3m}=\text{O}(u^{1m}),`$ (79)
with $`n=0,\mathrm{},4`$ and $`m=0,\mathrm{},2`$. This behavior can be understood in the following way: one can perform a construction equivalent to the one we presented, but for past-oriented light cones, parametrized by an advanced time $`v`$ coordinate. The components of the Weyl and Maxwell spinors will peel off with respect to the affine parameters of the null generators of these past-oriented cones (recall that boost-rotation symmetric spacetimes are time symmetric!). Now, close to $`I^+`$, one can show that $`r_v=u+\text{O}(u^1)`$, where $`r_v`$ is the affine parameter in the past-oriented system. Similarly, $`vr`$ in a neighbourhood of $`I`$, and $`l^a`$ and $`n^a`$ are “almost” parallel to $`n_v^a`$ and $`l_v^a`$ respectively. Thus, the leading coefficients should peel off with respect to $`r_v`$. Whence, they also have to peel off with $`u`$.
## 4 Asymptotic and radiative properties
We have already gathered enough information so as to study some the spacetime’s physical properties at $`I^+`$ in the late-time regime.
### 4.1 The electromagnetic charge
The total charge in the spacetime is given by
$`e={\displaystyle _{S^2}}\varphi _1^2(_0Y_{0,0})\text{d}S=\sqrt{\pi }{\displaystyle _{S^2}}\varphi _1^2\mathrm{sin}\theta \text{d}\theta =O(u^5),`$ (80)
i.e. it vanishes up to $`O(u^5)`$. Moreover, since the total electromagnetic charge in the spacetime is a conserved quantity, it should not contain any $`u`$ dependence. Therefore one can affirm that for our spacetimes
$$e=0.$$
(81)
This is in agreement with the interpretation put forward by Cornish & Utteley , who interpreted the charged C-metric as two charged black holes of opposite charges in hyperbolic motion.
### 4.2 The Bondi mass and the NUT mass
Of great physical relevance is the study of the mass loss of the system due to radiative processes. This can be done by evaluating the *Bondi mass*, which is given by:
$$M_B=\frac{1}{2}_{S^2}\text{Re }(\mathrm{\Psi }_2^3+\sigma _2\dot{\overline{\sigma }}_2)(_0,Y_{0,0})\text{d}S,$$
(82)
and can be shown to be non-increasing. Another related quantity, the NUT mass will also be of some interest due to the generic existence of nodal singularities in the interior of boost-rotation symmetric spacetimes . The NUT mass is essentially an imaginary version of the Bondi mass:
$$M_{NUT}=\frac{1}{2}_{S^2}\text{Im }(\mathrm{\Psi }_2^3+\sigma _2\dot{\overline{\sigma }}_2)(_0,Y_{0,0})\text{d}S,$$
(83)
The evaluation of these integrals for our solutions yields:
$`M_B=4\sqrt{\pi }\left({\displaystyle \frac{1}{15}}q_1\overline{q}_1u^5+{\displaystyle \frac{2}{5}}p_3\overline{p}_3u^7{\displaystyle \frac{16}{35}}\left(p_5\overline{p}_3+p_3\overline{p}_5\right)u^9+{\displaystyle \frac{32}{63}}p_5\overline{p}_5u^{11}\right)+𝒪(u^{13})`$ (84)
$`M_{NUT}={\displaystyle \frac{16\sqrt{\pi }}{35}}\left(p_3\overline{p}_5\overline{p}_3p_5\right)u^9+\text{O}(u^{13}).`$ (85)
If the electromagnetic field is absent then,
$$M_B=\frac{8\sqrt{\pi }}{5}p_3\overline{p}_3u^7+\text{O}(u^9),$$
(86)
thus $`m=p_3\overline{p}_3=|p_3|^2`$ can be interpreted as the mass of the system at a fiduciary retarded time (say $`u=1`$). Note as well that if $`u\mathrm{}`$ then $`M_B0`$, i.e. all the mass in the spacetime is carried away. This is consistent with the standard interpretation of boost-rotation symmetric as spacetimes describing particles in hyperbolic motion. These particles eventually pierce null infinity, thus leaving the (unphysical) spacetime. Consistent with these results, Dray has shown that the ADM mass of the C-metric is zero. His analysis shows that the reason for this is the particular kind of isometries this spacetime has. Therefore, it is very likely that this result also holds for the whole class of boost-rotation-symmetric spacetimes.
A similar discussion can be done with the NUT mass. The constant $`g=p_3\overline{p}_5\overline{p}_3p_5`$ can be interpreted as the measure of the “strength” of the nodal singularities joining the particles in hyperbolic motion at a fiduciary time, $`u=1`$.
### 4.3 The Newman-Penrose constants
Further insight on the physical meaning of the parameters $`q_1`$ and $`p_3`$ can be achieved by looking at the Newman-Penrose constants of the gravitational and electromagnetic field.
As shown by Exton et al. , there will be two sets of conserved quantities for the Einstein-Maxwell system under consideration: one for the electromagnetic field only, and one for the combined electromagnetic-gravitational system. The electromagnetic NP constants are given by:
$$_m=_{S^2}\varphi _0^4(_1\overline{Y}_{1,m})\text{d}S.$$
(87)
Note that due to the axial symmetry, all the constants, except for the one corresponding to $`m=0`$, will vanish identically (the $`m=0`$ spherical harmonic has no $`\phi `$ dependence).
The NP constants for the combined gravitational-electromagnetic system are somewhat more complicated as they involve the inverse operator $`ð^1`$:
$$𝒢_m=_{S^2}\{\mathrm{\Psi }_0^6+4\varphi _0^1\overline{ð}^1(\overline{\varphi }_1^2\overline{E})4\overline{E}\overline{ð}^1(\varphi _0^4F)\}(_2\overline{Y}_{2,m})\text{d}S,$$
(88)
where
$`E=e(_0Y_{0,0}),`$ (89)
$`F={\displaystyle \underset{m=1}{\overset{1}{}}}_m(_1Y_{1,m}),`$ (90)
and
$`\overline{ð}^1(\overline{\varphi }_1^2\overline{E})={\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \frac{(_1Y_{l,m})}{\sqrt{l(l+1)}}}{\displaystyle }\overline{\varphi }_1^2(_0\overline{Y}_{l,m})\text{d}S,`$ (91)
$`\overline{ð}^1(\overline{\varphi }_0^4\overline{F})={\displaystyle \underset{l=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}{\displaystyle \frac{(_2Y_{l,m})}{\sqrt{(l1)(l+2)}}}{\displaystyle }\overline{\varphi }_0^4(_1\overline{Y}_{l,m})\text{d}S.`$ (92)
The spacetimes under discussion are axially symmetric, and thus they only have one nonzero complex electromagnetic and one complex gravitational NP constants (those corresponding to the $`m=0`$ spherical harmonics). In the region below the “roof”, the spacetime is stationary, and therefore its electromagnetic Newman-Penrose constants will be of the form (mass monopole)$`\times `$(electric dipole)$``$(electric charge)$`\times `$(mass dipole), while the gravitational constants will be of the form (mass dipole)$`{}_{}{}^{2}`$(mass monopole)$`\times `$(mass quadrupole) .
Using equation (87) one finds that the electromagnetic NP constant is given by:
$$_0=\sqrt{\frac{3}{8\pi }}q_1.$$
(93)
Therefore $`q_1`$ is a product of a mass monopole and electric dipole, for there is no electromagnetic charge in the spacetime. Similarly, using equation (88) the evaluation of the gravitational NP constants yield:
$$𝒢_0=6\sqrt{\frac{5}{\pi }}\overline{p}_3.$$
(94)
Thus, $`p_3`$ has to be interpreted as a quantity of quadrupolar nature, even if strictly speaking the spacetime does not has a mass monopolar moment for the particle undergoing hyperbolic motion in a boost-rotation symmetric spacetime does not always are present in pairs.
Furthermore, let us recall the formulae for radiated power for the electromagnetic and the linearized gravitational field, the so called *dipole and Einstein’s quadrupole formulae* :
$`I_{elect}\ddot{\text{d}}^2,`$ (95)
$`I_{grav}\stackrel{2}{\stackrel{\mathrm{}}{\text{D}}},`$ (96)
where d is the dipolar moment of the charge distribution giving rise to the radiation field, and D is the quadrupolar moment of the mass distribution (a 3-dimensional tensor). ¿From the Bondi mass one can obtain by differentiation the flux of energy through null infinity. This is,
$$\dot{M}_B=4\sqrt{\pi }\left(\frac{1}{3}q_1\overline{q}_1u^6+\frac{14}{5}p_3\overline{p}_3u^8\right)+\text{O}(u^{10}).$$
(97)
Therefore the spacetimes under discussion radiate according to the dipole and Einstein’s quadrupole formula up to the leading terms. Higher order corrections are due to the damping of the electromagnetic field by the gravitational field, some other nonlinear effects like gravitational wave backscattering. Note as well how the flux of energy due to the gravitational field decays much faster than the flux due to the electromagnetic field. Thus at late times the main contribution to mass loss is electromagnetic in origin (Poynting vector).
### 4.4 Polarization states of the electromagnetic and gravitational waves
The discussion along this paper has focused on peeling electromagnetic and gravitational fields, and consequently the fields at null infinity are well defined. If one wishes to study the states of polarization of the electromagnetic or the gravitational fields, it is only necessary to consider the behavior of the type N part of the fields, that is, the leading coefficients of $`\varphi _2`$ and $`\mathrm{\Psi }_4`$ (i.e. $`\varphi _2^1`$ and $`\mathrm{\Psi }_4^1`$). The spinorial electromagnetic field is related to the components of the Maxwell field tensor via:
$`E_1iB_1=\varphi _0\varphi _2`$ (98)
$`E_2iB_2=i(\varphi _0+\varphi _2)`$ (99)
$`E_3iB_3=2\varphi _1`$ (100)
In the case here the electromagnetic field can be seen to have two polarization states, associated with the real part and imaginary parts of $`\varphi _2^1`$ respectively. One way to see this is to realize that if $`d_1`$ is set to zero then there will still be non vanishing components of the electromagnetic field. A similar phenomenon happens if $`c_1`$ is set to zero. This is nothing but a consequence of the existence of two different polarization states. If $`\varphi _2^1`$ is real, then the electric and magnetic fields —see equations (98)-(100)— lie along the $`e_1`$ and $`e_2`$ directions respectively. Note that $`e_1`$ and $`e_2`$ are orthonormal vectors transverse to the null generators of outgoing light cones, explicitly:
$`e_2^a={\displaystyle \frac{1}{\sqrt{2}}}i(m^a\overline{m}^a),`$ (101)
$`e_3^a={\displaystyle \frac{1}{\sqrt{2}}}(l^an^a).`$ (102)
Both 3-vectors are orthogonal, and have the same norm, as it should be expected from a plane wave. Let us call this polarization state $`P_x`$ (the polarization vector lies on the $`e_1`$ direction —x axis— ). When $`\varphi _2^1`$ is purely imaginary ($`P_y`$ polarization), the situation gets reversed and the electric field and magnetic lie along the $`e_2`$ and $`e_1`$ directions respectively (see figure 6.3).
By calculating the Poynting vectors for each state it can be checked that each configuration carries energy independently from the other, supporting again the assertion made on the existence of two independent degrees of freedom for the electromagnetic field.
For the gravitational field, the situation is fairly similar to that of the electromagnetic field. Again, there are two set of parameters: $`\{a_3,a_5\}`$ (the real parts of $`p_3`$ and $`p_5`$), and $`\{b_3,b_5\}`$ (the imaginary parts of $`p_3`$ and $`p_5`$). Both of them play an equivalent role in the expression for the Bondi mass. Hence both states of the gravitational field carry energy independently, i.e. they are true degrees of freedom (polarization states) of the gravitational field. The one associated to the real part of the Weyl tensor, which is given by the $`a`$-parameters, describes a $`(+)`$ polarization state. This can be seen by considering the geodesic deviation equation :
$$\delta \ddot{x}^a=\frac{1}{2}\text{Re }\mathrm{\Psi }_4^1\left(e_2^ae_{2c}e_1^ae_{1c}\right)\delta x^c,$$
(103)
therefore there will be tidal forces along the $`e_1`$ and $`e_2`$ directions. Similarly, if one considers the polarization state associated to the imaginary part of the Weyl tensor, one can perform a spin boost so that $`\mathrm{\Psi }_4`$ becomes real, so that one can use again equation (103). This polarization state will be a $`(\times )`$ one, with tidal forces rotated $`\pi /4`$ with respect to the $`e_1`$ and $`e_2`$ directions.
Note that if the Killing vectors are hypersurface orthogonal, then there is only one polarization state of the gravitational field ($`+`$). By contrast, the electromagnetic field can always have two polarizations states, as long as $`q_1`$ is complex (a magnetic dipole moment is then present!).
### 4.5 Conclusions
The main objective of this paper was to construct spacetimes that could be used as examples and test bench of different techniques and ideas used in the description of on asymptotic and radiative properties of the gravitational and electromagnetic fields.
Boost-rotation symmetric spacetimes are regarded as the best suited candidates for such a test bench since they contain the only examples known up to date of asymptotically flat radiative exact solutions. Perhaps the most remarkable result was that of showing how the Newman-Penrose constants allow one to consider, at least up to quadrupoles, multipolar moments of non-stationary spacetimes without having to resort to the linearized theory (cfr. the results on the mass loss). As seen, the complexity of the equations involved precludes the obtention of nice closed expressions; however, the asymptotic expressions for late times have allowed to extract most of physical properties of interest; this regardless of the natural limitations that this approach carries. |
warning/0003/cond-mat0003463.html | ar5iv | text | # Tunneling Spectroscopy of the Underdoped High-𝑇_𝑐 Superconductors
## I Introduction
The suppression of density of states around the Fermi level in the underdoped high-$`T_c`$ superconductors (HTSC) above $`T_c`$ is observed in various normal state experiments such as optical conductivity, dc-resistivity, angle resolved photoemission, NMR, tunneling spectroscopy, neutron scattering, specific heat, Raman spectroscopy etc , and this phenomena is termed as pseudogap (PG). Currently a consensus of the origin and nature of this PG is still lacking and undoubtedly it is a key issue to be resolved to make any progress toward a theory of high Tc superconductivity.
About the origin of the PG state only two possibilities are logically allowed. The first possibility is that the PG is somehow related with the superconducting correlation and it develops into a real superconducting gap below Tc. And the second possibility is that the PG has nothing to do with a superconducting gap but with something else. Along the first line of thinking, preformed pair scenario , pairing fluctuations scenario , etc. are proposed. For the second possibility, antiferromagnetic correlation, charge stripes, etc are considered as the origin of the PG. At present experimental evidences exist both for superconducting origin and for non-superconducting origin. Therefore it is necessary to design an experiment to identify distinct features among the proposed scenarios. Recently we proposed an experimental test using tunneling spectroscopy in the PG state, specifically, for the preformed pair scenario. Namely we claim that there should be an Andreev reflection signal even above Tc but below $`T^{}`$ (PG cross-over temperature) if there exist preformed Cooper pairs without long range phase coherence. Until now there are only one positive and one negative experiment reported on the observation of an Andreev signal in the PG region of the underdoped HTSC compounds.
In this paper we examine the second possibility for the PG state, i.e., that the PG is irrelevant from the SC gap. We calculate the tunneling conductance ($`dI/dV`$) at zero temperature when the PG coexists with a SC gap. Specifically we considered two cases: (1) the PG is a simple suppression of density of states of an unknown origin and the SC correlation has no direct interplay with the PG; (2) the PG is caused by an AFM correlation (approximated by SDW) and below Tc the AFM correlation and the SC correlation coexist and interplay with each other. For both cases, we use the BTK theory to calculate the tunneling conductance ($`dI/dV`$). The results show characteristic features of the tunneling density of states in each case and those features can be used to sort out the true origin of the PG in comparison with experiments.
The main results are: (1) for the first case, the basic line shape of the tunneling conductance is a simple superposition of a standard BTK conductance and an assumed PG density of states. When the SC gap size is smaller than the PG size, the SC gap feature shows up as a distinguishable peak inside the PG. However if the size of SC gap is larger than the PG, the PG feature is overwhelmed by the SC gap feature; (2) for the second case, there is an interesting interplay between the AFM and SC correlations. Irrespective of the relative sizes of the PG and SC gap, the tunneling conductance shows only one gap feature at $`E=\mathrm{\Delta }_{total}`$ ($`=\sqrt{\mathrm{\Delta }_{SC}^2+\mathrm{\Delta }_{PG}^2}`$). However depending on the relative sizes of the PG and SC gap, the line shape of $`dI/dV`$ looks very different. When the SC gap is bigger than the PG, it looks more like a conventional NIS (Normal metal-Insulator-Superconductor) junction. But for the other case the main feature of the $`dI/dV`$ curve is determined by the SDW correlation and shows no diverging density of states approaching the gap energy in contrast to the NIS junction. This difference comes from the fact that the tunneling density of states of the NISDW (Normal metal-Insulator-SDW state) junction obtained by the BTK theory (even at large barrier potential ($`Z`$) limit) is not the same as the actual density of states, which would have been obtained by the tunneling Hamiltonian method. The difference of the tunneling conductance by the BTK theory and by the tunneling Hamiltonian method in the SDW state, which is in contrast to the NIS (Normal metal-Insulator-Superconducor) junction, shows that the description of tunneling process by these two methods are not the same. The important and interesting questions are then which description is more proper description for actual tunneling process and why in the case of NIS junction those two methods seem to give consistent results for a $`Z>>1`$ limit. More details will be discussed in the later section.
## II PG model I: simple suppression of density of states
In this section we consider the case that the PG is a simple suppression of density of states, which has been developed above $`T_c`$ by an unknown origin. An assumption is that the SC gap develops below $`T_c`$ on top of the already developed PG density of states and there is no other correlation between the SC gap and the PG.
The tunneling conductance of NIS junction is calculated by the BTK theory as follows.
$`dI(eV)/dV`$ $`=`$ $`2ev_FA{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{𝑑𝐸}N_{\mathrm{𝑃𝐺}}(E){\displaystyle \frac{f(EV)}{V}}`$ (2)
$`[1+|A(E)|^2|R(E)|^2],`$
where
$$|A|^2=\{\begin{array}{cc}\frac{\mathrm{\Delta }_S^2}{(E+(1+2Z^2)\sqrt{E^2\mathrm{\Delta }_S^2})^2}\hfill & \text{for }E>\mathrm{\Delta }_S\hfill \\ \frac{\mathrm{\Delta }_S^2}{E^2+(1+2Z^2)^2(\mathrm{\Delta }_S^2E^2)}\hfill & \text{for }E<\mathrm{\Delta }_S\hfill \end{array}$$
(3)
$$|R|^2=\{\begin{array}{cc}\frac{4Z^2(E^2\mathrm{\Delta }_S^2)(Z^2+1)}{(E+(1+2Z^2)\sqrt{E^2\mathrm{\Delta }_S^2})^2}\hfill & \text{for }E>\mathrm{\Delta }_S\hfill \\ \frac{4Z^2(Z^2+1)(\mathrm{\Delta }_S^2E^2)}{E^2+(1+2Z^2)^2(\mathrm{\Delta }_S^2E^2)}\hfill & \text{for }E<\mathrm{\Delta }_S\hfill \end{array}$$
(4)
are the coefficients for the Andreev and normal reflections, respectively, and $`Z`$ is the strength of the insulating barrier ($`Z=\frac{mV}{\mathrm{}^2k_F}`$). The above equation is a standard formula for the tunneling conductance by the BTK theory with one modification, i.e., multiplied by the PG density of states $`N_{PG}(E)`$. We simulate $`N_{PG}(E)`$ by Dynes’ formula with the SC gap replaced by the PG ($`\mathrm{\Delta }_P`$).
$$N_{PG}(E)=2\pi N(0)Re\left[\frac{\omega +i\mathrm{\Gamma }}{\sqrt{(\omega +i\mathrm{\Gamma })^2\mathrm{\Delta }_P^2}}\right].$$
(5)
The numerical calculation has done for a tunneling into (1,0,0) direction toward the HTSC and assume that $`\mathrm{\Delta }_P`$ and $`\mathrm{\Delta }_S`$ has the maximum value at the same direction (1,0,0). To simulate the surface roughness we also present the angle averaged results by simply replacing $`\mathrm{\Delta }_{S,P}\mathrm{\Delta }_{S,P}^0\mathrm{cos}(2\theta )`$. All the presented tunneling conductance is normalized by the normal-state resistance $`R_N=(1+Z^2)/(2N(0)e^2v_FA)`$ and the bias voltage is measured in unit of $`\mathrm{\Delta }_P`$. Also for all the presented result we chose $`\mathrm{\Gamma }=0.3`$, which determines the shape of the PG density of states $`N_{PG}(E)`$ from Eq(4).
In Fig.1(a) we show the normalized tunneling conductance $`R_NdI/dV`$ at zero temperature with the superconducting gap $`\mathrm{\Delta }_S=0.5\mathrm{\Delta }_P`$ for different barrier potentials (Z=0,0.5, and 1). We also plot $`N_{PG}(E)`$ with $`\mathrm{\Gamma }=0.3`$ for comparison. The results can be trivially understood. As well known from the BTK theory, for small value of Z, we see the enhanced conductance below $`\mathrm{\Delta }_S`$ due to the Andreev reflection and this effect quickly disappears with increasing Z value. The main difference of our results from the conventional BTK calculation is that we modulate the conventional BTK tunneling conductance by multiplying with $`N_{PG}(E)`$. As a result the enhanced conductance below $`\mathrm{\Delta }_S`$ and the sharp peak structure at $`E=\mathrm{\Delta }_S`$ appear inside the PG density of states ($`\mathrm{\Delta }_S<\mathrm{\Delta }_P`$). In Fig.1(b) we show the same calculations but angle averaged to simulated the surface roughness in real tunneling experiment. The main features are the same as Fig.1(a) and the line shapes become more rounded off. In particular the result of Z=1 case in Fig.1(b) looks quite similar with the recent tunneling experiment by Krasnov et al.
In Fig.1(c) and Fig.1(d) we show the results of same calculations as Fig.1(a) and Fig.1(b), respectively but with $`\mathrm{\Delta }_S=\mathrm{\Delta }_P`$. When $`\mathrm{\Delta }_S`$ has the same value as $`\mathrm{\Delta }_P`$ the line shape of the tunneling conductance looks similar to the conventional NIS junction, in particular with a large Z value. The effect of the PG is just to enhance the over-all line shape of the conductance. For smaller Z value the conductance is still enhanced below $`\mathrm{\Delta }_S`$ due to the Andreev reflection. However even the enhanced conductance at $`E0`$ limit is far below than 2 (for a conventional NIS junction it approaches 2 for $`E<\mathrm{\Delta }_S`$ as Z $``$ 0) because of the reduced DOS of PG origin. Again the angle averaged results in Fig.1(d) shows a more rounded off line shapes. The angle averaged PG density of states (PDOS) is shown for comparison in Fig.1(b) and (d).
In Fig2.(a-d), we plot the normalized conductance with varying $`\mathrm{\Delta }_S(=`$0.5,1, and 1.5 $`\mathrm{\Delta }_P`$). When Z=1 (Fig.2(a) and Fig.2(b)) the results are easily understood. For $`\mathrm{\Delta }_S<\mathrm{\Delta }_P`$ the distinct peak structure due to SC appears inside the PG as explained above. And when $`\mathrm{\Delta }_S>\mathrm{\Delta }_P`$ there is no more distinct peak structure and the SC feature overwhelms the PG structure. When Z=0.1 (Fig.2(c) and Fig.2(d)) the conductance line shapes look more peculiar, but basically it can be understood as an overlap of the Andreev enhanced conductance below $`\mathrm{\Delta }_S`$ and the PG density of states $`N_{PG}(E)`$. G. Deutscher et al. reported that there are two energy scales observed in tunneling experiments, and depending on the barrier strength ($`Z`$) only one of the two energy scales dominates the conductance line shape. In Fig2.(c) the line with $`\mathrm{\Delta }_S=0.5\mathrm{\Delta }_P`$ (solid line) shows a similar feature to the experimental data in Fig.2 of Ref..
In view of the experiments of Krasnov et al. and G. Deutscher et al. , if our calculations have any relevance with underdoped HTSC, it would be the case of $`\mathrm{\Delta }_S<\mathrm{\Delta }_P`$ (Fig.1(a)(b) and Fig.2(c)(d)) and it will be interesting if more tunneling experiments with different Z parameters become available in near future.
## III PG model II:SDW
In this section we consider another possibility of the PG, namely, in which the correlation, which induces the PG, coexists with a SC correlation below $`T_c`$ and the PG correlation interplay with the SC correlation. In this case the key difference from the case (I) is that either $`\mathrm{\Delta }_S`$ or $`\mathrm{\Delta }_P`$ does not show up as separate features in the tunneling conductance but the total gap $`\mathrm{\Delta }_{total}=\sqrt{\mathrm{\Delta }_S^2+\mathrm{\Delta }_P^2}`$ shows up as a result of the interplay of two gaps. However more detailed line shape of the tunneling conductance reveals the existence of two gaps, $`\mathrm{\Delta }_S`$ and $`\mathrm{\Delta }_P`$.
Specifically, we assume the SDW correlation as the origin of PG. The reason for this assumption is of course that there is a strong antiferromagnetic (AFM) correlation in underdoped HTSC compounds. Particularly for the underdoped HTSC compounds there is clear experimental evidence from neutron scattering for the coexistence of the AFM correlation and superconductivity below $`T_c`$ and also the AFM correlation length becomes much larger than the SC correlation length. This enable us to treat the short range AFM correlation by SDW from the point of view of the SC correlation.
### A tunneling with SDW only
As a prelude we first generalize the BTK theory to the NISDW (Normal metal-Insulator-SDW) junction. For simplicity we consider only a one dimension model and assume a commensurate SDW state with $`Q=k_F=\pi /a`$ ($`a`$ is the lattice distance) and also neglect the Fermi surface mismatch between the normal metal and the SDW state.
The wave functions of the lefthand (normal metal) and the righthand (SDW) sides of the tunneling barrier is written as,
$`\mathrm{\Psi }_L(x)`$ $`=`$ $`e^{ikx}+Re^{ikx}`$ (6)
$`\mathrm{\Psi }_R(x)`$ $`=`$ $`T[a(k)e^{ikx}+b(k)e^{i(kQ)x}]`$ (7)
where $`a(k)^2=\frac{ϵ+\xi }{2ϵ}`$ and $`b(k)^2=\frac{ϵ\xi }{2ϵ}`$ are the Bogoliubov coefficients of the SDW state, $`ϵ=\sqrt{\xi ^2+\mathrm{\Delta }_{SDW}^2}`$, and $`\xi =\mathrm{}^2k^2/2m\mu `$. By matching the boundary conditions of the wave functions at $`x=0`$, we obtained $`R`$ and $`T`$ as follows.
$`R(E)`$ $`=`$ $`{\displaystyle \frac{(a+b)iZa}{(a+b)iZ+b}}`$ (8)
$`T(E)`$ $`=`$ $`{\displaystyle \frac{1}{(a+b)iZ+b}}`$ (9)
For $`ϵ^2>\mathrm{\Delta }_{SDW}^2`$, $`|R|^2`$ is easily calculated, and for $`ϵ^2<\mathrm{\Delta }_{SDW}^2`$, using $`a^2=\frac{ϵi\sqrt{\mathrm{\Delta }_{SDW}^2ϵ^2}}{2ϵ}`$ and $`b^2=\frac{ϵ+i\sqrt{\mathrm{\Delta }_{SDW}^2ϵ^2}}{2ϵ}`$, we can show $`|R|^2=1`$.
The tunneling conductance of NISDW junction with different Z (=0,0.3, and 1) are shown in Fig.3(a,b,c) as dash-dot lines ($`\mathrm{\Delta }_S=0`$ case). The tunneling conductance line shape is qualitatively different from the actual density of states $`N_{SDW}(E)`$, which would have been obtained by the tunneling Hamiltonian method. Here we clearly observe that the BTK theory and the tunneling Hamiltonian method give different results for tunneling conductance. And it has been already known that these two methods describe different physical processes for the tunneling phenomena . In particular, the tunneling Hamiltonian method, although physically appealing, has never been justified in regard with the tunneling transfer matrix element $`T`$ (it is always taken as a constant). Then another question may arise: why then for the BCS state does the BTK theory give a qualitatively similar result as the tunneling Hamiltonian methods at least for a large $`Z`$ limit where the Andreev reflection process is suppressed, but not for the SDW state ? The reason is that for the BCS state the quasiparticle is a superposition of a momentum $`k`$ particle and a momentum $`k`$ hole, and all together it carries the flux of the same momentum $`k`$, while for the SDW state the quasiparticle is a superposition of momenta $`k`$ and $`k+Q`$ particles and therefore the quasiparticle in SDW state doesn’t carry a single momentum. For the BTK theory of tunneling the main physics is the flux conservation described by Liouville’s theorem, therefore the BTK theory sensitively traces the correlation of different momentum states in the SC or SDW states, while the tunneling Hamiltonian method is completely blind of the momentum correlation. In the case of SC (BCS) state, we are in a fortune situation as described above, therefore momentum correlation doesn’t play an important role but particle-hole correlation plays an important role in the BTK theory, which provides the main difference between the BTK theory and the tunneling Hamiltonian method through the Andreev scattering. Now for the case of SDW state, the momentum correlation is important but not the particle-hole correlation; therefore it is expected that the two methods will give different results.
An important question is then which theory should be trusted. As described above, the tunneling Hamiltonian method has never been justified, while there is no logical flow in the BTK theory based on the flux conservation. Therefore we think unless the tunneling interface is very rough (of course we need to estimate how much rough is rough) the BTK theory becomes more trustable.
### B tunneling with SDW+SC
Now we will consider the tunneling junction of a normal metal-insulator-SDW+SC (NISDWS). Again for simplicity we consider only a one dimensional model. In reality, as in HTSC, the analysis of tunneling in the two dimensional SDW state is more complicated, in particular when SDW+SC is considered. However our one dimensional model is enough to provide a qualitative understanding of the interplay of SDW and SC correlations in tunneling process.
Once the SDW state is formed, there are two branches of quasiparticles created: $`\alpha _{+,k}=a(k)c_k+b(k)c_{k+Q}`$ and $`\alpha _{,k}=b(k)c_ka(k)c_{k+Q}`$, where $`a(k)`$ and $`b(k)`$ are the Bogoliubov coefficients as defined above. Now we assume that the superconducting pairing occurs between $`\alpha _{+,k}`$ and $`\alpha _{+,k}`$, and also between $`\alpha _{,k}`$ and $`\alpha _{,k}`$, but not between $`\alpha _{+,k}`$ and $`\alpha _{,k}`$, for example. Although this assumption is mainly for simplicity of the analysis, it is known that the pairing interaction between different branches are much weaker in Hubbard model. Also strictly speaking the SDW state with a commensurate wave vector ($`Q=\pi /a`$) is an insulator with a fully developed gap below $`\mathrm{\Delta }_{SDW}=\mathrm{\Delta }_P`$, in contrast to the PG state in HTSC where residual density of states still remain and the system remains a metal. Therefore it is clear that our SDW state only mimic a state with a short range AFM correlation of the underdoped HTSC and it is justified by the fact that the AFM correlation is effectively long range from the viewpoint of the SC correlation. A main drawback of this assumption is that the effect of the residual density of states below PG in the tunneling conductance is completely missing. In summary, the purpose of this section is to study a tunneling characteristics of a SDW+SC state due to the interplay of two correlations, but not all the details of real materials.
Now for the BTK theory the wave functions of the lefthand (normal metal) and righthand (SDW+SC) sides of the tunneling interface are written as,
$`\mathrm{\Psi }_L(E)`$ $`=`$ $`\left(\begin{array}{c}1\\ 0\end{array}\right)e^{ikx}+R(E)\left(\begin{array}{c}1\\ 0\end{array}\right)e^{ikx}+A(E)\left(\begin{array}{c}0\\ 1\end{array}\right)e^{ikx},`$ (16)
$`\mathrm{\Psi }_R(E)`$ $`=`$ $`C(E)\left(\begin{array}{c}u\\ v\end{array}\right)\alpha _+(k)+D(E)\left(\begin{array}{c}v\\ u\end{array}\right)\alpha _{}(k).`$ (21)
where $`E^2=\sqrt{ϵ^2+\mathrm{\Delta }_S^2}=\sqrt{\xi ^2+\mathrm{\Delta }_{SDW}^2+\mathrm{\Delta }_S^2}`$ and $`\xi =\mathrm{}^2k^2/2m\mu `$. All $`kk_F`$ approximation is taken; the error of the approximation is $`O(\mathrm{\Delta }_{total}/E_F)`$. Also a shorthand notation of the wave functions are
$`\alpha _+(k)`$ $`=`$ $`a(ϵ)e^{ikx}+b(ϵ)e^{i(kQ)x},`$ (22)
$`\alpha _{}(k)`$ $`=`$ $`b(ϵ)e^{ikx}a(ϵ)e^{i(kQ)x}.`$ (23)
As usual the SC Bogoliubov coefficients are $`u(E)^2=(E+ϵ)/2E`$ and $`u(E)^2=(Eϵ)/2E`$. By matching boundary conditions, which is a lengthy but straight forward calculation, we obtained the coefficients of $`A(E)`$ and $`R(E)`$ as follows.
$`A`$ $`=`$ $`{\displaystyle \frac{uv}{F}},`$ (24)
$`R`$ $`=`$ $`{\displaystyle \frac{(u^2v^2)[Z^2(b^2a^2)+iZ(b^2a^2+2ab)ab]}{F}},`$ (25)
$`F`$ $`=`$ $`(u^2v^2)[Z^2(b^2a^2)+iZ2ab]+b^2u^2+a^2v^2.`$ (26)
In order to calculate $`|A|^2`$ and $`|R|^2`$ we need to divide the region of $`E`$ into three regions. For region (I), where $`E>\mathrm{\Delta }_{total}(=\sqrt{\mathrm{\Delta }_P^2+\mathrm{\Delta }_S^2}`$), $`ϵ=\sqrt{E^2\mathrm{\Delta }_S^2}`$ and $`\xi =\sqrt{E^2\mathrm{\Delta }_{total}^2}`$. For region (II), where $`\mathrm{\Delta }_S<E<\mathrm{\Delta }_{total}`$, $`ϵ=\sqrt{E^2\mathrm{\Delta }_S^2}`$ and $`\xi =i\sqrt{\mathrm{\Delta }_{total}^2E^2}`$, and finally for region (III), where $`E<\mathrm{\Delta }_S`$, $`ϵ=i\sqrt{\mathrm{\Delta }_S^2E^2}`$ and $`\xi =i\sqrt{\mathrm{\Delta }_{total}^2E^2}`$. But in the final results of $`|A|`$ and $`|R|`$, $`ϵ`$ factors are all cancelled out, so that the region (II) and (III) are not distinguished. Therefore the tunneling characteristics shows distinct change only across $`E=\mathrm{\Delta }_{total}`$ but no change across either $`E=\mathrm{\Delta }_S`$ or $`E=\mathrm{\Delta }_P`$.
The final results of the Andreev ($`A`$) and normal reflection ($`R`$) coefficients are: for region (I),
$`|A|^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_S^2}{G}},`$ (27)
$`|R|^2`$ $`=`$ $`{\displaystyle \frac{(2Z^2\sqrt{E^2\mathrm{\Delta }_t^2}\mathrm{\Delta }_P)^2+4Z^2(\sqrt{E^2\mathrm{\Delta }_t^2}+\mathrm{\Delta }_P)^2}{G}},`$ (28)
$`G`$ $`=`$ $`[(2Z^2+1)\sqrt{E^2\mathrm{\Delta }_t^2}+E]^2+4Z^2\mathrm{\Delta }_P^2;`$ (29)
and for region (II) and (III),
$`|A|^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_S^2}{H}},`$ (30)
$`|R|^2`$ $`=`$ $`{\displaystyle \frac{(\mathrm{\Delta }_P+2Z\sqrt{\mathrm{\Delta }_t^2E^2})^2+4Z^2(Z\sqrt{\mathrm{\Delta }_t^2E^2}+\mathrm{\Delta }_P)^2}{H}},`$ (31)
$`H`$ $`=`$ $`E^2+[(2Z^2+1)\sqrt{\mathrm{\Delta }_t^2E^2}+2Z\mathrm{\Delta }_P]^2.`$ (32)
In Fig.3(a-c) we plot the numerical calculations of the normalized tunneling conductance as a function of bias voltage (in unit of $`\mathrm{\Delta }_P`$) with varying size of $`\mathrm{\Delta }_S`$ (=0,0.5,1, and 2) for different barrier potentials (Z=0, 0.3, and 1). The main features of the tunneling conductance are very different from a pure SC case. When $`\mathrm{\Delta }_S`$ is much bigger than $`\mathrm{\Delta }_P`$ (say, $`\mathrm{\Delta }_S/\mathrm{\Delta }_P=2`$ in Fig.3) the line shape is similar to the pure SC case. But even in this case the position of a gap in the conductance is determined by $`\mathrm{\Delta }_t`$ and the presence of SC gap $`\mathrm{\Delta }_S`$ only show up through an Andreev scattering coefficient, which enhances the conductance below the total gap. On the other hand when $`\mathrm{\Delta }_S\mathrm{\Delta }_P`$, the line shape looks closer to the pure SDW case, which is qualitatively different from a real density of states as explained in the previous section. Again the presence of SC gap $`\mathrm{\Delta }_S`$ show up only through an Andreev scattering and enhances the conductance below the total gap.
In Fig.4(a)(b) we plot the same calculations in different organization. Here we vary Z values with fixed $`\mathrm{\Delta }_S`$ to see more clearly the effect of varying strength of the tunneling barrier. Main features are already explained above. Both in Fig.4(a) and (b), the normalized conductance below $`\mathrm{\Delta }_{total}`$ is enhanced by the Andreev reflection and it increases as $`Z0`$. When Z=0, we can have a simple results of $`|A|^2=\frac{\mathrm{\Delta }_S^2}{\mathrm{\Delta }_t^2}`$ and $`|R|^2=\frac{\mathrm{\Delta }_P^2}{\mathrm{\Delta }_{total}^2}`$ for $`E<\mathrm{\Delta }_{total}`$. These expressions of $`|A|^2`$ and $`|R|^2`$ clearly show the contrasting role of the SDW gap (normal reflection) and the SC gap (Andreev reflection); the normalized conductance is given by $`R_NdI/dV(E<\mathrm{\Delta }_t)=\frac{2\mathrm{\Delta }_S^2}{\mathrm{\Delta }_t^2}`$ in this limit. Increasing Z from zero, it continues to suppress the Andreev reflection and enhances the normal reflection.
In view of current tunneling experiments, if our results of the NISDWS junction has any relevance with the underdoped HTSC, only possibility is the case of $`\mathrm{\Delta }_S>\mathrm{\Delta }_P`$ at low temperature. In order to realize this possibility starting with a sizable $`\mathrm{\Delta }_P`$ and a zero $`\mathrm{\Delta }_S`$ above $`T_c`$, we have to imagine that $`\mathrm{\Delta }_P`$ gradually decreases and $`\mathrm{\Delta }_S`$ gradually increases as temperature is lowered. Because this scenario is quite unlikely from the currently available experiments, we can conclude either that the PG is not caused by an AFM correlation, or that even if the origin of PG is an AFM correlation the AFM correlation is not strong enough to have any significant interplay with a SC correlation.
## IV Conclusion
We considered two phenomenological models of the PG state in underdoped HTSC and studied the characteristics of tunneling conductance at zero temperature by generalizing the BTK theory. In model I, we assumed that the PG is a simple suppression of density of states of a unknown origin and when the system goes to SC state there is no direct interplay between the PG correlation and the SC correlation. In this case the characteristics of tunneling conductance is a simple superposition of a standard tunneling conductance with SC gap $`\mathrm{\Delta }_S`$ and the PG density of states $`N_{PG}(E)`$. Despite the simplicity of the model, the results seem to explain some of the recent tunneling experiments by Krasnov et al. and G. Deutscher , which indicate non-superconducting origin of the PG. In the model II, we assumed that the PG is caused by an AFM correlation and it is simulated by SDW state. Below $`T_c`$ the SDW correlation and SC correlation show an interesting interplay. As a result, the tunneling gap is given by $`\mathrm{\Delta }_{total}=\sqrt{\mathrm{\Delta }_S^2+\mathrm{\Delta }_P^2}`$ and individual gaps, $`\mathrm{\Delta }_S`$ and $`\mathrm{\Delta }_P`$, do not show up explicitly. In particular when $`\mathrm{\Delta }_{P,SDW}>\mathrm{\Delta }_S`$, the line shape of the tunneling conductance looks qualitatively different from a conventional NIS junction. In view of available tunneling experiments in underdoped HTSC, the relevance of the model II is possible only when $`\mathrm{\Delta }_{P,SDW}<\mathrm{\Delta }_S`$ below $`T_c`$, which is quite unlikely at present. To clarify the issue of the PG more experiments on tunneling are essential and our study should serve as a useful benchmark.
We would like to thank H.J. Lee for invaluable discussion of his data and G. Deutscher for sending us their preprint prior to publication, respectively. This work was supported by the Korean Science and Engineering Foundation (KOSEF) through the Center for Strongly Correlated Materials Research (CSCMR) (2000)(YB) and through the Grant No. 1999-2-114-005-5 (YB and HYC). |
warning/0003/hep-ph0003236.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Two central issues in supersymmetric grand unified theories (SUSY-GUT’s) are a) mechanism of SUSY breaking (including the origin of the SUSY breaking scale) and mediation of SUSY breaking to the SM superpartners and b) mechanism of GUT symmetry breaking (down to the SM gauge group) and the origin of the GUT scale (denoted by $`M_{GUT}`$) $`2\times 10^{16}`$ GeV (the energy scale at which SM gauge couplings unify with the MSSM particle content).
In , we presented a model in which both these symmetries (SUSY and a simple GUT gauge group) are broken by the same modulus field (i.e., by the same scalar potential). A non-zero vev for the $`F`$-component of this field is generated dynamically breaking SUSY. The GUT scale which is the vev of the scalar ($`A`$-)component of the same field (at the minimum of the potential) is determined (dynamically) by an “inverted hierarchy” mechanism. Therefore the GUT scale is naturally both larger than the SUSY breaking scale (which is required for the vev of the scalar component to be calculable in perturbation theory) and smaller than the Planck scale. This is the first example of its kind in the literature. <sup>3</sup><sup>3</sup>3In also, SUSY and a GUT gauge group are broken by the same field. However, in , SUSY breaking is not dynamical so that a very small superpotential coupling is required to explain the smallness of the SUSY breaking scale compared to $`M_{Pl}`$ and also if all superpotential couplings are of the same order, then $`M_{GUT}M_{Pl}`$. In , the GUT gauge group is not simple so that gauge coupling unification is not a prediction of the model and also an assumption about a (non-calculable) Kähler potential is required for the model to work.
In this model, there are two comparable contributions to the MSSM scalar and gaugino masses – one is mediated by supergravity (SUGRA) and the other by gauge interactions. This is crucial to achieving a realistic sfermion mass spectrum since if we neglect the SUGRA contribution, then the (gauge mediated contribution to) the right-handed (RH) slepton (mass)<sup>2</sup> (at the weak scale) is negative. However, the SUGRA contribution to sfermion (mass)<sup>2</sup> is, in general, arbitrary in flavor space giving unacceptable rates for flavor changing neutral currents (FCNC’s). We will argue that, in general, this “problem” will be present in any model where SUSY and a GUT symmetry are broken by the same field.
In this paper, we show how the above model can be “improved”, i.e., how “flavor-conserving” and positive sfermion (mass)<sup>2</sup> can be attained using the recently proposed idea of gaugino mediated SUSY breaking .
## 2 Review of model and the problem
### 2.1 Model
We begin with a brief review of the model of . The gauge group is:
$$SU(6)_{GUT}\times SU(6)_S$$
(1)
and the field content is listed in Table 1.
The core part of the superpotential is:
$$W_1=\lambda _Q\mathrm{\Sigma }Q\overline{Q}+\frac{\lambda _\mathrm{\Sigma }}{3}\mathrm{\Sigma }^3.$$
(2)
Along the flat direction parametrized by tr $`\mathrm{\Sigma }^2`$, the vev of $`\mathrm{\Sigma }`$ is:
$$\mathrm{\Sigma }=\frac{v}{\sqrt{12}}\text{diag}[1,1,1,1,1,1]$$
(3)
which breaks $`SU(6)_{GUT}`$ to $`SU(3)\times SU(3)\times U(1)`$ at the scale $`v`$. We identify one unbroken $`SU(3)`$ with $`SU(3)_{color}`$ and show later how to break the other $`SU(3)\times U(1)`$ to $`SU(2)_{weak}\times U(1)_Y`$ at the same scale ($`v`$). Thus we identify the value of $`v`$ at the minimum of the potential (see later) with the GUT scale (i.e., the scale at which SM gauge couplings meet). The $`\mathrm{\Sigma }`$ vev gives mass to $`Q`$, $`\overline{Q}`$ so that below the scale $`v`$, the only massless fields are the flat direction $`v`$ (we will denote both the chiral superfield and its scalar component by $`v`$) and the pure gauge theory $`SU(6)_S`$ (the other components of $`\mathrm{\Sigma }`$ either get a mass through the $`\mathrm{\Sigma }^3`$ term or are eaten by the broken generators of $`SU(6)_{GUT}`$; see for details).
The pure $`SU(6)_S`$ gauge theory undergoes gaugino condensation and generates the superpotential:
$$W=6\mathrm{\Lambda }_L^3=\sqrt{3}\lambda _Q\mathrm{\Lambda }^2v,$$
(4)
where $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }_L`$ are the dynamical scales of the high and low energy $`SU(6)_S`$, respectively (they are related by one-loop matching at the scale $`v`$, assuming $`v\mathrm{\Lambda }`$). Thus, below the scale $`\mathrm{\Lambda }_L`$, the only massless field is $`v`$ with the above superpotential so that SUSY is broken since $`F_v\mathrm{\Lambda }^2`$. But, with the canonical (tree-level) Kähler potential ($`v^{}v`$), the vev of the scalar ($`A`$-)component of $`v`$ (i.e., the GUT scale) is undetermined since the scalar potential is flat $`\mathrm{\Lambda }^4`$. The dominant corrections to the Kähler potential (and hence the scalar potential) are given by the wavefunction renormalization of $`\mathrm{\Sigma }`$ (denoted by $`Z`$) so that:
$$V(v)\frac{\mathrm{\Lambda }^4}{Z(v)}.$$
(5)
Since $`v\mathrm{\Lambda }`$, we can compute $`Z(v)`$ using perturbation theory. In renormalization group (RG) evolution, at one-loop, $`Z(v)`$ receives contributions from the Yukawa coupling(s) ($`\lambda _{\mathrm{\Sigma },Q}`$) and the $`SU(6)_{GUT}`$ gauge coupling – the former (latter) tends to decrease (increase) $`Z(v)`$ as $`v`$ increases. Thus, if the gauge coupling dominates at small $`v`$ whereas the Yukawa coupling is larger at high scales (as is natural if $`SU(6)_{GUT}`$ is asymptotically free), then the potential can develop a minimum. Furthermore, the minimum can be (naturally) at a value of $`v\mathrm{\Lambda }`$ since $`Z`$ and both couplings depend logarithmically on $`v`$: at the minimum, we require that $`v10^{16}`$ GeV (to obtain the correct GUT scale) whereas $`\mathrm{\Lambda }10^{10}10^{11}`$ GeV so that MSSM sparticles have masses $`100\text{GeV}1`$ TeV (see later). <sup>4</sup><sup>4</sup>4This is only a local minimum since there is a supersymmetric minimum with $`\mathrm{\Sigma }Q\overline{Q}\mathrm{\Lambda }`$ . However since, at the local minimum $`v\mathrm{\Lambda }`$ the tunneling rate from the local minimum to the “true” vacuum is very small. Thus, this inverted hierarchy mechanism can generate a GUT scale, $`M_{GUT}`$, much larger than the SUSY breaking scale, $`\mathrm{\Lambda }`$ (which is also required for the perturbative calculation mentioned above to be valid) <sup>5</sup><sup>5</sup>5It should be possible to contruct models based on (say) $`SO(10)`$ along similar lines.. We can also view this as “generating” the GUT scale from the Planck scale ($`M_{Pl}2\times 10^{18}`$ GeV) as follows. We can choose $`Z=1`$ (canonical normalization) at $`M_{Pl}`$ and RG evolve $`Z`$ to lower energies. Since $`Z`$ and the gauge and Yukawa couplings vary logarithmically (“slowly”) with energy, (we can choose $`O(1)`$ couplings at $`M_{Pl}`$ such that) $`Z`$ reaches a maximum at an energy scale $`M_{GUT}`$ which is “much” smaller (i.e., by two orders of magnitude) than $`M_{Pl}`$ .
To break the other $`SU(3)\times U(1)`$ to $`SU(2)_{weak}\times U(1)_Y`$ (and to get the usual light Higgs doublets) we add:
$$W_2=\underset{i=1}{\overset{2}{}}S_i(H_i\overline{H}_i\mathrm{\Sigma }^2),$$
(6)
$$W_3=\underset{i=1}{\overset{2}{}}H_i(\mathrm{\Sigma }+X_i)\overline{h}_i+\underset{i=1}{\overset{2}{}}\overline{H}_i(\mathrm{\Sigma }+\overline{X}_i)h_i,$$
(7)
$$W_4=\frac{1}{M}\left(\left(H_1\overline{H}_1\right)\left(H_2\overline{H}_2\right)\left(H_1\overline{H}_2\right)\left(H_2\overline{H}_1\right)\right),$$
(8)
$$W_5=S_3\left(H_1\overline{H}_2H_2\overline{H}_1\right).$$
(9)
The role of these $`W`$’s is as follows (for details, see ). $`W_2`$ forces $`H`$, $`\overline{H}`$ to acquire vev’s. With $`H=\overline{H}v(1,0,0,0,0,0)`$, $`SU(3)\times U(1)`$ is broken to $`SU(2)\times U(1)`$. $`W_3`$ forces $`X`$, $`\overline{X}`$ to acquire vev’s such that only the triplets in $`H`$’s get a mass with those in $`\overline{h}`$’s (this is the “sliding singlet” mechanism). <sup>6</sup><sup>6</sup>6Any symmetries which allow the terms $`\mathrm{\Sigma }^3`$, $`S\mathrm{\Sigma }^2`$ and $`SH\overline{H}`$ (which we need to obtain the desired vev’s), also allow the term $`\mathrm{\Sigma }H\overline{H}`$ which spoils the above pattern of vev’s; in addition, there might be higher dimension operators (allowed by the same symmetries) which might spoil the sliding singlet mechanism and/or the pattern of vev’s. Thus, this model is only “technically” natural, i.e., the superpotential is not the most general one allowed by symmetries. A pair of doublets in $`H`$, $`\overline{H}`$ is eaten in the gauge symmetry breaking while a pair of doublets in $`h`$, $`\overline{h}`$ gets a mass with doublets in $`\mathrm{\Sigma }`$ (due to the $`H`$, $`\overline{H}`$ vev). This leaves two pairs of doublets massless: the one in $`H`$, $`\overline{H}`$ acquires mass through $`W_4`$ and the one in $`h`$, $`\overline{h}`$ is the usual pair of Higgs doublets. $`W_5`$ gives a required constraint between the vev’s of $`H_1`$, $`\overline{H}_1`$ and $`H_2`$, $`\overline{H}_2`$.
To complete the model, the SM quarks and leptons are obtained through:
$`W_6`$ $`=`$ $`N_i(\overline{P}_{1j}\overline{H}_1+\overline{P}_{2j}\overline{H}_2)+N_i(\overline{P}_{1j}\overline{h}_1+\overline{P}_{2j}\overline{h}_2)`$ (10)
$`+N_iN_jY+(X_1+X_2)Y\overline{Y}+\overline{Y}(H_1h_1H_2h_2).`$
For each generation, the vev’s of $`\overline{H}`$’s gives a mass to one combination of $`\overline{\mathrm{𝟓}}`$ (under $`SU(5)`$) in $`\overline{P}_{1,2}`$ with the $`\mathrm{𝟓}`$ (under $`SU(5)`$) of $`N`$ leaving one $`\overline{\mathrm{𝟓}}`$ in $`\overline{P}_{1,2}`$ and $`\mathrm{𝟏𝟎}`$ in $`N`$ massless – these are the quarks and leptons. The other terms in $`W_6`$ generate their Yukawa couplings.
### 2.2 Problem with sfermion spectrum
The MSSM scalars and gauginos acquire masses $`\left[\alpha _{GUT}/\pi \right]\left[F_v/M_{GUT}\right]10^2\left[F_v/M_{GUT}\right]`$ through gauge mediation (GM) by coupling at one or two-loops to two sources <sup>7</sup><sup>7</sup>7The “loop” factor for these masses is $`\alpha _{GUT}/(4\pi )`$, but there is usually an enhancement from group theory and/or large number of messengers which effectively makes the loop factor $`\alpha _{GUT}/\pi 10^2`$ (for $`\alpha _{GUT}0.04`$) – this estimate suffices for comparison to SUGRA mediated masses (see later). The precise expressions for the gauge mediated masses are given in Eqs. (11) and (12). : one is the “matter” messengers – the $`Q`$, $`\overline{Q}`$ fields and the heavy components of $`H`$, $`\overline{H}`$, $`h`$ and $`\overline{h}`$ fields which as usual have a SUSY breaking mass spectrum due to the coupling to $`v`$. <sup>8</sup><sup>8</sup>8With the addition of $`W_{2,3}`$, the flat direction $`v`$ is a combination of $`\mathrm{\Sigma }`$, $`H`$, $`\overline{H}`$, $`X`$ and $`\overline{X}`$. The other source, usually referred to as “gauge” messengers, is the heavy part of the $`SU(6)_{GUT}`$ gauge multiplet which also has a SUSY breaking spectrum since the field(s) breaking the GUT gauge group ($`\mathrm{\Sigma }`$, $`H`$ and $`\overline{H}`$) have a non-zero $`F`$-component.
Using the technique of , the MSSM gaugino masses are given by ($`A`$ denotes the gauge group):
$$M_A(\mu )\frac{\alpha _A(\mu )}{4\pi }\frac{F_v}{v}\left(b_Ab_6\right),$$
(11)
where $`b_A`$’s are the beta functions of the SM gauge groups below the GUT scale and $`b_6`$ is the beta function of the $`SU(6)_{GUT}`$ above the GUT scale. The gaugino masses are non-universal since the messengers (both gauge and matter) are not in complete GUT multiplets.
The matter messengers give a positive contribution to the scalar (mass)<sup>2</sup> as usual, but the gauge messenger contribution is typically negative (for all scalars) so that most scalars have negative (mass)<sup>2</sup> at the GUT scale. Of course, in RG scaling to the weak scale, the sfermion (mass)<sup>2</sup> get a positive contribution from the gaugino masses. The gauge mediated MSSM sfermion (other than stop) (mass)<sup>2</sup> at the scale $`\mu `$ are given by (again using the technique of ):
$`m_i^2(\mu )`$ $``$ $`{\displaystyle \frac{1}{16\pi ^2}}\left({\displaystyle \frac{F_v}{v}}\right)^2\times `$ (12)
$`\left({\displaystyle \underset{A}{}}{\displaystyle \frac{2C_A^i}{b_A}}\left(\alpha _A^2(\mu )\left(b_6b_A\right)^2b_6^2\alpha _6^2\right)+2C_6^ib_6\alpha _6^2\right),`$
where $`C_A^i`$ is the quadratic Casimir invariant for the scalar $`i`$ under the gauge group $`A`$, i.e., $`4/3`$, $`3/4`$ for fundamentals of $`SU(3)_c`$, $`SU(2)_L`$ respectively and $`3/5Y^2`$ for $`U(1)_Y`$. $`C_6^i=35/12`$ for fields in $`\overline{\mathrm{𝟓}}`$ of $`SU(5)`$ ($`\overline{\mathrm{𝟔}}`$ of $`SU(6)_{GUT}`$) and $`14/3`$ for fields in $`\mathrm{𝟏𝟎}`$ of $`SU(5)`$ ($`\mathrm{𝟏𝟓}`$ of $`SU(6)_{GUT}`$). The beta function for $`SU(N_c)`$ group is defined as $`3N_cN_{f,eff}`$, where $`N_{eff}`$ is the “effective” number of flavors. $`\alpha _6`$ is the $`SU(6)`$ coupling at the GUT scale. With the particle content in Table 1, we get $`b_6=11`$.
In this model, it turns out that the gauge mediated contribution to RH slepton (mass)<sup>2</sup>, Eq. (12), is negative at the weak scale, i.e., the positive bino mass contribution (which is not an independent parameter) is not enough to make the RH slepton (mass)<sup>2</sup> positive while all other sfermion (mass)<sup>2</sup> are positive due to the larger wino/gluino mass contribution.
So far, the SUGRA contribution to MSSM scalar and gaugino masses has been neglected. The SUGRA contribution to MSSM sfermion (mass)<sup>2</sup> from the operators:
$$d^4\theta c_1\frac{X^{}X\mathrm{\Phi }_i^{}\mathrm{\Phi }_j}{M_{Pl}^2},$$
(13)
where $`X`$ is any SUSY breaking field ($`\mathrm{\Sigma }`$, $`H_i`$, $`X_i`$ etc.) and $`\mathrm{\Phi }`$ is a MSSM matter field, is $`\left[F_v/M_{Pl}\right]^210^4\left[F_v/M_{GUT}\right]^2`$ ( assuming $`c_1O(1)`$ and $`M_{Pl}2\times 10^{18}`$ GeV). Thus, the SUGRA and gauge mediated contributions (to sfermion masses) are comparable so that the combined RH slepton (mass)<sup>2</sup> can still be positive (provided the SUGRA contribution is positive and a bit larger than the GM contribution) and a realistic spectrum can be achieved. However, the SUGRA contribution (Eq. (13)) violates flavor – in general, the off-diagonal terms in the sfermion (mass)<sup>2</sup> matrices (in flavor space) will be $`O(\left[F_v/M_{Pl}\right]^2)`$ which clearly result in too large SUSY contributions to FCNC’s. There is also a SUGRA contribution to MSSM gaugino masses ($`F_v/M_{Pl}`$) comparable to the GM contribution (see ($`4D`$ equivalent of) Eq. (14)) which, in turn, contributes to scalar masses in RG scaling to the weak scale; this contribution to sfermion (mass)<sup>2</sup> is positive and flavor-conserving.
From the above discussion, it is clear that generic models (i.e., not just the one in ) in which the same field breaks SUSY and a GUT symmetry will have the same problem – at the scale $`M_{GUT}`$, there will be a contribution (which is typically negative) to the scalar (mass)<sup>2</sup> from the heavy gauge multiplet so that, at least for the RH slepton, the gauge mediated (mass)<sup>2</sup> at the weak scale might be negative. Also, even if the gauge mediation contribution is positive, the SUGRA contribution (Eq. (13)) to the sfermion (mass)<sup>2</sup> is comparable which leads to flavor violation, in general. It is clear that the latter is a problem also in models of GM with messenger scale close to the GUT scale (having nothing to do with GUT symmetry breaking), for example the model of .
## 3 Improved model using an extra dimension
We now show how to improve the above model, i.e., how to obtain positive and (at the same time) flavor-conserving sfermion (mass)<sup>2</sup> using the framework of gaugino mediated SUSY breaking ($`\stackrel{~}{g}`$MSB) .
Consider the following embedding of this model in a $`5`$-dimensional ($`5D`$) theory. The MSSM matter fields (i.e., $`N`$ and $`\overline{P}_{1,2}`$) are localized on a “$`3`$-brane” (“MSSM matter” brane) whereas the $`SU(6)_{GUT}`$ gauge fields and $`Y`$, $`\overline{Y}`$, $`h`$, $`\overline{h}`$ fields propagate in the extra dimension. SUSY breaking fields, i.e., $`\mathrm{\Sigma }`$, $`H`$, $`\overline{H}`$, $`X`$ and $`\overline{X}`$ and $`Q`$, $`\overline{Q}`$ are localized on a different $`3`$-brane (“SUSY breaking” brane) which is separated from the MSSM matter brane by a distance $`R3M^1`$ in the extra dimension, where $`M`$ is the “fundamental” ($`5D`$) Planck scale. For simplicity we also assume that $`R`$ is the size of the extra dimension. The $`S_{1,2,3}`$ fields and the $`SU(6)_S`$ gauge fields can either propagate in the bulk or be localized on the SUSY breaking brane.
In the effective $`5D`$ field theory below $`M`$ ($`10^{18}`$ GeV, see later), direct couplings between fields on the matter brane and SUSY breaking brane are forbidden since such couplings are not local . Of course, such operators might be generated by integrating out bulk states with mass $`M`$, but such effects will suppressed by a Yukawa factor $`e^{RM}`$ due to the (position space) propagator of the bulk state . In particular, the coefficient of the operator in Eq. (13) is $`\stackrel{<}{}10^2`$. Thus, the SUGRA contributions to the MSSM squark and slepton (mass)<sup>2</sup> from contact Kähler terms are suppressed by this factor relative to other contributions (see later).
The superpotential couplings written in section 2.1 are all allowed, except for the $`N\overline{P}_{1,2}\overline{H}_{1,2}`$ coupling in Eq. (10) ($`N`$, $`\overline{P}_{1,2}`$ and $`\overline{H}`$ fields are localized on different branes). This coupling in the model of section 2.1 gives an $`O(M_{GUT})`$ mass to $`\mathrm{𝟓}`$ (under $`SU(5)`$) of $`N`$ with a $`\overline{\mathrm{𝟓}}`$ (under $`SU(5)`$) in $`\overline{P}_{1,2}`$ while the massless $`\overline{\mathrm{𝟓}}`$ in $`\overline{P}_{1,2}`$ and $`\mathrm{𝟏𝟎}`$ in $`N`$ are the usual quarks and lepton fields. Thus, in the above framework, a $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$ (for each generation) is massless in addition to the usual quarks and leptons. The beta-functions of SM gauge group ($`b_A`$’s) and hence the GM contribution to scalar and gaugino masses (in the $`4D`$ theory; see later) depends on whether these addtional fields are light or not. We assume that the extra $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$’s acquire mass $`O(M_{GUT})`$ by some mechanism. <sup>9</sup><sup>9</sup>9For example, the extra $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$’s might couple to additional fields or another possibility is that the operator $`d^2\theta \left(N\overline{P}\overline{H}X/M\right)`$ is generated with a coupling $`e^{RM}`$ by exchange of bulk (“string”) states – this operator will gives a mass $`O(e^3M_{GUT}/M)10^{12}`$ GeV to the extra $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$. In the latter case, there is an additonal messenger scale for GM $`10^{12}`$ GeV since the $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$ in $`N`$, $`\overline{P}_{1,2}`$ have a non-supersymmetric spectrum due to the coupling to $`\overline{H}`$, $`X`$ – in the $`4D`$ model of section 2.1 these messengers were at the scale $`M_{GUT}`$.
We require the (usual Higgs doublets in) $`h`$, $`\overline{h}`$ fields to couple to matter fields so that Yukawa couplings can be generated. These Higgs fields should also couple to the GUT symmetry breaking fields (and hence SUSY breaking fields, in this case), which are localized on a different brane, so that the usual Higgs doublet-triplet splitting can be achieved. Thus, Higgs fields $`h`$, $`\overline{h}`$ have to propagate in the bulk.
The MSSM gauginos get a mass from the following ($`5D`$) SUGRA interactions:
$$d^2\theta \left(c_2\frac{\mathrm{\Sigma }W_\alpha W^\alpha }{M^2}+c_3\frac{(X+\overline{X})W_\alpha W^\alpha }{M^2}\right)+h.c.,$$
(14)
where $`W_\alpha `$ is a $`5D`$ gauge field and $`c_{2,3}O(1)`$. The contribution to MSSM gaugino masses from the singlet is universal while that from $`\mathrm{\Sigma }`$ is non-universal: $`M_1:M_2:M_3=1:\mathrm{\hspace{0.33em}5}:5`$. When these operators and also the operators $`d^4\theta X^{}X/M^5W_\alpha D^2W^\alpha `$ (which generate a SUSY breaking gaugino wavefunction) are inserted in one-loop diagrams, we get contributions to sfermion (mass)<sup>2</sup> $`g_{5D}^2/(4\pi ^2)\left[F_v/M^2\right]^21/R^3`$ and $`g_{(5D)}^2/(4\pi ^2)F_v^2/M^5\mathrm{\hspace{0.33em}1}/R^4`$, respectively ($`g_{(5D)}`$ is the $`5D`$ gauge coupling) which are finite due to the spatial separation of the SUSY breaking and the MSSM matter branes in the $`5D`$ theory. <sup>10</sup><sup>10</sup>10In other words, these contributions to sfermion (mass)<sup>2</sup> are the effect of integrating out the extra dimension, i.e., the Kaluza-Klein excitations of the gauge fields. Using Eqs. (15) and (16) (see below), these one-loop contributions to sfermion (mass)<sup>2</sup> are of order $`\alpha _{(4D)}/\pi \left[F_v/M_{Pl}\right]^2\mathrm{\hspace{0.33em}1}/(MR)`$ and $`\alpha _{(4D)}/\pi \left[F_v/M_{Pl}\right]^2\mathrm{\hspace{0.33em}1}/(MR)^2`$, respectively and thus are negligible compared to GM contribution at $`M_{GUT}`$ and gaugino mass contribution in RG scaling to the weak scale (see below).
When the extra dimension is compactified (i.e., the Kaluza-Klein (KK) excitations of supergravity and other bulk states are integrated out), we get an effective $`4D`$ field theory below the scale $`R^11/3M`$. The $`4D`$ and $`5D`$ Planck scales are related by
$$M_{Pl}^2M^3R3M^2$$
(15)
so that $`MM_{Pl}/\sqrt{3}10^{18}`$ GeV and $`R^1M_{Pl}/54\times 10^{17}`$ GeV. <sup>11</sup><sup>11</sup>11For simplicity, we assume that any other extra dimensions have size $`M^1`$ so that $`M`$ is the fundamental quantum gravity (“string”) scale and the inverted hierarchy mechanism generates the hierarchy between $`M`$ and $`M_{GUT}`$ (à la the hierarchy between $`M_{Pl}`$ and $`M_{GUT}`$ in $`4D`$). It is easy to extend this framework to one with more than one extra dimensions of size slightly larger than the fundamental length scale. However, in that case, the fundamental (say, $`(4+n)D`$) Planck scale, $`M`$, might be smaller than $`10^{18}`$ GeV (since $`M_{Pl}^2M^{n+2}R^n`$ for $`n`$ extra dimensions of size $`R`$) so that the motivation for the inverted hierarchy mechanism (to explain $`M_{GUT}/M`$) is a bit weaker. In what follows, we assume $`MM_{Pl}`$. It was shown in that the exchange of supergravity Kaluza-Klein (KK) excitations also does not lead to contact Kähler terms of order $`1/M_{Pl}^2`$ (of the form Eq. (13)). The $`4D`$ and $`5D`$ gauge couplings are related by
$$g_{(4D)}^2\frac{g_{(5D)}^2}{R}.$$
(16)
The $`4D`$ (zero-mode) gaugino mass is given by $`c_{2,3}F_v/M^2\mathrm{\hspace{0.33em}1}/Rc_{2,3}F_v/(\sqrt{3}M_{Pl})`$, i.e., $`O(F_v/M_{Pl})`$ .
The anomaly mediated contribution to MSSM scalar and gaugino masses is
$`F_v/M_{Pl}\alpha /\pi `$ and thus can be neglected in comparison to SUGRA contribution to gaugino masses, GM contributions (at $`M_{GUT}`$) to gaugino and scalar masses and gaugino mass contribution to scalar masses in RG scaling to the weak scale (see below). There is also a contribution to scalar (mass)<sup>2</sup> at one-loop $`1/\left(16\pi ^2\right)F_v^2/(R^2M_{Pl}^4)`$ from integrating out SUGRA KK modes ; for $`MR3`$, this is comparable to the anomaly mediated contribution and thus can be neglected. <sup>12</sup><sup>12</sup>12 Both these contributions to sfermion masses and also the one-loop gaugino contribution in the $`5D`$ theory mentioned above are flavor-conserving.
At the scale $`M_{GUT}10^{16}`$ GeV <sup>13</sup><sup>13</sup>13As before, the GUT scale is determined by the one-loop corrections to the wavefunction ($`Z`$) of $`\mathrm{\Sigma }`$. Between the energy scales $`M`$ and $`R^1`$, the $`SU(6)_{GUT}`$ gauge coupling (rather $`\alpha _6^1`$) (and similarly $`Z`$, $`\lambda _{Q,\mathrm{\Sigma }}`$) “runs” with a power of energy since the gauge theory is $`5D`$ whereas below $`R^14\times 10^{17}`$ GeV, we have the usual ($`4D`$) RG scaling. In any case, the inverted hierarchy mechanism can result in a minimum of the potential at $`v10^{16}`$ GeV $`M`$; most of the RG scaling of $`Z`$, $`g_6`$ and $`\lambda _{\mathrm{\Sigma },Q}`$ from $`M`$ to $`10^{16}`$ GeV is the usual $`4D`$ evolution., the heavy gauge multiplet and the matter messengers are integrated out (as in section 2.2) generating the contributions to the scalar (mass)<sup>2</sup> (at two-loops) and gaugino masses (at one-loop), Eqs. (11) and (12). <sup>14</sup><sup>14</sup>14 In computing these “threshold” corrections, the SUGRA contribution $`F_v/M_{Pl}`$ to the SUSY breaking masses of the heavy gauge multiplet and the matter messengers ($`Q`$, $`\overline{Q}`$ etc.) can be neglected in comparison to the contribution $`F_v/M_{GUT}`$ from direct coupling to $`v`$. There are also contributions to sfermion and gaugino masses obtained by replacing the zero-mode ($`4D`$) gauge fields in the above loop diagrams by KK gauge fields which have masses $`kR^1`$. The one-loop diagrams with KK gauge fields and gauge messengers give (zero-mode) gaugino masses $`_k\alpha /\pi F_vv/\left(kR^1\right)^2`$ (since the $`R`$-symmetry breaking scale is $`v`$) and the two-loop diagrams with KK gauge fields (and either matter or gauge messengers) give sfermion (mass)$`{}_{}{}^{2}_k(\alpha /\pi )^2F_v^2/\left(kR^1\right)^2`$. <sup>15</sup><sup>15</sup>15Strictly speaking, since this is the effect of integrating out the KK gauge fields, these contributions to sfermion and gaugino masses appear at/above the scale $`R^1`$. Since $`_k\mathrm{\hspace{0.33em}1}/k^2`$ is convergent for the case of one extra dimension, we see that these contributions are suppressed by a factor $`1/\left(M_{GUT}R\right)^2O(100)`$ compared to Eqs. (11) and (12). As mentioned earlier, GM contributions to gaugino and scalar masses (at $`M_{GUT}`$) (Eqs. (11) and (12)) and the SUGRA contribution to gaugino mass are all of order $`F_v/M_{Pl}`$. In RG scaling to the weak scale, the gaugino masses give an additional (positive) contribution to the sfermion (mass)<sup>2</sup> $`\alpha /\pi \mathrm{ln}\left(M_{GUT}/m_Z\right)\left[F_v/M_{Pl}\right]^2\left[F_v/M_{Pl}\right]^2`$.
It is clear that if $`R3M^1`$, then the SUGRA contribution to the sfermion (mass)<sup>2</sup> at the high scale from contact Kähler terms (Eq. (13) with a coefficient $`e^{MR}`$) can be neglected in comparison to the GM contribution at $`M_{GUT}`$ and the contribution generated by gaugino masses in RG scaling. This means that sfermion masses conserve flavor, i.e., sfermion with the same gauge quantum numbers are degenerate at the $`O(e^3)`$ percent level. <sup>16</sup><sup>16</sup>16This degeneracy is sufficient to evade limits from $`CP`$-conserving flavor-violating processes, for example, $`\mu e\gamma `$, $`\mathrm{\Delta }m_K`$ etc. But, if the SUGRA mediated sfermion (mass)<sup>2</sup> in Eq. (13) have $`O(1)`$ phases, then we need $`R6M^1`$ to obtain degeneracy at the $`0.1\%`$ level so that SUSY contributions to $`CP`$ and flavor-violating processes are suppressed.
As mentioned earlier, if the SUGRA contribution to sfermion and gaugino masses is neglected, then the (gauge mediated contribution to) RH slepton (mass)<sup>2</sup> at the weak scale is negative. However, in this ($`5D`$) framework, while the SUGRA contribution to RH slepton mass at the high scale (Eq. (13)) is suppressed, we have to include the SUGRA contribution ($`F_v/M_{Pl}`$) to the bino mass which, in turn, generates in RG scaling to the weak scale an additional (compared to the pure GM case) positive contribution $`\left[F_v/M_{Pl}\right]^2`$ to the RH slepton (mass)<sup>2</sup>. This can result in a positive RH slepton (mass)<sup>2</sup> at the weak scale, i.e., unlike the pure GM case, due to the SUGRA contribution to the bino mass, the bino mass contribution is (effectively) independent of (and of the same order as) the GM contribution to the RH slepton mass. <sup>17</sup><sup>17</sup>17Thus, the role of the extra dimension here (as in ) is to suppress the (flavor-violating) SUGRA contribution to sfermion masses while allowing SUGRA contribution to gaugino masses (in particular, in this case, bino mass). It is clear that any other framework which provides these boundary conditions also suffices.
It is obvious that the same framework (extra dimension of size $`R3M^1`$) can be used to suppress (flavor-violating) SUGRA contribution to sfermion masses (Eq. (13)) in any model where the GM and SUGRA contributions are comparable, for example, a model of GM with messenger scale close to the GUT scale (where the GM contribution to sfermion (mass)<sup>2</sup> is, say, positive) . As mentioned earlier, any model with GUT and SUSY breaking by the same field is likely to have negative GM contribution to sfermion (say, RH slepton) (mass)<sup>2</sup>; it is clear that (in addition to suppressing the SUGRA contribution to sfermion masses) the above idea (the SUGRA contribution to bino mass) can be used to obtain positive RH slepton (mass)<sup>2</sup>.
A slightly modified version of the above model is obtained by gauging only the $`SU(5)`$ subgroup of the $`SU(6)_{GUT}`$ symmetry . The superpotential is given by $`W_1+W^{}`$ with
$$W^{}=h\left(\lambda _1\mathrm{\Sigma }_1+\lambda _{24}\mathrm{\Sigma }_{24}\right)\mathrm{\Sigma }_{\overline{5}}+\overline{h}\left(\overline{\lambda }_1\mathrm{\Sigma }_1+\overline{\lambda }_{24}\mathrm{\Sigma }_{24}\right)\mathrm{\Sigma }_5,$$
(17)
where $`\mathrm{\Sigma }_r`$’s refer to components of $`\mathrm{\Sigma }`$ transforming as $`𝐫`$ under $`SU(5)_{local}`$ and $`h`$, $`\overline{h}`$ form a $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$ of $`SU(5)_{local}`$. $`\mathrm{\Sigma }`$ and $`Q`$, $`\overline{Q}`$ are as usual localized on the SUSY breaking brane with $`h`$, $`\overline{h}`$ fields in the bulk. In the absence of $`W^{}`$, $`\mathrm{\Sigma }`$ has a pair of massless color triplets which are Nambu-Goldstone fields since the full $`SU(6)_{GUT}`$ is not gauged. $`W^{}`$ gives masses to these triplets with those in $`h`$, $`\overline{h}`$. Since $`\mathrm{\Sigma }_1\text{diag}[1,1,1,1,1]`$ whereas $`\mathrm{\Sigma }_{24}\text{diag}[3,3,2,2,2]`$ in $`SU(5)`$ space, we can (fine-)tune the $`\lambda _{1,24}`$ couplings so that the weak doublets in $`h`$, $`\overline{h}`$ are massless; these will be the usual Higgs doublets.
In this version of the model, doublet-triplet splitting (although “technically natural”) is fine-tuned. Also, there is a global $`SU(6)`$ symmetry on the SUSY breaking brane, i.e., $`W_1`$ is $`SU(6)`$ symmetric, but the couplings of the bulk fields, $`h`$, $`\overline{h}`$ to $`\mathrm{\Sigma }`$ (Eq. (17)) break this to $`SU(5)_{local}`$. The nice feature compared to the earlier model is that quarks and leptons are contained in the usual $`(\overline{\mathrm{𝟓}}+\mathrm{𝟏𝟎})`$ of $`SU(5)`$ (localized on the matter brane) and so, unlike the gauged $`SU(6)`$ model, “splitting” of matter superfields is not required. As mentioned earlier, in the gauged $`SU(6)`$ model, quarks and leptons are contained in $`(\mathrm{𝟏𝟓}+\overline{\mathrm{𝟔}}+\overline{\mathrm{𝟔}})`$ of $`SU(6)`$ and a $`(\overline{\mathrm{𝟓}}+\mathrm{𝟓})`$ (under $`SU(5)`$) are made heavy by coupling to $`\overline{H}`$ in the $`4D`$ model – this coupling is not allowed (at the renormalizable level) in the $`5D`$ model since $`\overline{H}`$ and matter fields are localized on different branes. In the $`SU(5)`$ model, there are SUGRA contributions to MSSM gaugino masses from both singlet $`(\mathrm{\Sigma }_1)`$ and adjoint ($`\mathrm{\Sigma }_{24}`$) SUSY breaking fields (in addition to the GM contribution). Thus, as in the gauged $`SU(6)`$ model, the MSSM gaugino masses ($`M_1`$, $`M_2`$ and $`M_3`$) are free parameters (see section 4).
#### Comments on other models:
Some of the models of gauge mediation in the literature () also have gauge messengers and hence negative MSSM scalar (mass)<sup>2</sup> at the messenger scale, $`M_{mess}`$ (the SUGRA contribution is much smaller than the GM contribution if $`M_{mess}\stackrel{<}{}10^{15}`$ GeV). As usual, the gaugino masses give a positive contribution in RG scaling to the weak scale; however, at least RH sleptons still have negative (mass)<sup>2</sup> . The framework of $`\stackrel{~}{g}`$MSB can also be used to “resurrect” these models as follows. <sup>18</sup><sup>18</sup>18 This was partly hinted in . Suppose the SUSY breaking fields are localized on a brane different than the MSSM matter brane in an extra dimension with a compactification scale ($`R^1`$) slightly smaller than $`M_{mess}`$ (with gauge fields in the bulk). <sup>19</sup><sup>19</sup>19As before, we assume that the two branes are maximally separated in the extra dimension. Then the two-loop (negative) gauge mediated contribution to sfermion (mass)<sup>2</sup> at $`M_{mess}`$ is suppressed (in the $`5D`$ theory) by a factor $`1/(RM_{mess})^21/25`$ (if $`R5M_{mess}^1`$) relative to the $`4D`$ result, i.e., it is $`1/25\left[\alpha /\pi \right]^2\left[F/M_{mess}\right]^2`$ . The gaugino masses are the same as in $`4D`$ ($`\alpha /\pi F/M_{mess}`$) and, in turn, give a positive $`O(\alpha /\pi F/M_{mess})^2`$ contribution to sfermion (mass)<sup>2</sup> in RG scaling to low scales (provided $`M_{mess}m_Z`$, i.e, RG logarithm is large enough). Since the (negative) sfermion (mass)<sup>2</sup> at $`M_{mess}`$ is small compared to this RG contribution, the sfermion (including RH slepton) (mass)<sup>2</sup> at the weak scale can be positive.
Above the scale $`M_{mess}`$, the theory is $`5D`$ with MSSM matter fields on a brane and gauge fields in bulk. If Higgs fields are also in the bulk, then in this framework gauge coupling unification is (approximately) preserved (with lower unification scale) .
We can use a similar idea to suppress (negative) gauge mediated sfermion (mass)<sup>2</sup> (at $`M_{GUT}`$) in the GUT model, i.e, we can invoke an extra dimension slightly larger than (inverse of) the GUT scale (which is $`M_{mess}`$ in this case). This will suppress both the SUGRA and (negative) GM contributions to sfermion (mass)<sup>2</sup> at the high scale. But, with (say) $`R5M_{GUT}^1(4\times 10^{15}\text{GeV})^1`$, we get the $`5D`$ Planck scale $`M10M_{GUT}`$ (using $`M_{Pl}^2M^3R`$) so that the motivation for the inverted hierarchy mechanism (to generate $`M_{GUT}`$ smaller than the fundamental Planck scale by a factor of $`100`$ as before) is a bit weaker. <sup>20</sup><sup>20</sup>20An even larger extra dimension (in which only gravity propagates) can be used to lower $`M`$ all the way to $`M_{GUT}`$ so that there is no hierarchy between the fundamental Planck scale and GUT scale – of course, in this case (as in the case with $`R5M_{GUT}^1`$) one has to explain the “hiererchy” between $`R^1`$ and $`M`$. Here (as in ), instead we would like to “explain” $`M_{GUT}10^2M`$ using the inverted hierarchy mechanism. In addition a new “hierarchy”, $`R50M^1`$, would have to be explained. <sup>21</sup><sup>21</sup>21In this model, the RG scaling of SM gauge couplings above the energy scale $`R^14\times 10^{15}`$ GeV is $`5D`$ – unification of SM gauge couplings still occurs (Higgs fields also propagate in $`5D`$) but at a scale lower than the usual $`M_{GUT}`$ by a factor of $`2`$ . For illustrative purposes, we assumed above that the unification scale is still $`M_{GUT}2\times 10^{16}`$ GeV. In general, we can choose the compactification scale, $`R^1`$, to be much smaller than $`10^{15}`$ GeV so that unification occurs at a scale, $`M_{GUT}^{}`$, which depends on $`R^1`$ and which is much smaller than the usual $`M_{GUT}10^{16}`$ GeV . However, in this case, one has to “explain” why the compactification scale, $`R^1`$, is correlated with the GUT scale ($`M_{GUT}^{}`$) which (in the context of the model in this paper), in turn, is determined by a modulus field (i.e, which is not a fundamental scale). Also, if $`RM^1`$, then the $`5D`$ gauge couplings (assuming that $`4D`$ gauge couplings are $`O(1)`$) might become non-perturbative (larger than their strong coupling value) (see Eq. (16)). Of course, one faces a similar issue(s) in trying to “save” the GM models above except that in that case the correlation between $`M_{mess}`$ (which is presumably fixed by a modulus field also) and $`R^1`$ is weaker since we only require $`R\stackrel{>}{}5M_{mess}^1`$. In contrast, in the GUT model with $`R3M^1`$ there is only a “modest” hierarchy between $`R^1`$ and the $`5D`$ Planck scale $`M`$ (which, as mentioned earlier, is assumed to be fundamental). In any case, in the model with $`R5M_{GUT}^1`$, the SUGRA and GM contributions to gaugino masses will be (roughly) same as in the model with $`R3M^1`$, i.e., $`M_{1,2,3}`$ are free parameters which will generate positive sfermion (mass)<sup>2</sup> in RG scaling to the weak scale – the only difference is that in the model with $`R3M^1`$, there is a (negative) GM contribution to scalar (mass)<sup>2</sup> at $`M_{GUT}`$. The phenomeonlogy of these two models should be similar (see section 4).
It is also clear from the discussion in section 2.2 that, in general, in models with (dominant) SUSY breaking in a GUT non-singlet (denoted by $`\mathrm{\Sigma }`$), there is a contribution to MSSM gaugino masses at one-loop (and to scalar (mass)<sup>2</sup> at two-loop) from the coupling to gauge messengers – to repeat, these are the heavy gauge multiplets (with mass $`M_{GUT}`$) which have a non-supersymmetric spectrum since the SUSY breaking field ($`\mathrm{\Sigma }`$) transforms under the GUT gauge group. The (contributions to) MSSM gaugino masses generated at one-loop by integrating out gauge messengers (at the scale $`M_{GUT}`$) are generically (i.e., barring accidental cancellations) given by $`\alpha /\pi F_\mathrm{\Sigma }M_\overline{)}R/M_{GUT}^2`$, where $`M_\overline{)}R`$ is the $`R`$-symmetry breaking scale – thus the size of this contribution is model-dependent due to $`M_\overline{)}R`$. The point is that if the field $`\mathrm{\Sigma }`$ also breaks the GUT symmetry (down to the SM gauge group), i.e., if the vev of the scalar component of $`\mathrm{\Sigma }`$ ($`v_\mathrm{\Sigma }`$) is $`O(M_{GUT})`$ (as in our model), then $`M_\overline{)}RM_{GUT}`$. Therefore, in a model with $`F_\mathrm{\Sigma }0`$, if $`v_\mathrm{\Sigma }`$ (or, in general, $`M_\overline{)}R`$) is $`O(M_{GUT})`$, then this gauge messenger contribution to MSSM gaugino masses is comparable to (and independent of) the SUGRA contribution (from the operator $`d^2\theta \mathrm{\Sigma }/M_{Pl}W_\alpha W^\alpha +h.c.`$) <sup>22</sup><sup>22</sup>22We assume that $`\mathrm{\Sigma }`$ is in a representation which appears in the symmetric product of two adjoints. $`F_\mathrm{\Sigma }/M_{Pl}`$ (and is also non-universal, in general). <sup>23</sup><sup>23</sup>23There might also be other GM contribution to gaugino (and scalar) masses from, say, “matter” messengers (as in our model). Thus, as in our $`SU(6)`$ GUT model, the MSSM gaugino mass relations are modified from that expected with only SUGRA contribution – in fact the model becomes less predictive (as far as gaugino masses are concerned) since there is an extra parameter corresponding to the gauge messenger contribution.
There have been some recent studies of MSSM gaugino masses in a scenario with SUSY breaking by a GUT non-singlet so that SUGRA contribution to gaugino masses is non-universal – in these studies, the above gauge messenger contribution has not been mentioned. In the first reference in , it is assumed that $`v_\mathrm{\Sigma }0`$, i.e., the SUSY breaking field has a small vev in its scalar component. In this case, it is possible that $`M_\overline{)}RM_{GUT}`$ so that the gauge messenger contribution to MSSM gaugino masses is small compared to the SUGRA contribution. Nonetheless, (in general) in order to be sure that the gauge messenger contribution to MSSM gaugino masses is smaller than the SUGRA contribution, the complete model has to be analysed to check that $`v_\mathrm{\Sigma }(\text{or}M_\overline{)}R)M_{GUT}`$. Also, a contribution to MSSM scalar (mass)$`{}_{}{}^{2}[\alpha /\pi ]^2[F_\mathrm{\Sigma }/M_{GUT}]^2`$ is generated at two-loops by integrating out gauge messengers – there is no suppression due to $`M_\overline{)}R`$ unlike in the case of MSSM gaugino masses (since scalar (mass)<sup>2</sup> do not break $`R`$-symmetry). Thus, the gauge messenger contribution to scalar (mass)<sup>2</sup> is comparable to (and again, independent of) the SUGRA contribution to scalar (mass)$`{}_{}{}^{2}[F_\mathrm{\Sigma }/M_{Pl}]^2`$ (due to contact Kähler terms) even if $`v_\mathrm{\Sigma }0`$ (i.e., the size of this contribution is model-independent). Furthermore, the gauge messenger contribution depends on gauge quantum numbers and thus it is different for squarks and sleptons (as in our $`SU(6)`$ GUT model), although it is flavor-conserving.
## 4 Sparticle spectrum
We now present briefly a sample sparticle spectrum in the $`5D`$ model.
### 4.1 Parameters of the model
In this model, the MSSM sfermion masses (at $`M_{GUT}`$) and gaugino masses are determined by three parameters – $`F_v/v(vM_{GUT})`$, $`c_2`$ and $`c_3`$. The GM contribution (Eqs. (11) and (12)) can be written in terms of $`F_v/v`$ (assuming that the beta-functions <sup>24</sup><sup>24</sup>24As mentioned earlier, even though the $`N\overline{P}\overline{H}`$ coupling is not allowed (at the renormalizable level) in the $`5D`$ model, we assume that the additional $`(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$ in $`N`$, $`\overline{P}`$’s have a mass of $`M_{GUT}`$ and so the values of the beta-functions are the same as in section 2.1 (the $`4D`$ model), i.e., $`b_6=11`$ and $`b_A`$’s are the usual MSSM beta-functions. and gauge couplings are fixed) and the SUGRA contrbution to MSSM gaugino masses, Eq. (14), depends on $`F_v/v`$ and $`c_{2,3}`$ (any uncertainty in the ratio of $`vM_{GUT}`$ and $`M`$ (or $`M_{Pl}`$) can be absorbed into $`c_{2,3}`$). For convenience, we choose the gaugino masses at the GUT scale, $`M_1`$, $`M_2`$ and $`M_3`$ (which are combinations of these parameters) to be the free parameters. <sup>25</sup><sup>25</sup>25Note that if only the adjoint field (and not singlets) breaks SUSY, then $`c_3=0`$ in Eq. (14) and one has only two free parameters. Using Eqs. (11) and (14), we get:
$$\frac{1}{36}(5M_1+3M_2+2M_3)=\frac{\alpha _6}{4\pi }\frac{F_v}{v}$$
(18)
which parametrizes the GM contribution. This relation is used to determine the GM contribution to scalar (mass)<sup>2</sup> at the GUT scale in terms of $`M_A`$’s using Eq. (12) with $`\mu M_{GUT}`$. The scalar masses are then evolved to the weak scale using the one-loop RG equations.
We note that the SUSY-GUT prediction for $`\mathrm{sin}^2\theta _w`$ (in terms of $`\alpha _s`$ and $`\alpha _{em}`$) is affected by the operator $`W_\alpha W^\alpha \mathrm{\Sigma }/M^2`$ in Eq. (14) – in this model, this effect (which is at the $`M_{GUT}/M_{Pl}`$, i.e., percent level) is related to SUGRA contribution to gaugino mass. The method used to give mass to extra doublets in $`H`$, $`\overline{H}`$ (see section 2.1) also affects the prediction for $`\mathrm{sin}^2\theta _w`$ – in this case this pair of doublets gets a mass of $`O(M_{GUT}^2/M)<M_{GUT}`$ (from the superpotential $`W_4`$) which shifts the $`\mathrm{sin}^2\theta _w`$ prediction by about a percent. Since we wish to illustrate the main features of the spectrum in this paper, we will neglect these effects (which might cancel each other).
We also neglect RG scaling of sfermion masses (at one-loop due to SUGRA contribution to gaugino masses and at two-loops due to SUGRA contribution ($`F_v/M_{Pl}`$) to $`Q`$, $`\overline{Q}`$ etc. masses) between the ($`5D`$) Planck scale $`M`$ and the GUT scale (sfermion masses are negligible at the scale $`M`$) <sup>26</sup><sup>26</sup>26This contribution is positive and flavor-conserving – thus it makes sfermions with same gauge quantum numbers more degenerate and in particular the RH slepton heavier (see later). – the RG logarithm $`\mathrm{ln}\left(M/M_{GUT}\right)`$ is much smaller than that for RG scaling between GUT and weak scale (of course, the larger group theory factors above $`M_{GUT}`$ might compensate for the smaller RG logarithm as emphasized in the context of extra dimensional models in ). <sup>27</sup><sup>27</sup>27A detailed study of the phenomenology in this GUT model (including the effects of RG scaling between $`M`$ and $`M_{GUT}`$) is in progress. The (SUGRA mediated) gaugino masses also run between $`M`$ and $`M_{GUT}`$; however since $`M_A/\alpha _A`$ is RG-invariant (at one-loop), the ratio of the MSSM gaugino masses remains the same in this RG scaling since the gauge couplings are unified (of course, at $`M_{GUT}`$, the MSSM gauginos get additional contribution to their masses (Eq. (11)).
Since the usual Higgs doublets propagate in the extra dimension, soft SUSY breaking Higgs (mass)<sup>2</sup> and also $`B\mu `$ and $`\mu `$ (Giudice-Masiero mechanism ) are generated by the SUGRA interactions:
$``$ $``$ $`{\displaystyle }d^4\theta ({\displaystyle \frac{1}{M^3}}h^{}h[\mathrm{\Sigma }^{}\mathrm{\Sigma }+..]+{\displaystyle \frac{1}{M^3}}\overline{h}\overline{h}^{}[X^{}X+..])+`$ (19)
$`{\displaystyle }d^4\theta (h\overline{h}{\displaystyle \frac{1}{M^2}}[\mathrm{\Sigma }^{}+X^{}+\overline{X}^{}]+h\overline{h}{\displaystyle \frac{1}{M^3}}[\mathrm{\Sigma }^{}\mathrm{\Sigma }+..]+h.c.),`$
where $`h`$, $`\overline{h}`$ are $`5D`$ fields. The couplings of zero-modes of $`h`$, $`\overline{h}`$ (which correspond to the light $`4D`$ Higgs fields) are suppressed by a factor of $`\sqrt{R}`$ (to account for the normalization of the zero mode) so that $`m_{H_{u,d},SUGRA}^2F_v^2/M^3\mathrm{\hspace{0.33em}1}/R`$, $`\mu F_v/M^2\mathrm{\hspace{0.33em}1}/R`$ and $`B\mu F_v^2/M^3\mathrm{\hspace{0.33em}1}/R`$. Since $`MR3`$ and $`M_{Pl}^2M^3R`$, we see that $`m_{H_{u,d},SUGRA}^2`$, $`\mu ^2`$, $`B\mu `$ are all of order $`\left(F_v/M_{Pl}\right)^2`$, i.e, of the same order as, but independent of, gaugino and sfermion masses (they are also independent of each other). Of course, the Higgs doublets also get soft (mass)<sup>2</sup> from GM (Eq. (12), which are related to the other sfermion and gaugino masses. We choose the Higgs soft (mass)<sup>2</sup> at the GUT scale (which are the sum of the SUGRA and GM contributions) to be free parameters, denoted by $`m_{H_u}^2`$ and $`m_{H_d}^2`$.
In this model, the gauginos of the heavy gauge multiplet have a SUSY breaking mass $`F_v/M_{GUT}`$ since the SUSY breaking field is an adjoint under the GUT gauge group. This generates trilinear (MSSM) scalar terms of $`O(\alpha /\pi F_v/M_{GUT})F_v/M_{Pl}`$ at the GUT scale (when the heavy gauginos are integrated out). The exact expression is :
$`V`$ $``$ $`{\displaystyle \underset{i}{}}A_iQ_i_{Q_i}W(Q),`$ (20)
$`A_i(\mu _{RG})`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}Z_{Q_i}(v^{},v,\mu _{RG})}{\mathrm{ln}v}}{\displaystyle \frac{F_v}{v}}.`$
In this case, we have
$`A_i(M_{GUT})`$ $`=`$ $`{\displaystyle \frac{F_v}{v}}{\displaystyle \frac{\alpha _6}{4\pi }}\left(2C_6^i2{\displaystyle \underset{A}{}}C_A^i\right).`$ (21)
We neglect all Yukawa couplings other than the top quark coupling and so only the coupling $`\lambda _tA_tH_u\stackrel{~}{Q}_3\stackrel{~}{t}^c`$ is non-zero and is given by
$`A_t(M_{GUT})`$ $`=`$ $`15.3{\displaystyle \frac{F_v}{v}}{\displaystyle \frac{\alpha _6}{4\pi }}.`$ (22)
Thus, the $`A`$-term at the GUT scale depends only on $`F_v/v`$ (and gauge couplings), i.e., it is not an independent parameter – in particular, there is no SUGRA contribution to $`A_t`$ since the top squark and the SUSY breaking fields are localized on separate branes. <sup>28</sup><sup>28</sup>28As mentioned earlier, we neglect RG scaling between GUT and Planck scales; this effect does generate (due to non-vanishing gaugino masses) a small $`A_t`$ term at the GUT scale which depends on the SUGRA contribution to gaugino masses, i.e., $`F_v/v`$, $`c_2`$ and $`c_3`$.
Thus, the fundamental parameters in this model are: $`F_v/v`$, $`c_{2,3}`$, $`m_{H_{u,d}}^2`$, $`B\mu `$, $`\mu `$ and $`\lambda _t`$ (top quark Yukawa coupling). Two of these parameters are fixed by the observed values of $`m_Z`$ and $`m_t`$ so that the free (i.e., input) parameters can be chosen to be $`M_{1,2,3}`$ (which, as explained earlier are combinations of $`F_v/v`$ and $`c_{2,3}`$), $`m_{H_{u,d}}^2`$ and $`\mathrm{tan}\beta `$; $`\mu `$ and $`B\mu `$ can then be determined in terms of these parameters as usual using the minimization conditions for the Higgs potential.
### 4.2 Sample sparticle spectrum
In Table 2, a sample spectrum is presented for the input parameters $`M_1=M_2=300`$ GeV and $`M_3=150`$ GeV, $`m_{H_u}^2=(150\text{GeV})^2`$, $`m_{H_d}^2=(300\text{GeV})^2`$ and $`\mathrm{tan}\beta =5`$. We have included the electroweak $`D`$-term and Fayet-Illiopoulos (FI) $`D`$-term contributions to scalar (mass)<sup>2</sup> which are given by $`m_Z^2\mathrm{cos}2\beta \left(T_3Q\mathrm{sin}^2\theta _w\right)`$ and $`0.053Y\left(m_{H_u}^2m_{H_d}^2\right)`$ <sup>29</sup><sup>29</sup>29The FI $`D`$-term contribution in $`\stackrel{~}{g}`$MSB was emphasized in ., respectively. <sup>30</sup><sup>30</sup>30We assume for simplicity that there is no $`D`$-term contribution from the breaking of the extra $`U(1)`$ (of $`SU(6)`$) at $`M_{GUT}`$. The mixing between the top squarks and the one-loop corrections to the effective Higgs potential (due to top quark and top squark masses only) have been included.
Some of the characteristic features of the spectrum are as follows.
Gaugino masses are non-universal in general since both the GM contribution and a part of the SUGRA contribution (due to the operator $`d^2\theta c_2\mathrm{\Sigma }/M^2W_\alpha W^\alpha `$) are non-universal (the relative GM contributions to gaugino masses (Eq. (11)) are model-dependent due to dependence on $`b_6`$). Models with non-universal gaugino masses (at the GUT/Planck scale) have been studied earlier . In most of these models, sfermion masses are independent parameters whereas in our GUT model the sfermion masses are determined in terms of the gaugino masses ($`M_{1,2,3}`$). <sup>31</sup><sup>31</sup>31In “$`D`$-brane” models , it is possible that $`M_2M_3`$ if $`SU(3)_c`$ and $`SU(2)_w`$ originate from different $`D`$-brane sectors. Since quark doublets transform under all three SM gauge groups, $`U(1)_Y`$ has to orginate in either of these two sectors – this implies that $`M_1=M_2`$ or $`M_1=M_3`$, unlike the GUT model where it is possible that all three gaugino masses are different.
No-scale SUGRA models have vanishing scalar masses at $`M_{Pl}`$ and gaugino masses (which can be non-universal) as usual drive sfermion (mass)<sup>2</sup> positive in RG scaling to the weak scale ($`\stackrel{~}{g}`$MSB with non-universal gaugino masses has similar boundary conditions). However, in the GUT model studied here, there is a GM contribution to scalar masses at the GUT scale and also there are SUGRA contributions to Higgs soft masses (which, in turn, result in a FI $`D`$-term contribution to sfermion masses at the weak scale). Thus, the GUT model can (in principle) be distinguished from no-scale SUGRA models with non-universal gaugino masses by precision sparticle spectroscopy.
The (GM contribution to) RH slepton (mass)<sup>2</sup> at the GUT scale is negative (Eq. (12)) and therefore RH slepton (its (mass)<sup>2</sup> is driven positive by bino mass) and $`\chi _1^0`$ (which is roughly the bino) are close in mass. <sup>32</sup><sup>32</sup>32In this sample spectrum, the FI $`D`$-term (positive for RH slepton) makes the RH slepton (slightly) heavier than $`\chi _1^0`$ (for one sign of $`\mu `$). For the same reason, the mass splitiing between the left-handed slepton and the RH slepton is large <sup>33</sup><sup>33</sup>33If the GM contribution at $`M_{GUT}`$ is neglected, then we get $`m_{\stackrel{~}{e}_R}140`$ GeV while $`m_{\stackrel{~}{e}_L}`$ remains about the same. even though we have chosen $`M_1=M_2`$ for the above sample spectrum, i.e., usually we expect $`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}`$ to be large only if $`M_2>M_1`$ (due to gaugino mass contributions in RG scaling). As mentioned earlier, if we include RG scaling between $`M`$ and $`M_{GUT}`$, then all sfermions will be heavier leading to a larger mass splitting between RH slepton and $`\chi _1^0`$ (which will be the lightest supersymmetric particle (LSP)) whereas the large mass splitting between $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$ (which is due to the negative GM contribution to $`m_{\stackrel{~}{e}_R}^2`$ at $`M_{GUT}`$) will remain about the same.
The lower limit on the RH slepton mass $`90`$ GeV fixes (roughly) a minimum value for $`M_1`$. But since the three gaugino masses are independent parameters, $`M_3`$ can be smaller than $`M_{1,2}`$ so that there is not much of a hierarchy (in masses) between squarks/gluino and sleptons as seen in Table 2. Also, since $`M_3`$ can be smaller than $`M_{1,2}`$ and also (GM mediated) stop (and other squark) (mass)<sup>2</sup> are small and negative at the GUT scale, $`|m_{H_u}^2|`$ at the weak scale (as usual the up-type Higgs (mass)<sup>2</sup> is driven negative by stop (mass)<sup>2</sup> and gluino mass) and hence $`\mu `$ can be small (in this case $`180`$ GeV), thus reducing the fine-tuning in electroweak symmetry breaking – small $`\mu `$ also results in “light” chargino/neutralino. <sup>34</sup><sup>34</sup>34As mentioned earlier, RG scaling between $`M_{GUT}`$ and $`M`$ will give a positive contribution to scalar (mass)<sup>2</sup> due to (SUGRA mediated) gaugino masses – this contribution is about the same for all scalars (squarks, sleptons and Higgs) because of the unified gauge symmetry, unlike the case of RG scaling below the GUT scale where, since $`\alpha _3>\alpha _{1,2}`$ the gluino mass contribution (to squark masses) is larger (assuming universal gaugino masses). Thus, the above features, i.e., the “small” hierarchy between squarks and sleptons and small $`\mu `$ will not be affected by the inclusion of this effect.
This should be compared to $`\stackrel{~}{g}`$MSB with universal gaugino mass where a minimum value for $`M_1`$ fixed by RH slepton mass implies a minimum value for $`M_3`$ (typically $`\stackrel{>}{}200`$ GeV) which results in a larger hierarchy betwen slepton and squark/gluino masses and also larger $`|m_{H_u}^2|`$ and hence fine tuning (due to larger $`\mu `$). In a minimal gauge mediation model also, there is a large hierarchy between squark/gluino masses and slepton masses (since masses are proportional to gauge couplings). In minimal SUGRA mediated SUSY breaking (with universal scalar mass, $`m_0`$, and universal gaugino mass, $`M_{1/2}`$) it is possible to have small hierarchy between sleptons and squarks/gluino. In any case, non-universal gaugino masses distinguishes these models from our model.
Of course, we expect that with extra parameters as compared to a minimal model (especially non-universal gaugino masses) such a spectrum can be attained. However in this model these extra parameters are not ad hoc, but are well-motivated – they are justified by the way SUSY is broken in the model.
## 5 Summary
A model in which the same scalar potential breaks SUSY and a GUT symmetry was presented in – this model has dynamical origins for both SUSY breaking and GUT scales. In this model, the SUGRA and gauge mediated contributions to scalar and gaugino masses are comparable – this enables a viable spectrum to be attained since the gauge mediated contribution to RH slepton (mass)<sup>2</sup> by itself is negative. But, the flip side is that the SUGRA contribution to sfermion masses (from non-renormalizable contact Kähler terms) results in flavor violation.
In this paper, we suggested that this “problem” will be present in any model in which the same field breaks SUSY and a GUT symmetry and demonstrated that, using an extra spatial dimension, positive and (at the same time) flavor-conserving sfermion (mass)<sup>2</sup> can be obtained in this model. The model has non-universal gaugino masses and sfermion masses are predicted in terms of gaugino masses. The hierarchy between squark/gluino masses and slepton masses can be small and (typically) a large mass splitting between RH and LH slepton is expected.
Acknowledgments
The author thanks Markus Luty and Nir Polonsky for suggestions and Neal Weiner for comments. This work is supported by DOE Grant DE-FG03-96ER40969. |
warning/0003/math0003211.html | ar5iv | text | # Cauchy-Riemann Geometry and Contact Topology in Three Dimensions
## I Introduction
We study low-dimensional problems in topology and geometry via a study of contact and Cauchy-Riemann ($`CR`$) structures. Let us start with a closed (compact without boundary) oriented three-manifold $`M`$. A contact structure (or bundle) $`\xi `$ on $`M`$ is a completely non-integrable rank 2 subbundle of $`TM`$. It is well known that there are no local invariants for contact structures according to a classical theorem of Darboux. Also, two nearby contact structures on a closed manifold are isotopy-equivalent by Gray’s theorem (Gray, 1959; Hamilton, 1982). Therefore, a contact structure has no continuous moduli. In this sense, it is a kind of geometric structure even softer than a complex structure. The isotopy classes are distinguished by so-called tight or overtwisted contact structures (Eliashberg, 1992). The existence of contact structures on a closed oriented three-manifold is known from the work of Martinet (1971) and Lutz (1977). (See also Altschuler (1995) for an analytic proof using the so-called linear contact flow.)
Given a contact structure, we can consider a $`CR`$-structure, i.e., a “complex structure” defined on a contact bundle. Different from the usual complex structure, a $`CR`$-structure does have local invariants. Thus, analysis is needed. In Section II, we give a brief introduction to $`CR`$-geometry and an application in K$`\ddot{a}`$hler geometry. In Section III, we introduce a global $`CR`$-invariant $`\mu _\xi `$ and discuss its behavior on the moduli space of $`CR`$-structures. Also, we argue that the contractibility of our $`CR`$ moduli space for $`S^3`$ confirms the so-called Smale conjecture.
In Section IV, we discuss spherical $`CR`$-structures: the critical points of $`\mu _\xi `$. To distinguish among “$`CR`$ lens” spaces, we propose a possible $`CR`$-invariant defined for spherical $`CR`$-structures, which is a contact-analogue of Ray-Singer’s analytic torsion. In Section V, we give a heuristic argument for how our understanding of $`\mu _\xi `$ can be applied to the problem of counting the number of complex structures on a closed four-manifold. In Section VI, we propose the study of a certain kind of monopole equation for contact three-manifolds.
## II Basics in $`CR`$-Geometry
A $`CR`$-structure $`J`$ compatible with the contact structure $`\xi `$ is a complex structure on $`\xi `$, i.e., a bundle endomorphism $`J:\xi \xi `$, such that $`J^2=Identity`$. Natural examples come from boundaries of strictly pseudoconvex domains $`D`$ in $`C^2`$. Let $`J_{C^2}`$ denote the multiplication by $`i`$ in $`C^2`$. Let our three-manifold $`M=D`$, the boundary of $`D`$. The contact structure $`\xi `$ is considered to be the intersection of $`TM`$ and $`J_{C^2}TM`$, the tangent subspaces invariant under $`J_{C^2}`$. In addition, our $`CR`$-structure is taken to be a restriction of $`J_{C^2}`$ on $`\xi `$. This $`CR`$-structure is usually called the $`CR`$-structure induced from $`C^2`$.
In his famous theorem, Fefferman (1974) asserts that two strictly pseudoconvex domains with smooth boundaries in $`C^{n+1}`$ are biholomorphic to each other if and only if their boundaries are $`CR`$-equivalent. Therefore the $`CR`$-structure on the boundary reflects the complex structure of the inside domain. It is well known that we have the Riemann mapping theorem in $`C^1`$. However, this theorem is no longer true for higher dimensions. Indeed, we do have local invariants for our $`CR`$ manifold $`(M,\xi ,J)`$ (e.g., Cartan (1932) and Chern and Moser (1974)).
First, choose eigenvectors $`Z_1`$, $`Z_{\overline{1}}`$ of $`J`$ with eigenvalues $`i`$, $`i`$, respectively. Let $`\{\theta ^1,\theta ^{\overline{1}}\}`$ be a set of complex one-forms dual to $`\{Z_1,Z_{\overline{1}}\}`$. Then, choose a local one-form $`\theta `$ annihilating $`\xi `$ (called contact form) so that
$$d\theta =ih_{1\overline{1}}\theta ^1\theta ^{\overline{1}}+\theta \varphi $$
for some real one-form $`\varphi `$ and positive $`h_{1\overline{1}}`$. (We will use $`h_{1\overline{1}}`$ and $`h^{1\overline{1}}=\left(h_{1\overline{1}}\right)^1`$ to raise or lower indices.) Now, for a different choice of coframe $`(\stackrel{~}{\theta },\stackrel{~}{\theta }^1,\stackrel{~}{\theta }^{\overline{1}};\stackrel{~}{\varphi })`$ satisfying the above equation, we have the following transformation relation:
$$\{\begin{array}{ccc}\stackrel{~}{\theta }=u\theta \hfill & & \\ \stackrel{~}{\theta }^1=u_{1}^{}{}_{}{}^{1}\theta ^1+v^1\theta \hfill & & \\ \stackrel{~}{\varphi }=\frac{du}{u}+\varphi +2Re\left(iu^1v^1u_{\overline{1}1}\theta ^{\overline{1}}\right)+s\theta \hfill & & \end{array}$$
(1)
for positive $`u`$ and some real function $`s`$. Differentiating $`\theta ^1`$, $`\varphi `$ gives the first structural equations:
$$\{\begin{array}{cc}d\theta ^1=\theta ^1\varphi _{1}^{}{}_{}{}^{1}+\theta \varphi ^1\hfill & \\ d\varphi =2Re\left(i\theta _{\overline{1}}\varphi ^{\overline{1}}\right)+\theta \psi \hfill & \end{array}$$
(2)
for the connection forms $`\varphi _{1}^{}{}_{}{}^{1},\varphi ^1,\psi `$. Differentiating the connection forms again and requiring certain trace conditions (e.g., Chern and Moser (1974)), we obtain the second set of structural equations:
$$\{\begin{array}{ccc}d\varphi _{1}^{}{}_{}{}^{1}i\theta _1\varphi ^1+2i\varphi _1\theta ^1+\frac{1}{2}\psi \theta =0\hfill & & \\ d\varphi ^1\varphi \varphi ^1\varphi ^1\varphi _{1}^{}{}_{}{}^{1}+\frac{1}{2}\psi \theta ^1=Q_{}^{1}{}_{\overline{1}}{}^{}\theta ^{\overline{1}}\theta \hfill & & \\ d\psi \varphi \psi 2i\varphi ^1\varphi _1=\left(R_1\theta ^1+R_{\overline{1}}\theta ^{\overline{1}}\right)\theta ,\hfill & & \end{array}$$
(3)
in which $`Q_{}^{1}{}_{\overline{1}}{}^{}`$ or $`Q_{11}`$ is called the Cartan (curvature) tensor, and $`R_1,R_{\overline{1}}`$ are determined by means of suitable covariant derivatives of $`Q_{}^{1}{}_{\overline{1}}{}^{}`$ (Cheng, 1987). The normalization condition: $`\varphi \varphi _{1}^{}{}_{}{}^{1}\varphi _{\overline{1}}^{}{}_{}{}^{\overline{1}}=0`$ and the above structural equations Eqs.(2) and (3) uniquely determine the connection forms $`\varphi _{1}^{}{}_{}{}^{1},\varphi ^1,\psi `$. Under the change of coframe Eq.(1), the Cartan tensor is transformed as follows:
$$Q_{11}=\stackrel{~}{Q}_{11}u\left(u_{1}^{}{}_{}{}^{1}\right)^2.$$
The fundamental theorem of 3-dimensional $`CR`$-geometry due to Cartan (1932a,1932b) asserts that $`Q_{11}=0`$ if and only if $`(M,\xi ,J)`$ is locally $`CR`$-equivalent to $`(S^3,\widehat{\xi },\widehat{J})`$, where $`(\widehat{\xi },\widehat{J})`$ denotes the standard $`CR`$-structure on the unit 3-sphere $`S^3`$, induced from $`C^2`$.
$`\mathrm{𝐃𝐞𝐟𝐢𝐧𝐢𝐭𝐢𝐨𝐧}`$. We call a $`CR`$ manifold $`(M,\xi ,J)`$ or just $`J`$ spherical if it is locally $`CR`$-equivalent to $`(S^3,\widehat{\xi },\widehat{J})`$. Quantitatively, a $`CR`$-structure is spherical if $`Q_{11}=0`$ according to Cartan’s theorem.
$`\mathrm{𝐀𝐧}\mathrm{𝐀𝐩𝐩𝐥𝐢𝐜𝐚𝐭𝐢𝐨𝐧}\mathrm{𝐢𝐧}𝐊\ddot{𝐚}\mathrm{𝐡𝐥𝐞𝐫}\mathrm{𝐆𝐞𝐨𝐦𝐞𝐭𝐫𝐲}`$
Let $`N`$ be an $`n`$-dimensional K$`\ddot{a}`$hler manifold. Suppose we have a holomorphic line bundle $`L`$ with the first Chern class being the K$`\ddot{a}`$hler class so that a suitable circle bundle $`ML`$ with the induced $`CR`$-structure is closely related to the K$`\ddot{a}`$hler geometry of $`N`$ (Webster, 1977). It turns out that we can identify (up to a constant) the Cartan tensor $`Q_{11}`$ of $`M`$ with $`R_{,11}`$, the covariant derivative of the scalar curvature $`R`$ of $`N`$ in the (1,0)-direction twice. When $`n2`$, we can identify the Chern tensor (Chern and Moser, 1974, 1983) in higher dimensional $`CR`$-geometry with the Bochner tensor of $`N`$. In 1977, Sid Webster applied $`CR`$-geometry to obtain the following result:
Let $`N`$ be a simply-connected closed K$`\ddot{a}`$hler manifold of dimension $`n`$. Suppose that $`N`$ admits a Hodge metric for which the Bochner tensor vanishes if $`n2`$ or for which $`R_{,11}`$ vanishes if $`n=1`$. Then, $`N`$ is holomorphically isometric to complex projective space $`CP^n`$ with a standard Fubini-Study metric. (Webster, 1977)
Next, relative to a special coframe $`\left(\theta ,\theta ^1,\theta ^{\overline{1}};\varphi =0\right)`$ satisfying $`d\theta =ih_{1\overline{1}}\theta ^1\theta ^{\overline{1}}`$, we can define the so-called pseudohermitian connection $`\omega _{1}^{}{}_{}{}^{1}`$, torsion $`A_{11}`$, and curvature $`𝒲`$, called the Tanaka-Webster curvature (Tanaka, 1975; Webster, 1977). These data are uniquely determined by the following equations:
$$\{\begin{array}{ccc}d\theta ^1=\theta ^1\omega _{1}^{}{}_{}{}^{1}+A_{}^{1}{}_{\overline{1}}{}^{}\theta \theta ^{\overline{1}}\hfill & & \\ d\omega _{1}^{}{}_{}{}^{1}=𝒲\theta ^1\theta ^{\overline{1}}\left(mod\theta \right)\hfill & & \\ \omega _{1}^{}{}_{}{}^{1}+\omega _{\overline{1}}^{}{}_{}{}^{\overline{1}}=h^{1\overline{1}}dh_{1\overline{1}}.\hfill & & \end{array}$$
The torsion $`A_{11}`$ and the Tanaka-Webster curvature $`𝒲`$ are not “tensorial” under the change of contact form $`\stackrel{~}{\theta }=u\theta `$ (Lee, 1986), but are “tensorial” under the change $`\stackrel{~}{\theta }^1=u_{1}^{}{}_{}{}^{1}\theta ^1`$. The Cartan tensor can be expressed in terms of these data (Cheng and Lee, 1990):
$$Q_{11}=\frac{1}{6}𝒲_{,11}+\frac{i}{2}𝒲A_{11}A_{11,0}\frac{2i}{3}A_{11,\overline{1}1}.$$
Here, covariant derivatives are taken with respect to the pseudohermitian connection $`\omega _{1}^{}{}_{}{}^{1}`$, and “0” means the $`T`$\- direction. (The tangent vector field $`T`$ is uniquely determined by $`\theta \left(T\right)=1`$ and $`L_T\theta =0`$.) Before going on, another result should be noted:
The boundary of a circular domain in $`C^{n+1}`$ is $`CR`$-equivalent to the unit sphere $`S^{2n+1}C^{n+1}`$ with the standard induced $`CR`$-structure if and only if the Tanaka-Webster curvature $`𝒲constant`$ (with respect to a suitable choice of contact form)(Unpublished paper by J. Bland and P. M. Wang).
The proof of the above result in the original draft contains a gap which can be remedied by the following result:
Let $`N`$ be a closed complex manifold with two K$`\ddot{a}`$hler metrics $`g`$, $`\stackrel{~}{g}`$. Suppose the Bochner tensor of $`g`$ vanishes and the scalar curvature of $`\stackrel{~}{g}`$ is a constant. Then, the fact that the K$`\ddot{a}`$hler class of $`g`$ is cohomologous to the K$`\ddot{a}`$hler class of $`\stackrel{~}{g}`$ implies that $`(N,g)`$ and $`(N,\stackrel{~}{g})`$ are isometric to each other (Chen and Lue, 1981).
## III The $`\mu _\xi `$-Invariant and the Moduli Space
First, we will construct an energy functional on the space of $`CR`$-structures so that the critical points consist of spherical $`CR`$-structures. Let $`\mathrm{\Pi }`$ denote the $`su(2,1)`$-valued Cartan connection form defined by
$$\mathrm{\Pi }=\left(\begin{array}{ccc}\frac{1}{3}\left(\varphi _{1}^{}{}_{}{}^{1}+\varphi \right)& \theta ^1& 2\theta \\ i\varphi _1& \frac{1}{3}\left(2\varphi _{1}^{}{}_{}{}^{1}\varphi \right)& 2i\theta _1\\ \frac{1}{4}\psi & \frac{1}{2}\varphi ^1& \frac{1}{3}\left(\varphi +\varphi _{\overline{1}}^{}{}_{}{}^{\overline{1}}\right)\end{array}\right).$$
The curvature form $`\mathrm{\Omega }`$ is defined as usual by $`\mathrm{\Omega }=d\mathrm{\Pi }\mathrm{\Pi }\mathrm{\Pi }`$. The transgression $`TC_2\left(\mathrm{\Pi }\right)`$ of the second Chern form is given by
$`TC_2\left(\mathrm{\Pi }\right)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}[tr(\mathrm{\Pi }\mathrm{\Omega }+{\displaystyle \frac{1}{3}}tr(\mathrm{\Pi }\mathrm{\Pi }\mathrm{\Pi })]`$
$`=`$ $`{\displaystyle \frac{1}{24\pi ^2}}tr\left(\mathrm{\Pi }\mathrm{\Pi }\mathrm{\Pi }\right)`$
$`\left(sincetr\left(\mathrm{\Pi }\mathrm{\Omega }\right)=0\right).`$
We can verify that the 3-form $`TC_2\left(\mathrm{\Pi }\right)`$ is invariant under the change of contact form and invariant up to an exact form under the coframe change Eq.(1). In the late 1980’s, Burns and Epstein (1988)(also Cheng and Lee (1990)) discovered that the integral of $`TC_2\left(\mathrm{\Pi }\right)`$, denoted as $`\mu _\xi `$, is a global $`CR`$-invariant (assuming trivial holomorphic tangent bundle as in Burns and Epstein (1988); extended to arbitrary $`M`$ by a relative version of the invariant in Cheng and Lee (1990)):
$`\mu _\xi \left(J\right)`$ $`=`$ $`{\displaystyle \frac{1}{24\pi ^2}}{\displaystyle _M}tr\left(\mathrm{\Pi }\mathrm{\Pi }\mathrm{\Pi }\right)`$
$`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _M}\left[2Re\left(i\theta ^1\varphi ^{\overline{1}}\varphi _{1}^{}{}_{}{}^{1}\right)+{\displaystyle \frac{1}{2}}\theta \psi \varphi 2i\theta \varphi ^1\varphi ^{\overline{1}}{\displaystyle \frac{1}{2}}d\left(\theta \psi \right)\right]`$
$`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _M}\left[\left({\displaystyle \frac{1}{6}}𝒲^2+2\left|A_{11}\right|^2\right)\theta d\theta +{\displaystyle \frac{2}{3}}\omega _{1}^{}{}_{}{}^{1}d\omega _{1}^{}{}_{}{}^{1}\right]`$
$`\left(intermsofpseudohermitiangeometry\right).`$
It is remarkable that the above integral is independent of the choice of contact form, and that the integrand involves only the second and lower-order derivatives (relative to a coframe field) while the lowest order of local invariants is of order 4 as indicated by the Cartan tensor $`Q_{11}`$.
Next, we will discuss the moduli space of $`CR`$-structures. Let $`𝒥_\xi `$ denote the space of all $`CR`$-structures compatible with $`\xi `$. Let $`𝒞_\xi `$ denote the group of contact diffeomorphisms with respect to $`\xi `$. Clearly, $`𝒞_\xi `$ acts on $`𝒥_\xi `$ by pulling back. The invariant $`\mu _\xi `$ is actually defined on the moduli space $`𝒥_\xi /𝒞_\xi `$.
Given a $`CR`$-structure $`J`$ in $`𝒥_\xi `$, we call a “submanifold” $`S`$ passing through $`J`$ a local slice if it is transverse to the orbit of $`𝒞_\xi `$-action, so that any element in $`𝒥_\xi `$ near $`J`$ can be pulled back to an element of $`S`$ by means of a certain contact diffeomorphism. In the early 1990’s, Jack Lee and the author proved the following:
Local slices always exist for all cases (Cheng and Lee, 1995).
As a corollary, the standard spherical $`CR`$-structure $`\left[\widehat{J}\right]`$ in $`𝒥_{\widehat{\xi }}/𝒞_{\widehat{\xi }}`$ for $`S^3`$ is a strict local minimum for $`\mu _{\widehat{\xi }}`$ (Cheng and Lee, 1995).
Let $`Q_J=2Re\left[iQ_{1}^{}{}_{}{}^{\overline{1}}\theta ^1Z_{\overline{1}}\right]`$. It is a straightforward computation to obtain the first variation formula: $`\delta \mu _\xi \left(J\right)=\frac{1}{8\pi ^2}Q_J`$. Consider the downward gradient flow for $`\mu _\xi `$:
$$_tJ_{\left(t\right)}=Q_{J_{\left(t\right)}}.$$
(4)
Since $`\delta Q_J`$ is subelliptic modulo the action of our symmetry group $`𝒞_\xi `$, we can play a suitable “De-Turck trick” to break the symmetry and imitate the usual $`L^2`$-theory for elliptic operators to obtain the short time solution of Eq.(4)(Cheng and Lee, 1990). However, we can not prove the long term solution and convergence even for $`M=S^3`$. This is related to the so-called Smale conjecture as first pointed out by Eliashberg.
The Smale conjecture asserts that the diffeomorphism group of $`S^3`$ is homotopy-equivalent to the orthogonal group $`O\left(4\right)`$. Suppose we have the long term solution and convergence of Eq.(4) for $`M=S^3`$. Then, any starting $`J`$ must converge to $`\widehat{J}`$, the unique spherical $`CR`$-structure on $`S^3`$ (up to symmetry). Therefore, the (certain marked) $`CR`$ moduli space $`𝒥_{\widehat{\xi }}^{^{}}/𝒞_{\widehat{\xi }}^{^{}}`$ is contractible. But $`𝒥_{\widehat{\xi }}^{^{}}`$ is contractible, too. It follows that $`𝒞_{\widehat{\xi }}^{^{}}`$ is contractible. Then, with the aid of contact geometry, we can confirm the Smale conjecture.
To learn more analytic techniques which can be used to tackle Eq.(4), we have been working on some comparatively easier flows like the $`CR`$ Calabi flow and the $`CR`$ Yamabe flow. For the $`CR`$ Yamabe flow, S.-C. Chang and the author deformed a contact form in the direction of the Tanaka-Webster curvature:
$$_t\theta _{\left(t\right)}=𝒲\theta _{\left(t\right)}.$$
(5)
In their present work, Chang and Cheng obtain a Harnack estimate and (possibly) the long term solution for Eq.(5).
## IV The Moduli Space of Spherical $`CR`$-Structures
Let $`𝒮_\xi `$ denote the space of all spherical $`CR`$-structures compatible with $`\xi `$. Since the linearization of the Cartan tensor is subelliptic modulo the action of $`𝒞_\xi `$, the virtual dimension of $`𝒮_\xi /𝒞_\xi `$: the moduli space of spherical $`CR`$-structures is finite. Let $`M`$ be a circle bundle over a closed surface of genus $`g>1`$ with the Euler class $`e\left(M\right)<0`$. Let $`Pic(g,c_1)`$, the universal Picard variety, denote the space of all pairs $`(L,N)`$ in which $`L`$ is a holomorphic line bundle over a Riemann surface $`N`$ of genus $`g>1`$ with $`c_1\left(L\right)=e\left(M\right)`$ modulo an equivalence relation defined by diffeomorphisms. In 1996 and 1997, I-Hsun Tsai and the author studied the relation between $`𝒮_\xi /𝒞_\xi `$ and $`Pic(g,c_1)`$. We found the following:
For an above-mentioned circle bundle $`M`$, there is a diffeomorphism between $`𝒮_\xi /^{^{}}𝒞_\xi `$ and $`Pic(g,c_1)^{}`$. (The prime means a suitably modified version.) Moreover, $`Pic(g,c_1)^{}`$ is a complex manifold of dimension $`4g3`$ (Cheng and Tsai, 2000).
Our above result is similar to describing a Teichmuller space by means of conformal classes. It is known in Teichmuller theory that we can pick up a unique hyperbolic metric as a representative for each conformal class. A similar situation occurs for our spherical $`CR`$ manifolds. In fact, our theory for the universal Picard variety has counterparts in Teichmuller theory as shown in Table 1.
Table 1. Comparison of two theories
| Teichmuller space | universal Picard variety |
| --- | --- |
| conformal classes | spherical $`CR`$ circle bundles |
| Riemannian hyperbolic metrics | pseudohermitian hyperbolic geometries |
$`\mathrm{𝐋𝐨𝐜𝐚𝐥}\mathrm{𝐑𝐢𝐠𝐢𝐝𝐢𝐭𝐲}\mathrm{𝐨𝐟}\mathrm{𝐒𝐩𝐡𝐞𝐫𝐢𝐜𝐚𝐥}\mathrm{𝐂𝐑}\mathrm{𝐒𝐭𝐫𝐮𝐜𝐭𝐮𝐫𝐞𝐬}`$
$`\left(\mathrm{𝐃𝐢𝐬𝐜𝐫𝐞𝐭𝐞}\mathrm{𝐌𝐨𝐝𝐮𝐥𝐢}:dim𝒮_\xi /𝒞_\xi =0\right)`$
Let $`Aut_{CR}\left(S^3\right)`$ denote the $`CR`$-automorphism group of $`(S^3,\widehat{\xi },\widehat{J})`$, which is known to be isomorphic to $`SU(2,1)/center`$. Let $`\mathrm{\Gamma }`$ denote a fixed point free finite subgroup of $`Aut_{CR}\left(S^3\right)`$. Then, $`\mathrm{\Gamma }\backslash S^3`$ inherits both contact and (spherical) $`CR`$-structures from $`(S^3,\widehat{\xi },\widehat{J})`$. This induced spherical $`CR`$-structure on $`\mathrm{\Gamma }\backslash S^3`$ is locally rigid; i.e. it has no nontrivial deformation. (The algebraic reason is that $`H^1(\mathrm{\Gamma },su(2,1))=0`$, in which the group cohomology has coefficients in the holonomy representation: developing map composed with the adjoint representation) (Burns and Shnider, 1976). On the other hand, note that $`\mathrm{\Gamma }\backslash S^3`$ has positive constant Tanaka-Webster curvature and zero torsion. Now, generalizing using an analytical method, we obtain the following:
Let $`(M,J)`$ be a closed spherical $`CR`$ three-manifold. Suppose there is a contact form such that the torsion $`A_{11}=0`$ and $`𝒲>0`$, $`4𝒲(5𝒲^2+3\mathrm{\Delta }_b𝒲)3|_b𝒲|_\theta ^2>0`$. Then, $`J`$ is locally rigid (Cheng, 1999).
Next, we want to compare two $`\mathrm{\Gamma }\backslash S^3`$. Suppose $`\mathrm{\Gamma }_1\backslash S^3`$ and $`\mathrm{\Gamma }_2\backslash S^3`$ are diffeomorphic. How can we distinguish one spherical $`CR`$-structure from the other one? (They have the same $`\mu _\xi `$-value.) To deal with this problem, we borrow ideas from quantum physics. If we view $`\mu _\xi `$ as a Lagrangian (action, more accurately) in $`2+1`$ dimensions, spherical $`CR`$-structures are just classical fields. Therefore, “quantum fluctuations” should give us refined invariants. In practice, we compute the partition function heuristically:
$`𝒵_k`$ $`=`$ $`{\displaystyle _{𝒥_\xi /𝒞_\xi }}𝒟\left[J\right]e^{ik\mu _\xi \left(\left[J\right]\right)}`$
$`=`$ $`k^{\frac{dim}{2}}\left(𝒵_{sc}+O\left(k^1\right)\right)\left(klarge\right),`$
in which $`𝒵_{sc}`$ is called the semi-classical approximation. Note that only classical fields make contributions to $`𝒵_{sc}`$. By imitating the finite dimensional case, we can compute the modulus of $`𝒵_{sc}`$ (Cheng, 1995):
$`\left|𝒵_{sc}\right|`$ $`=`$ $`lim_k\mathrm{}k^{\frac{dim}{2}}\left|𝒵_k\right|`$
$`=`$ $`\mathrm{\Sigma }_{J:spherical}\left|{\displaystyle \frac{det\mathrm{}_J}{det^{}\delta Q_J}}\right|^{\frac{1}{2}},`$
in which $`\mathrm{}_J`$ is a fourth-order subelliptic self-adjoint operator related to the $`𝒞_\xi `$-action, and $`\delta Q_J`$, the second variation of $`\mu _\xi `$, is also a fourth-order subelliptic self-adjoint operator modulo the $`𝒞_\xi `$-action. We can regularize two determinants via zeta functions. ($`det^{}`$ means taking a regularized determinant under a certain gauge-fixing condition.)
$`\mathrm{𝐂𝐨𝐧𝐣𝐞𝐜𝐭𝐮𝐫𝐞}`$: If $`J`$ is spherical,
$$Tor\left(J\right)\stackrel{def}{=}\left|\frac{det\mathrm{}_J}{det^{}\delta Q_J}\right|^{\frac{1}{2}}$$
is independent of any choice of contact form, i.e., a $`CR`$ invariant.
We expect to use $`Tor\left(J\right)`$ to distinguish among “contact lens” (or “$`CR`$ lens”) spaces $`\left\{\mathrm{\Gamma }\backslash S^3\right\}`$. Also, we note that $`Tor\left(J\right)`$ is a contact-analogue of Ray-Singer’s analytic torsion while no contact-analogue is known for the Reidemeister torsion.
## V Counting the Number of Complex Structures
This is another “quantum level” problem in our ongoing project. We will discuss the problem of counting the number of complex structures on a closed (compact without boundary) four-manifold. We hope to view this number as the partition function of a certain 3+1 quantum field theory (QFT in short).
Let us begin with a 0+1 theory, i.e., a particle moving in a closed manifold $`N`$. The Hamiltonian of such a theory with supersymmetry is the Laplace-Beltrami operator $`\mathrm{\Delta }`$. All quantum ground states or vacua are cohomology classes of $`N`$, represented by harmonic forms (=zero eigenforms of $`\mathrm{\Delta }`$). Now suppose $`f`$ is a Morse function on $`N`$. Consider $`\mathrm{\Delta }_{tf}`$ in which $`d`$ is replaced by $`e^{tf}de^{tf}`$. When $`t\mathrm{}`$, the harmonic forms of $`\mathrm{\Delta }_{tf}`$ are concentrated near the critical points of f. These are the classical ground states (Witten, 1982).
The harmonic form corresponding to a critical point $`P`$ has a small correction due to another critical point $`Q`$ via the trajectories of $`f`$ from $`P`$ to $`Q`$. This is quantum mechanical tunnelling, which describes the probability of the transition $`PQ`$. The boundary operator of Witten’s chain complex (See Witten (1982) or Atiyah (1988) for a clear explanation.) is interpreted in terms of such tunnelling. (The homology of Witten’s chain complex can be shown to identify with the homology of $`N`$.) Witten’s idea was later adopted by Floer (1989) and applied to the infinite-dimensional case of the manifold of connections.
Next, we will give a brief introduction to the Donaldson-Floer theory. It is a 3+1 QFT. A “field” when restricted to the three-space $`M`$ in this theory is a connection (or gauge field) of a certain, say, $`SU\left(2\right)`$ bundle over $`M`$. The Morse function as mentioned above is the Chern-Simons functional defined on the space of connections in this case. The critical points consist of flat connections which are the classical ground states. Through consideration of the associated Witten complex, we obtain the so-called Floer homology or cohomology group $`HF\left(M\right)`$. This is the space of quantum ground states or vacua for this theory. Now, suppose we decompose a closed 4-manifold $`X`$ along $`M`$ (say, a homology 3-sphere) as shown in Fig.1.
where $`X=X^+_MX^{}`$. Let $`\mathrm{\Sigma }^+`$($`\mathrm{\Sigma }^{}`$, respectively) denote the set of restrictions on $`M`$ of all instantons on $`X^+`$($`X^{}`$, respectively). Then, $`\mathrm{\Sigma }^+`$, $`\mathrm{\Sigma }^{}`$ form cycles in $`HF\left(M\right)`$. The intersection number represents the algebraic number of instantons on $`X`$, (assuming it is finite) the Donaldson invariant, denoted as $`Z\left(X\right)`$. We can write
$$Z\left(X\right)=<vac\left(X^+\right)|vac\left(X^{}\right)>,$$
in which the vacuum $`vac\left(X^+\right)=\left[\mathrm{\Sigma }^+\right]`$ and the vacuum $`vac\left(X^{}\right)=\left[\mathrm{\Sigma }^{}\right]`$ are both elements of $`HF\left(M\right)`$. Also $`<|>`$ denotes the middle-dmension intersection number. In Witten (1988), Witten presented a Lagrangian for this theory so that $`Z\left(X\right)`$ identifies with its partition function.
Now, we can describe our $`3+1QFT`$. We put an auxiliary contact structure $`\xi `$ on our closed oriented three-manifold $`M`$. A “field” is a complex structure with the restriction on $`M`$ being a $`CR`$-structure compatible with $`\xi `$. Our Morse function is the $`\mu _\xi `$ which we introduce in $`\mathrm{\S }`$3. Spherical $`CR`$-structures which are critical points of $`\mu _\xi `$ are our classical ground states in this theory.
Let $`\mathrm{\Sigma }^+`$($`\mathrm{\Sigma }^{}`$, respectively) denote the set of all $`CR`$-structures compatible with $`\xi `$ on $`M`$, which can be extended to a complex structure on $`X^+`$$`(X^{}`$, respectively). Now, what is the associated “Floer” homology group $`HF(M,\xi )`$, i.e., the space of quantum vacua, for this theory? Since the Hessian $`\delta ^2\mu _\xi `$ at a spherical $`J`$ is subelliptic modulo $`𝒞_\xi `$, the dimension of its negative eigenspace is finite. Therefore, the Morse index is well defined. (We do not need the relative Morse index as in the case of the Donaldson-Floer theory.) As usual, $`\mathrm{\Sigma }^\pm `$ form cycles $`\left[\mathrm{\Sigma }^\pm \right]`$ in $`HF(M,\xi )`$ by pushing along the gradient flow of $`\mu _\xi `$ and seeing which critical points they “hang” on (Atiyah, 1988). The vacuum $`vac\left(X^+\right)`$($`vac\left(X^{}\right)`$, respectively) is defined as the homology class $`\left[\mathrm{\Sigma }^+\right]`$ ($`\left[\mathrm{\Sigma }^{}\right]`$, respectively) in $`HF(M,\xi )`$. Moreover, we define the quantity $`Z_\xi \left(X\right)`$ as
$`Z_\xi \left(X\right)`$ $`\stackrel{def}{=}`$ $`<vac\left(X^+\right)|vac\left(X^{}\right)>`$
$`\stackrel{def}{=}`$ $`intersectionnumberof\left[\mathrm{\Sigma }^+\right]and\left[\mathrm{\Sigma }^{}\right].`$
The sum of $`Z_\xi \left(X\right)`$ over the isomorphism classes of tight contact structures, denoted as $`Z\left(X\right)`$, can be interpreted as the (algebraic) number of complex structures on $`X`$. We propose the following “physical” problem:
$`\mathrm{𝐏𝐫𝐨𝐛𝐥𝐞𝐦}\mathbf{\hspace{0.25em}1}`$. Find a Lagrangian for the above theory so that its partition function identifies with $`Z\left(X\right)`$.
There are topological obstructions for $`M`$ to admit spherical $`CR`$-structures (Goldman, 1983). For instance, the three-torus $`T^3`$ does not admit any spherical $`CR`$-structure (compatible with any given contact structure $`\xi `$). Therefore, $`HF_{}(T^3,\xi )=0`$ for any $`\xi `$, and we can propose the following problem for “nonexistence”:
$`\mathrm{𝐏𝐫𝐨𝐛𝐥𝐞𝐦}\mathbf{\hspace{0.25em}2}`$. Suppose $`X=X^+_{T^3}X^{}`$. Find conditions on $`X`$ and, perhaps, $`X^\pm `$ such that $`Z\left(X\right)=0`$.
We still need to investigate the relation between $`Z\left(X\right)=0`$ and the nonexistence of complex structures. Another situation occurs when $`M`$ is the standard contact 3-sphere $`(S^3,\widehat{\xi })`$. This admits only one compatible spherical $`CR`$-structure, namely, the standard one $`\widehat{J}`$, which is a strict local minimum for $`\mu _{\widehat{\xi }}`$ modulo symmetry as mentioned in Section III. It follows that $`HF_0(S^3,\widehat{\xi })=Z`$ and $`HF_k(S^3,\widehat{\xi })=0`$ for $`k0`$. Therefore, we can propose the following problem concerning “global rigidity”:
$`\mathrm{𝐏𝐫𝐨𝐛𝐥𝐞𝐦}\mathbf{\hspace{0.25em}3}`$. Suppose $`X=X^+_{S^3}X^{}`$. Find conditions on $`X`$ and, perhaps, $`X^\pm `$ such that $`Z\left(X\right)=1`$.
Note that any tight contact structure on $`S^3`$ is isotopy-equivalent to $`\widehat{\xi }`$ according to Eliashberg (1992). Therefore, $`Z\left(X\right)`$ in Problem 3 is just $`Z_{\widehat{\xi }}\left(X\right)`$.
## VI Monopoles and Contact Structures
Recently, Kronheimer and Mrowka (1997) studied contact structures on 3-manifolds via the 4-dimensional Seiberg-Witten monopole theory. Here, we will outline another approach by Cheng and Chiu (1999).
Given a contact 3-manifold $`(M,\xi )`$ and a background pseudohermitian structure $`(J,\theta )`$, we can discuss a canonical $`spin^c`$-structure $`c_\xi `$ on $`\xi ^{}`$. With respect to $`c_\xi `$, we will consider the equations for our “monopole” $`\mathrm{\Phi }`$ coupled to the “gauge field” $`A`$. Here, $`A`$, the $`spin^c`$-connection, is required to be compatible with the pseudohermitian connection on $`M`$. The Dirac operator $`D_\xi `$ relative to $`A`$ is identified with a certain boundary $`\overline{}`$-operator $`\sqrt{2}\left(\overline{}_b^a+\left(\overline{}_b^a\right)^{}\right)`$. In terms of the components $`(\alpha ,\beta )`$ of $`\mathrm{\Phi }`$, our equations read as
$$\{\begin{array}{c}\left(\overline{}_b^a+\left(\overline{}_b^a\right)^{}\right)\left(\alpha +\beta \right)=0\\ \left(or\alpha _{,\overline{1}}^a=0,\beta _{\overline{1},1}^a=0\right)\\ da(e_1,e_2)𝒲=\left|\alpha \right|^2\left|\beta _{\overline{1}}\right|^2,\end{array}$$
(6)
where $`A=A_{can}+iaI`$ and $`𝒲`$ denotes the Tanaka-Webster curvature. Our first step in understanding Eq.(6) is as follows:
Suppose the torsion $`A_{11}=0`$. Also, suppose $`\xi `$ is symplectically semifillable, and that the Euler class $`e(\xi )`$ is not a torsion class. Then, Eq.(6) has nontrivial solutions (i.e., $`\alpha `$ and $`\beta `$ are not identically zero simultaneously)(Cheng and Chiu, 1999).
On the other hand, the Weitzenbock-type formula gives a nonexistence result for $`𝒲>0`$. Together with the above existence result, we can conclude the following:
Suppose the torsion $`A_{11}=0`$ and the Tanaka-Webster curvature $`𝒲>0`$. Then, either $`\xi `$ is not symplectically semifillable, or $`e(\xi )`$ is a torsion class (Cheng and Chiu, 1999).
We note that Rumin (1994) proved that $`M`$ must be a rational homology sphere under the conditions given above using a different method. Also, we do not know how to deal with the solution space of Eq.(6) in general although we hope that further study of Eq.(6) will produce invariants of contact structures.
$`\mathrm{𝐀𝐜𝐤𝐧𝐨𝐰𝐥𝐞𝐝𝐠𝐦𝐞𝐧𝐭}`$
In preparing this article the author benefited from a number of conversations with I-Hsun Tsai and Chin-Lung Wang. Also, the author would like to thank Professors Chuu-Lian Terng and I-Hsun Tsai for inviting him to lecture at the NCTS conference. The main part of this article is based on the author’s notes for that talk. The research was partially supported by the National Science Council under grant NSC 88-2115-M-001-015 (R.O.C.).
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warning/0003/hep-th0003199.html | ar5iv | text | # Consistent interactions in the Hamiltonian BRST formalism
## 1 Introduction
The analysis of consistent interactions that can be introduced among fields with gauge freedom without changing the number of gauge symmetries has been transposed lately at the level of the deformation of the master equation from the antifield-BRST formalism . This cohomological deformation technique has been applied, among others, to Chern-Simons models , Yang-Mills theories and two-form gauge fields . In this light, the antifield-BRST method was proved to be an elegant tool for investigating the problem of consistent interactions.
Recently, a Hamiltonian analysis of anomalies has been given . Moreover, in a very interesting paper , there has been established the precise relation of the local BRST cohomologies in both Lagrangian and Hamiltonian formalisms (see Theorem 6 from this reference). The procedures developed within these papers strongly stimulate a Hamiltonian BRST approach to other interesting problems.
In this letter we analyze the problem of constructing consistent interactions among fields with gauge freedom in the framework of the Hamiltonian BRST formalism , . Our strategy includes two main steps: (i) initially, we show that the problem of introducing consistent interactions among fields with gauge freedom can be reformulated as a problem of deforming the BRST charge and the BRST-invariant Hamiltonian with respect to a given “free” theory, and consequently we deduce the general equations that govern these two types of deformations; (ii) next, on behalf of the relationship between the Hamiltonian and antifield BRST formalisms for constrained systems we prove that the general equations possess solution. In the sequel, we reformulate the general equations in a manner that accounts for locality, and subsequently illustrate our general procedure in the case of three-dimensional Chern-Simons models. Finally, we remark that our method combined with the results in may simplify the computation of local Lagrangian BRST cohomologies in some cases of interest.
## 2 General equations of the Hamiltonian deformation approach
We begin with a system described by the canonical variables $`z^A`$, subject to the first-class constraints
$$G_{a_0}\left(z^A\right)0,a_0=1,\mathrm{},M_0,$$
(1)
which are assumed to be $`L`$-stage reducible
$$G_{a_0}Z_{a_1}^{a_0}=0,a_1=1,\mathrm{},M_1,$$
(2)
$$Z_{a_{k1}}^{a_{k2}}Z_{a_k}^{a_{k1}}0,a_k=1,\mathrm{},M_{k,}k=2,\mathrm{},L,$$
(3)
and suppose that there are no second-class constraints in the theory. The Grassmann parities of the canonical variables and first-class constraints are respectively denoted by $`\epsilon \left(z^A\right)=\epsilon _A`$ and $`\epsilon \left(G_{a_0}\right)=\epsilon _{a_0}`$. We denote the first-class Hamiltonian by $`H_0`$, such that the gauge algebra is expressed by
$$[G_{a_0},G_{b_0}]=G_{c_0}C_{a_0b_0}^{c_0},[H_0,G_{a_0}]=G_{b_0}V_{a_0}^{b_0}.$$
(4)
It is known that a constrained Hamiltonian system can be described by the action
$$S_0[z^A,u^{a_0}]=\underset{t_1}{\overset{t_2}{}}𝑑t\left(a_A\left(z\right)\dot{z}^AH_0G_{a_0}u^{a_0}\right),$$
(5)
where the Grassmann parities of the Lagrange multipliers are given by $`\epsilon \left(u^{a_0}\right)=\epsilon _{a_0}`$. In (5), $`a_A\left(z\right)`$ is the one-form potential that gives the symplectic two-form $`\omega _{AB}=()^{\epsilon _A+1}\frac{^La_A}{z^B}()^{\epsilon _B\left(\epsilon _A+1\right)}\frac{^La_B}{z^A}`$, whose inverse, $`\omega ^{AB}`$, corresponds to the fundamental Dirac brackets $`[z^A,z^B]=\omega ^{AB}`$. Action (5) is invariant under the gauge transformations
$$\delta _ϵz^A=[z^A,G_{a_0}]ϵ^{a_0},\delta _ϵu^{a_0}=\dot{ϵ}^{a_0}V_{b_0}^{a_0}ϵ^{b_0}C_{b_0c_0}^{a_0}ϵ^{c_0}u^{b_0}Z_{a_1}^{a_0}ϵ^{a_1}.$$
(6)
In order to generate consistent interactions at the Hamiltonian level, we deform the action (5) by adding to it some interaction terms
$$S_0\stackrel{~}{S}_0=S_0+g\underset{0}{\overset{(1)}{S}}+g^2\underset{0}{\overset{(2)}{S}}+\mathrm{},$$
(7)
and modify the gauge transformations (6) (to be denoted by $`\stackrel{~}{\delta }_ϵz^A`$, $`\stackrel{~}{\delta }_ϵu^{a_0}`$) in such a way that the deformed gauge transformations leave invariant the new action
$$\frac{\delta \stackrel{~}{S}_0}{\delta z^A}\stackrel{~}{\delta }_ϵz^A+\frac{\delta \stackrel{~}{S}_0}{\delta u^{a_0}}\stackrel{~}{\delta }_ϵu^{a_0}=0.$$
(8)
Consequently, the deformation of the action (5) and of the gauge transformations (6) produces a deformation of the first-class constraints, first-class Hamiltonian and structure functions like
$$G_{a_0}\gamma _{a_0}=G_{a_0}+g\underset{a_0}{\overset{(1)}{\gamma }}+g^2\underset{a_0}{\overset{(2)}{\gamma }}+\mathrm{},$$
(9)
$$H_0H=H_0+g\stackrel{(1)}{H}+g^2\stackrel{(2)}{H}+\mathrm{},$$
(10)
$$V_{b_0}^{a_0}\stackrel{~}{V}_{b_0}^{a_0}=V_{b_0}^{a_0}+g\underset{b_0}{\overset{a_0}{\stackrel{(1)}{V}}}+g^2\underset{b_0}{\overset{a_0}{\stackrel{(2)}{V}}}+\mathrm{},$$
(11)
$$C_{b_0c_0}^{a_0}\stackrel{~}{C}_{b_0c_0}^{a_0}=C_{b_0c_0}^{a_0}+g\underset{b_0c_0}{\overset{a_0}{\stackrel{(1)}{C}}}+g^2\underset{b_0c_0}{\overset{a_0}{\stackrel{(2)}{C}}}+\mathrm{},$$
(12)
such that the deformed gauge algebra becomes
$$[\gamma _{a_0},\gamma _{b_0}]=\gamma _{c_0}\stackrel{~}{C}_{a_0b_0}^{c_0},[H,\gamma _{a_0}]=\gamma _{b_0}\stackrel{~}{V}_{a_0}^{b_0}.$$
(13)
In the meantime, we deform the reducibility relations, but we do not explicitly write down these relations.
As the BRST charge and BRST-invariant Hamiltonian contain all the information on the gauge structure of a given gauge theory, we can reformulate the problem of introducing consistent interactions within the Hamiltonian BRST context in terms of these two essential compounds. Indeed, if the interaction can be consistently constructed, then the BRST charge of the undeformed theory, $`\stackrel{(0)}{\mathrm{\Omega }}`$, can be deformed such as to be the BRST charge of the deformed theory, i.e.,
$$\stackrel{(0)}{\mathrm{\Omega }}\mathrm{\Omega }=\stackrel{(0)}{\mathrm{\Omega }}+g\stackrel{(1)}{\mathrm{\Omega }}+g^2\stackrel{(2)}{\mathrm{\Omega }}+\mathrm{},$$
(14)
$$[\mathrm{\Omega },\mathrm{\Omega }]=0.$$
(15)
At the same time, the deformation of the BRST charge induces the deformation of the BRST-invariant Hamiltonian of the undeformed theory, $`\underset{B}{\overset{(0)}{H}}`$,
$$\underset{B}{\overset{(0)}{H}}H_B=\underset{B}{\overset{(0)}{H}}+g\underset{B}{\overset{(1)}{H}}+g^2\underset{B}{\overset{(2)}{H}}+\mathrm{},$$
(16)
in such a way that $`H_B`$ is the BRST-invariant Hamiltonian of the interacting theory, i.e.
$$[H_B,\mathrm{\Omega }]=0.$$
(17)
The equations (15) and (17) split accordingly the deformation parameter as
$$[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(0)}{\mathrm{\Omega }}]=0,[\underset{B}{\overset{(0)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]=0,$$
(18)
$$2[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]=0,[\underset{B}{\overset{(0)}{H}},\stackrel{(1)}{\mathrm{\Omega }}]+[\underset{B}{\overset{(1)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]=0,$$
(19)
$$2[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(2)}{\mathrm{\Omega }}]+[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]=0,[\underset{B}{\overset{(0)}{H}},\stackrel{(2)}{\mathrm{\Omega }}]+[\underset{B}{\overset{(1)}{H}},\stackrel{(1)}{\mathrm{\Omega }}]+[\underset{B}{\overset{(2)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]=0,$$
(20)
$$\mathrm{}$$
Equations (1820) stand for the general equations of our deformation procedure. The equations (18) are checked by hypothesis. Then, it appears the natural question whether the next equations possess or not solution. This will be investigated in the next section.
## 3 Solution to the general equations
In order to prove that the equations (1920), etc. possess solution, we use the link between the antifield and Hamiltonian BRST formalisms for constrained Hamiltonian systems . First-class constrained Hamiltonian systems can be approached from the point of view of the BRST formalism in two different manners. One is based on the antibracket-antifield formulation , while the other relies on the standard Hamiltonian BRST treatment , . The starting point of the antibracket-antifield formalism is represented by the invariance of the action (5) under the gauge transformations (6). In agreement with the general prescriptions of the antibracket-antifield procedure, we introduce the ghosts $`(\eta ^{a_{k1}},u^{a_k})`$, with $`k=1,\mathrm{},L`$ and
$$\epsilon \left(\eta ^{a_k}\right)=\epsilon _{a_k}+k+1mod\mathrm{\hspace{0.33em}2},gh\left(\eta ^{a_k}\right)=k+1,k=0,\mathrm{}L,$$
(21)
$$\epsilon \left(u^{a_k}\right)=\epsilon _{a_k}+kmod\mathrm{\hspace{0.33em}2},gh\left(u^{a_k}\right)=k,k=1,\mathrm{}L,$$
(22)
where $`gh`$ denotes the ghost number. The antifields associated with the fields $`(z^A,u^{a_0},\eta ^{a_{k1}},u^{a_k})`$ are denoted by $`(z_A^{},u_{a_0}^{},\eta _{a_{k1}}^{},u_{a_k}^{})`$ and display the properties $`\epsilon \left(antifield\right)=\epsilon \left(field\right)+1`$, $`gh\left(antifield\right)=gh\left(field\right)1`$. Up to terms that are quadratic in the antifields, the solution to the master equation reads as
$`\stackrel{(0)}{S}={\displaystyle \underset{t_1}{\overset{t_2}{}}}dt(a_A\left(z\right)\dot{z}^A+{\displaystyle \underset{k=0}{\overset{L}{}}}u_{a_k}^{}\dot{\eta }^{a_k}H_0G_{a_0}u^{a_0}+z_A^{}[z^A,G_{a_0}]\eta ^{a_0}`$
$`u_{a_0}^{}V_{b_0}^{a_0}\eta ^{b_0}+()^{\epsilon _{b_0}+1}u_{a_0}^{}C_{b_0c_0}^{a_0}\eta ^{c_0}u^{b_0}+{\displaystyle \frac{1}{2}}()^{\epsilon _{b_0}}\eta _{a_0}^{}C_{b_0c_0}^{a_0}\eta ^{c_0}\eta ^{b_0}+`$
$`{\displaystyle \underset{k=0}{\overset{L1}{}}}\eta _{a_k}^{}Z_{a_{k+1}}^{a_k}\eta ^{a_{k+1}}{\displaystyle \underset{k=1}{\overset{L1}{}}}u_{a_{k1}}^{}Z_{a_k}^{a_{k1}}u^{a_k}+\mathrm{}).`$ (23)
The Hamiltonian point of view is based on extending the phase-space through introducing the canonical pairs ghost-antighost $`(\eta ^{a_k},𝒫_{a_k})`$, with $`[\eta ^{a_k},𝒫_{a_k}]=\delta _{b_k}^{a_k}`$ and $`\epsilon \left(𝒫_{a_k}\right)=\epsilon _{a_k}+k+1`$, $`gh\left(𝒫_{a_k}\right)=k+1`$. The BRST charge starts like
$$\stackrel{(0)}{\mathrm{\Omega }}=G_{a_0}\eta ^{a_0}+\frac{1}{2}()^{\epsilon _{b_0}}𝒫_{a_0}C_{b_0c_0}^{a_0}\eta ^{c_0}\eta ^{b_0}+\underset{k=0}{\overset{L1}{}}𝒫_{a_k}Z_{a_{k+1}}^{a_k}\eta ^{a_{k+1}}+\mathrm{},$$
(24)
such that $`[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(0)}{\mathrm{\Omega }}]=0`$. The BRST-invariant extension of $`H_0`$
$$\underset{B}{\overset{(0)}{H}}=H_0+𝒫_{a_0}V_{b_0}^{a_0}\eta ^{b_0}+\mathrm{},$$
(25)
satisfies the equation $`[\underset{B}{\overset{(0)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]=0`$. By employing the identifications
$$u_{a_k}^{}=𝒫_{a_k},k=0,\mathrm{},L,$$
(26)
and extending the Dirac bracket such that $`[\eta ^{a_k},u_{a_k}^{}]=\delta _{b_k}^{a_k}`$, we get that
$`{\displaystyle \frac{1}{2}}(\stackrel{(0)}{S},\stackrel{(0)}{S})={\displaystyle \underset{t_1}{\overset{t_2}{}}}dt({\displaystyle \frac{d}{dt}}\stackrel{(0)}{\mathrm{\Omega }}[\underset{B}{\overset{(0)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]+{\displaystyle \frac{1}{2}}z_A^{}[z^A,[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(0)}{\mathrm{\Omega }}]]+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=0}{\overset{L}{}}}\eta _{a_k}^{}[\eta ^{a_k},[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(0)}{\mathrm{\Omega }}]]+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=0}{\overset{L}{}}}[[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(0)}{\mathrm{\Omega }}],u_{a_k}^{}]u^{a_k}).`$ (27)
The deformations (14) and (16) induce a deformation of the solution to the master equation
$$\stackrel{(0)}{S}S=\stackrel{(0)}{S}+g\stackrel{(1)}{S}+g^2\stackrel{(2)}{S}+\mathrm{},$$
(28)
such that the equation (3) for the deformed theory becomes
$`{\displaystyle \frac{1}{2}}(S,S)={\displaystyle \underset{t_1}{\overset{t_2}{}}}dt({\displaystyle \frac{d}{dt}}\mathrm{\Omega }[H_B,\mathrm{\Omega }]+{\displaystyle \frac{1}{2}}z_A^{}[z^A,[\mathrm{\Omega },\mathrm{\Omega }]]+`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=0}{\overset{L}{}}}\eta _{a_k}^{}[\eta ^{a_k},[\mathrm{\Omega },\mathrm{\Omega }]]+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=0}{\overset{L}{}}}[[\mathrm{\Omega },\mathrm{\Omega }],u_{a_k}^{}]u^{a_k}).`$ (29)
The equation (3) splits accordingly the deformation parameter as (3) and
$`(\stackrel{(0)}{S},\stackrel{(1)}{S})={\displaystyle \underset{t_1}{\overset{t_2}{}}}dt({\displaystyle \frac{d}{dt}}\stackrel{(1)}{\mathrm{\Omega }}[\underset{B}{\overset{(0)}{H}},\stackrel{(1)}{\mathrm{\Omega }}][\underset{B}{\overset{(1)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]+z_A^{}[z^A,[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]]+`$
$`{\displaystyle \underset{k=0}{\overset{L}{}}}\eta _{a_k}^{}[\eta ^{a_k},[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]]+{\displaystyle \underset{k=0}{\overset{L}{}}}[[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}],u_{a_k}^{}]u^{a_k}),`$ (30)
$`(\stackrel{(0)}{S},\stackrel{(2)}{S})+{\displaystyle \frac{1}{2}}(\stackrel{(1)}{S},\stackrel{(1)}{S})={\displaystyle \underset{t_1}{\overset{t_2}{}}}dt({\displaystyle \frac{d}{dt}}\stackrel{(2)}{\mathrm{\Omega }}[\underset{B}{\overset{(0)}{H}},\stackrel{(2)}{\mathrm{\Omega }}][\underset{B}{\overset{(1)}{H}},\stackrel{(1)}{\mathrm{\Omega }}]`$
$`[\underset{B}{\overset{(2)}{H}},\stackrel{(0)}{\mathrm{\Omega }}]+z_A^{}[z^A,[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(2)}{\mathrm{\Omega }}]+{\displaystyle \frac{1}{2}}[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]]+`$
$`{\displaystyle \underset{k=0}{\overset{L}{}}}\eta _{a_k}^{}[\eta ^{a_k},[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(2)}{\mathrm{\Omega }}]+{\displaystyle \frac{1}{2}}[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]]+`$
$`{\displaystyle \underset{k=0}{\overset{L}{}}}[[\stackrel{(0)}{\mathrm{\Omega }},\stackrel{(2)}{\mathrm{\Omega }}]+{\displaystyle \frac{1}{2}}[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}],u_{a_k}^{}]u^{a_k}),`$ (31)
$$\mathrm{}$$
The last equations emphasize that the existence of $`\stackrel{(1)}{S}`$ guarantees the existence of $`\stackrel{(1)}{\mathrm{\Omega }}`$ and $`\underset{B}{\overset{(1)}{H}}`$, the existence of $`\stackrel{(2)}{S}`$ guarantees the existence of $`\stackrel{(2)}{\mathrm{\Omega }}`$ and $`\underset{B}{\overset{(2)}{H}}`$, and so on. Moreover, the equations (1920), etc. are equivalent to the equations $`(\stackrel{(0)}{S},\stackrel{(1)}{S})=0`$, $`(\stackrel{(0)}{S},\stackrel{(2)}{S})+\frac{1}{2}(\stackrel{(1)}{S},\stackrel{(1)}{S})=0`$, etc. modulo imposing some appropriate boundary conditions for $`\mathrm{\Omega }`$ . On the other hand, the last equations possess solution. The existence of such solutions was proved in on behalf of the triviality of the antibracket in the cohomology. Thus, the existence of the solutions in the antibracket proves the existence of the solutions to (1920), etc. In conclusion, we can construct consistent interactions by means of the equations (1920), etc.
In practical applications, as commonly required, the deformation should be local, i.e., $`\stackrel{(1)}{\mathrm{\Omega }}`$, $`\stackrel{(2)}{\mathrm{\Omega }}`$, $`\underset{B}{\overset{(1)}{H}}`$, $`\underset{B}{\overset{(2)}{H}}`$, etc. should be local functionals. Let $`F_1=d^{D1}xf_1`$ and $`F_2=d^{D1}xf_2`$ be two local functionals. Then, $`[F_1,F_2]`$ is local, namely, there exists a local $`[f_1,f_2]`$ (but defined up to a $`\left(D1\right)`$-dimensional divergence), such that $`[F_1,F_2]=d^{D1}x[f_1,f_2]`$. Thus, the equations (1920), etc. can be written as
$$2\stackrel{(0)}{s}\stackrel{(1)}{\omega }=^k\underset{k}{\overset{(1)}{j}},\stackrel{(0)}{s}\underset{B}{\overset{(1)}{h}}+[\underset{B}{\overset{(0)}{h}},\stackrel{(1)}{\omega }]=^k\underset{k}{\overset{(1)}{m}},$$
(32)
$$2\stackrel{(0)}{s}\stackrel{(2)}{\omega }+[\stackrel{(1)}{\omega },\stackrel{(1)}{\omega }]=^k\underset{k}{\overset{(2)}{j}},\stackrel{(0)}{s}\underset{B}{\overset{(2)}{h}}+[\underset{B}{\overset{(1)}{h}},\stackrel{(1)}{\omega }]+[\underset{B}{\overset{(0)}{h}},\stackrel{(2)}{\omega }]=^k\underset{k}{\overset{(2)}{m}},$$
(33)
$$\mathrm{}$$
in terms of the integrands $`\underset{B}{\overset{(k)}{h}}`$ and $`\stackrel{(k)}{\omega }`$. In the above, $`\stackrel{(0)}{s}`$ stands for the undeformed BRST symmetry. The formalism developed so far does not guarantee locality. For instance, even if $`[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]`$ is $`\stackrel{(0)}{s}`$-exact, it is not granted that it is the BRST variation of a local functional. Such locality problems appear also in the Lagrangian deformation procedure . However, in the case of most important applications , , the Lagrangian BRST deformation procedure leads to local interactions. Thus, we expect that the Hamiltonian BRST deformation treatment also outputs local vertices in practical applications.
## 4 Example
Let us exemplify the prior procedure in the case of abelian Chern-Simons model in three dimensions. We start with the Lagrangian action
$$S_0\left[A_\mu ^a\right]=\frac{1}{2}d^3x\epsilon ^{\mu \nu \rho }k_{ab}A_\mu ^aF_{\nu \rho }^b,$$
(34)
where $`k_{ab}`$ is a non-degenerate symmetric and constant matrix, while $`F_{\nu \rho }^b=_\nu A_\rho ^b_\rho A_\nu ^b_{[\nu }A_{\rho ]}^b`$. Performing the canonical analysis and eliminating the second-class constraints (the independent variables are $`A_0^a`$, $`\pi _a^0`$ and $`A_k^a`$), we infer the first-class constraints $`G_{1a}\pi _a^00`$, $`G_{2a}\frac{1}{2}\epsilon ^{0ik}k_{ab}F_{ik}^b0`$ and the first-class Hamiltonian $`H_0=2d^2xA_0^aG_{2a}`$. The non-vanishing fundamental Dirac brackets read as $`[A_0^a,\pi _b^0]=\delta _b^a`$, $`[A_k^a,A_j^b]=\frac{1}{2}\epsilon _{0kj}k^{ab}`$, hence the BRST charge takes the simple form
$$\stackrel{(0)}{\mathrm{\Omega }}=d^2x(\pi _a^0\eta _1^a\frac{1}{2}\epsilon ^{0ik}k_{ab}F_{ik}^b\eta _2^a),$$
(35)
where $`k^{ab}`$ is the inverse of $`k_{ab}`$, and $`(\eta _1^a,\eta _2^a)`$ stand for the fermionic ghost number one ghosts. Thus, the BRST operator $`\stackrel{(0)}{s}`$ splits as $`\stackrel{(0)}{s}=\delta +\gamma `$, where $`\delta `$ is the Koszul-Tate differential and $`\gamma `$ represents the longitudinal exterior derivative along the gauge orbits. Then, we have
$$\delta A_0^a=0,\delta \pi _a^0=0,\delta A_k^a=0,\delta \eta _1^a=\delta \eta _2^a=0,$$
(36)
$$\delta 𝒫_{1a}=\pi _a^0,\delta 𝒫_{2a}=\frac{1}{2}\epsilon ^{0ik}k_{ab}F_{ik}^b,$$
(37)
$$\gamma A_0^a=\eta _1^a,\gamma \pi _a^0=0,\gamma A_k^a=\frac{1}{2}_k\eta _2^a,\gamma \eta _1^a=\gamma \eta _2^a=0,$$
(38)
$$\gamma 𝒫_{1a}=\gamma 𝒫_{2a}=0.$$
(39)
Now, we solve the former equation from (32). In view of this, we develop $`\stackrel{(1)}{\omega }`$ accordingly the antighost number
$$\stackrel{(1)}{\omega }=\underset{0}{\overset{(1)}{\omega }}+\underset{1}{\overset{(1)}{\omega }}+\mathrm{}+\underset{\mathrm{\Delta }}{\overset{(1)}{\omega }},antigh\left(\underset{\mathrm{\Delta }}{\overset{(1)}{\omega }}\right)=\mathrm{\Delta },gh\left(\underset{\mathrm{\Delta }}{\overset{(1)}{\omega }}\right)=1,$$
(40)
where the last term in (40) can be assumed to be annihilated by $`\gamma `$. As $`pgh\left(\underset{\mathrm{\Delta }}{\overset{(1)}{\omega }}\right)=\mathrm{\Delta }+1`$, we can represent $`\underset{\mathrm{\Delta }}{\overset{(1)}{\omega }}`$ under the form
$$\underset{\mathrm{\Delta }}{\overset{(1)}{\omega }}=\mu _{a_1\mathrm{}a_{\mathrm{\Delta }+1}}\eta _2^{a_1}\mathrm{}\eta _2^{a_{\mathrm{\Delta }+1}}.$$
(41)
With this choice, it results that the $`\gamma `$-invariant coefficient $`\mu _{a_1\mathrm{}a_{\mathrm{\Delta }+1}}`$ belongs to $`H_\mathrm{\Delta }\left(\delta |\stackrel{~}{d}\right)`$, i.e., is solution to the equation
$$\delta \mu _{a_1\mathrm{}a_{\mathrm{\Delta }+1}}+_kb_{a_1\mathrm{}a_{\mathrm{\Delta }+1}}^k=0,$$
(42)
for some $`b_{a_1\mathrm{}a_{\mathrm{\Delta }+1}}^k`$, where $`\stackrel{~}{d}=dx^i_i`$. Using the result from adapted to the Hamiltonian context, it follows that $`H_\mathrm{\Delta }\left(\delta |\stackrel{~}{d}\right)`$ vanish for $`\mathrm{\Delta }2`$ , hence $`\stackrel{(1)}{\omega }=\underset{0}{\overset{(1)}{\omega }}+\underset{1}{\overset{(1)}{\omega }}`$, with $`\underset{1}{\overset{(1)}{\omega }}=\frac{1}{2}\mu _{ab}\eta _2^a\eta _2^b`$, where $`\mu _{ab}`$ from $`H_1\left(\delta |\stackrel{~}{d}\right)`$. A general representative of $`H_1\left(\delta |\stackrel{~}{d}\right)`$ is of the type $`\mu _{ab}=C_{ab}^c𝒫_{2c}`$, where $`C_{ab}^c`$ are some constants, antisymmetric in the lower indices, $`C_{ab}^c=C_{ba}^c`$. The necessity for $`C_{ab}^c`$ to be constant results from the equation that must be satisfied by $`\mu _{ab}`$, namely, $`\delta \mu _{ab}=_k\left(C_{ab}^c\epsilon ^{0kj}k_{cd}A_j^d\right)`$. In this way, we obtained that $`\underset{1}{\overset{(1)}{\omega }}=\frac{1}{2}C_{ab}^c𝒫_{2c}\eta _2^a\eta _2^b`$. The former equation from (32) at antighost number zero reads as $`\delta \underset{1}{\overset{(1)}{\omega }}+\gamma \underset{0}{\overset{(1)}{\omega }}=_km^k`$, which further yields $`\underset{0}{\overset{(1)}{\omega }}=C_{ad}^ck_{cb}\epsilon ^{0kj}A_k^aA_j^d\eta _2^b`$. In this manner, we inferred $`\stackrel{(1)}{\omega }=C_{ab}^c(\frac{1}{2}𝒫_{2c}\eta _2^a\eta _2^b+k_{cd}\epsilon ^{0kj}A_k^aA_j^b\eta _2^d)`$. Simple computation leads to
$`[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]={\displaystyle }d^2x({\displaystyle \frac{1}{3}}C_{[nc}^mC_{ab]}^c𝒫_{2m}\eta _2^a\eta _2^b\eta _2^n`$
$`\epsilon ^{0ij}k_{ad}C_{[bc}^dC_{ne]}^c\eta _2^a\eta _2^bA_i^nA_j^e).`$ (43)
The last relation shows that $`[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]`$ cannot be written like a $`\stackrel{(0)}{s}`$-exact modulo $`\stackrel{~}{d}`$ local functional. For this reason it is necessary to have $`[\stackrel{(1)}{\mathrm{\Omega }},\stackrel{(1)}{\mathrm{\Omega }}]=0`$. This condition takes place if and only if $`C_{[nc}^mC_{ab]}^c=0`$, so if and only if the constants verify the Jacobi identity. This further implies $`\stackrel{(k)}{\mathrm{\Omega }}=0`$$`k2`$. The deformed BRST charge takes the final form
$`\mathrm{\Omega }={\displaystyle }d^2x(\pi _a^0\eta _1^a\epsilon ^{0ik}k_{ca}({\displaystyle \frac{1}{2}}F_{ik}^cgC_{bd}^cA_i^bA_k^d)\eta _2^a+`$
$`{\displaystyle \frac{1}{2}}gC_{ab}^c𝒫_{2c}\eta _2^a\eta _2^b),`$ (44)
so it is clearly a local functional.
Now, we derive the deformed BRST-invariant Hamiltonian. The BRST-invariant Hamiltonian for the free theory is given by $`\underset{B}{\overset{(0)}{H}}=H_0+2d^2x\eta _1^a𝒫_{2a}`$. Consequently, we find
$$[\underset{B}{\overset{(0)}{h}},\stackrel{(1)}{\omega }]=2C_{ab}^ck_{cd}\epsilon ^{0ij}A_j^b\left(\eta _1^dA_i^a+\eta _2^d_iA_0^a\right)2C_{ab}^c𝒫_{2c}\eta _2^a\eta _1^b.$$
(45)
Then, the solution of the latter equation in (32) reads as
$$\underset{B}{\overset{(1)}{h}}=2C_{ab}^c(k_{cd}\epsilon ^{0ij}A_0^dA_i^aA_j^b+A_0^b𝒫_{2c}\eta _2^a).$$
(46)
Straightforward computation leads to $`[\underset{B}{\overset{(1)}{H}},\stackrel{(1)}{\mathrm{\Omega }}]=0`$, hence the latter equation from (33) is satisfied with $`\underset{B}{\overset{(2)}{h}}=0`$. Therefore, the higher-order deformation equations for the BRST-invariant Hamiltonian are verified with $`\underset{B}{\overset{(3)}{H}}=\underset{B}{\overset{(4)}{H}}=\mathrm{}=0`$. Thus, the complete deformed BRST invariant Hamiltonian reads as
$$H_B=2d^2x\left(A_0^a\epsilon ^{0ik}k_{ca}\left(\frac{1}{2}F_{ik}^cgC_{bd}^cA_i^bA_k^d\right)+\left(\eta _1^agC_{cb}^aA_0^b\eta _2^c\right)𝒫_{2a}\right),$$
(47)
and is a local functional, too. With the help of (4) and (47) we identify the new gauge theory. From the antighost-independent terms in (4) we observe that the deformation of the BRST charge implies the deformed first-class constraints
$$\gamma _{2a}\epsilon ^{0ik}k_{ca}\left(\frac{1}{2}F_{ik}^cgC_{bd}^cA_i^bA_k^d\right)0,$$
(48)
the remaining constraints being undeformed. The term $`\frac{1}{2}gC_{ab}^c𝒫_{2c}\eta _2^a\eta _2^b`$ shows that the new constraint functions form a Lie algebra, i.e.,
$$[\gamma _{2a},\gamma _{2b}]=C_{ab}^c\gamma _{2c}.$$
(49)
On the other hand, the antighost-independent piece in (47)
$$H=2d^2xA_0^a\epsilon ^{0ik}k_{ca}\left(\frac{1}{2}F_{ik}^cgC_{bd}^cA_i^bA_k^d\right),$$
(50)
is precisely the first-class Hamiltonian of the deformed theory. The components linear in the antighosts from (47) indicate that the Dirac brackets among the new first-class Hamiltonian and the new constraint functions are modified as $`[H,\gamma _{2a}]=C_{ab}^cA_0^b\gamma _{2c}`$. Thus, the resulting first-class theory is nothing but the nonabelian version of the Chern-Simons model in three dimensions, described by the local Lagrangian action
$$\overline{S}_0\left[A_\mu ^a\right]=d^3x\epsilon ^{\mu \nu \rho }A_\mu ^a\left(\frac{1}{2}k_{ab}F_{\nu \rho }^b\frac{2}{3}gC_{abc}A_\nu ^bA_\rho ^c\right),$$
(51)
where $`C_{abc}=C_{[bc}^dk_{a]d}`$. As the first-class constraints generate gauge transformations, from the deformations (48) and (49) we can conclude that the added interactions involved with (50) modify both the gauge transformations and their gauge algebra.
## 5 Conclusion
To conclude with, in this letter we have presented a Hamiltonian BRST approach to the construction of consistent interactions among fields with gauge freedom. Our procedure reformulates the problem of constructing Hamiltonian consistent interactions as a deformation problem of the BRST charge and BRST-invariant Hamiltonian of a given “free” theory. We have derived the general equations from the Hamiltonian BRST deformation method, and proved that they possess solution. Next, we have written down the local version of these equations and discussed on the locality of their solutions. Finally, the general theory was exemplified in the case of the Chern-Simons model in three dimensions. We think that our approach together with the general results in might be successfully applied to computing local BRST cohomologies for those theories whose Lagrangian version is more intricate than the Hamiltonian one.
## Acknowledgment
This work has been supported by a Romanian National Council for Academic Scientific Research (CNCSIS) grant. |
warning/0003/math0003161.html | ar5iv | text | # Factorization of Combinatorial 𝑅 matrices and Associated Cellular Automata
## 1 Introduction
The box-ball systems are important examples of soliton cellular automata. They are discrete dynamical systems whose time evolution is expressed as a certain motion of balls along the one dimensional array of boxes. Their integrability has been understood by the ultradiscretization of classical integrable systems (soliton equations). In the recent works <sup>,</sup> it was revealed that the box-ball systems may also be viewed as quantum integrable systems at $`q=0`$. Here by quantum integrable systems we mean the ones whose integrability is guaranteed by the Yang-Baxter equation , and $`q`$ is the deformation parameter in the relevant quantum group. In fact, the box-ball systems are identified with a $`q0`$ limit of some two-dimensional solvable vertex models, where the role of time evolution is played by the action of a row transfer matrix. The simplest example is the original Takahashi-Satsuma automaton , whose classical origin is the discrete Lotka-Volterra equation and the quantum origin is the fusion six-vertex model. Here is an example of the automaton time evolution.
$`\mathrm{}1112211211111111\mathrm{}`$
$`\mathrm{}1111122121111111\mathrm{}`$
$`\mathrm{}1111111212211111\mathrm{}`$
$`\mathrm{}1111111121122111\mathrm{}`$
One regards $`1`$ as an empty box and $`2`$ as a box containing a ball. At each time step one moves every ball once starting from the leftmost one. The rule is that the ball goes to the nearest right empty box. One easily finds that the sequence of $`l`$ balls propagate stably to the right with velocity $`l`$ unless it interacts with other balls. By regarding such patterns as (ultradiscrete) solitons, the above figure illustrates how the larger soliton overtakes the smaller one with a phase shift.
In terms of the fusion six vertex model at $`q=0`$, the above figure corresponds to the configuration:
The left and right boundaries are to be understood as $`1`$ or $`111`$ everywhere. This is a configuration of the fusion six vertex model in which the quantum space is spin $`1/2`$ ($`1`$ or $`2`$) and the auxiliary space is spin $`3/2`$ ($`111,112,122`$ or $`222`$). At $`q=0`$ only some selected vertex configurations have non-zero Boltzmann weights and the transfer matrix yields a deterministic evolution of the spins on one row to another. The vertex configurations in the above figure are the non-zero ones, and form an example of the combinatorial ($`q=0`$) $`R`$ matrix, which will be a main subject in this paper. Let $`T_M`$ denote the row transfer matrix at $`q=0`$ corresponding to the spin $`M/2`$ auxiliary space. The above example corresponds to $`T_3`$. Actually it can be shown that the ball-moving algorithm coincides with the action of $`T_M`$ with sufficiently large $`M`$. In the above example $`T_M=T_3`$ holds for any $`M3`$.
The coincidence of the ultradiscrete limit of soliton equations and the $`q0`$ limit of vertex models is an interesting phenomenon in various respects. From a statistical mechanical point of view, it roughly means that in those solvable vertex models, the profile of low-lying excitations over the ferromagnetic ground state at $`q=0`$ admits an exact description by (ultradiscrete) soliton equations. From a mathematical point of view it leads to a systematic generalization by means of the quantum affine algebras and the crystal base theory . See section 4 of the reference for several examples of the scattering of ultradiscrete solitons.
Now we turn to a general setting, where the six vertex model is replaced with a solvable vertex model associated with the quantum affine algebra $`U_q^{}(\widehat{𝔤}_n)`$. In this paper we treat the non-exceptional case $`\widehat{𝔤}_n=A_n^{(1)},A_{2n1}^{(2)},A_{2n}^{(2)},B_n^{(1)},C_n^{(1)},D_n^{(1)}`$ and $`D_{n+1}^{(2)}`$. The box-ball systems correspond to $`\widehat{𝔤}_n=A_n^{(1)}`$. The Boltzmann weights are trigonometric functions satisfying the Yang-Baxter equation. The row transfer matrix is specified by the auxiliary space $`V_M`$ and the quantum space $`\mathrm{}V_{l_j}V_{l_{j+1}}\mathrm{}`$. Here $`V_M`$ denotes the $`M`$-fold symmetric fusion of the vector representation. We suppose a ferromagnetic boundary condition, namely, the spins in the distance $`|j|1`$ are all equal to some prescribed element in $`V_{l_j}`$. At $`q=0`$ the transfer matrix yields a deterministic evolution of the spins on one row to another.
To analyze such a situation we invoke the crystal base theory . Let $`B_l`$ be the crystal of $`V_l`$. It is a finite set listed in Appendix A endowed with the action of Kashiwara operators $`\stackrel{~}{e}_i,\stackrel{~}{f}_i:B_lB_l\{0\}`$ for $`0in`$. Let $`\delta _l[a]B_l`$ be the special element as in (12). The states in the automaton are the elements $`\mathrm{}b_jb_{j+1}\mathrm{}\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}`$ obeying the boundary condition $`b_j=\delta _{l_j}[a_k],|j|1`$. Here $`a_k`$ is specified in Table II with (15), and $`k`$ is a label of the boundary condition at our disposal. Let $`T_M`$ denote the $`q=0`$ transfer matrix, or $`M`$-th time evolution in the automaton; $`T_M:\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}`$. The map $`T_M(\mathrm{}b_jb_{j+1}\mathrm{})=\mathrm{}b_j^{}b_{j+1}^{}\mathrm{}`$ is induced by the isomorphism of crystals $`R:B_MB\stackrel{}{}BB_M`$ with $`B=\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}`$ according to (36). We call $`R`$ the combinatorial $`R`$ matrix. It is obtained by successive applications of the elementary ones $`B_MB_{l_j}\stackrel{}{}B_{l_j}B_M`$. The boundary condition matches the known properties like $`\delta _M[a_k]\delta _{l_j}[a_k]\stackrel{}{}\delta _{l_j}[a_k]\delta _M[a_k]`$. When $`M`$ gets large, $`T_M`$ stabilizes to a certain map, which we denote by $`T`$. In the box-ball system terminology ($`\widehat{𝔤}_n=A_n^{(1)}`$ case), this corresponds to the boxes with inhomogeneous capacities $`\{l_j\}`$ and the carrier with infinite capacity.
The main result of this note is Theorem 2, which states that the isomorphism $`R:B_MB\stackrel{}{}BB_M`$ with $`B=B_{l_1}\mathrm{}B_{l_N}`$ is expressed as
$$R=(\sigma _B\sigma )PS_{i_{k+d}}\mathrm{}S_{i_{k+2}}S_{i_{k+1}}$$
in the domain $`B_M[a_k]BB_MB`$ (16) if $`M`$ is sufficiently large. Here $`S_i`$ is the Weyl group operator (8) acting on $`B_MB`$. $`P(uv)=vu`$ is the transposition, $`\sigma `$ and $`\sigma _B`$ are the operators corresponding to the Dynkin diagram automorphism described around Proposition 1. The above result on $`R`$ reveals the factorization of the time evolution in the automaton: (Corollary 14)
$$T=\sigma _BS_{i_{k+d}}\mathrm{}S_{i_{k+2}}S_{i_{k+1}},$$
where $`S_i`$ here is the one acting on $`\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}`$. See the explanation after Corollary 14. Such a decomposition is by no means evident from the defining relation (36). Note that $`T`$ is a translation in the sense that the product $`\sigma r_{i_{k+d}}\mathrm{}r_{i_{k+1}}`$ ($`r_i`$ is a simple reflection) is so in the extended affine Weyl group, if $`\sigma `$ is interpreted as the Diagram automorphism acting on the weight lattice. According to the factorization, one can consider a finer time evolution $`𝒯_m`$ (37) that includes the original one as $`𝒯_{k+td}=T^t(t0)`$. For $`\widehat{𝔤}_n=A_n^{(1)}`$, the change from $`𝒯_m(p)`$ to $`𝒯_{m+1}(p)`$ agrees with the original ball-moving algorithm in the box-ball systems, where one touches only the balls with a fixed color. In particular when $`l_j=1`$, our Definition 11 provides a representation theoretical interpretation of the earlier observation. For $`\widehat{𝔤}_nA_n^{(1)}`$, the automaton corresponding to $`l_j=1`$ with $`a_k=1`$ has been introduced previously . The formula (37) in principle provides a simple algorithm to compute the refined time evolution for general $`ł_j`$ and $`a_k`$ in an analogous way to the ball-moving procedure for the $`A_n^{(1)}`$ case. The data $`d_1`$ and $`i_kI`$ are specified in Table II. Curiously, they have stemmed from the study of Demazure crystals . It will be interesting to investigate the present result in the light of the works .
## 2 Crystals
Let $`\widehat{𝔤}_n=A_n^{(1)}(n1),A_{2n1}^{(2)}(n3),A_{2n}^{(2)}(n2),B_n^{(1)}(n3),C_n^{(1)}(n2),D_n^{(1)}(n4)`$ and $`D_{n+1}^{(2)}(n2)`$. For each $`\widehat{𝔤}_n`$ and $`l_1`$, the $`U_q^{}(\widehat{𝔤}_n)`$ crystal $`B_l`$ has been constructed except for $`C_n^{(1)}`$ with $`l`$ odd. As for $`C_n^{(1)}`$, $`B_l`$ here is $`B^{1,l}`$ in the paper . The finite set $`B_l`$ and the actions of $`\stackrel{~}{e}_i,\stackrel{~}{f}_i:B_lB_l\{0\}`$ for $`iI=\{0,1,\mathrm{},n\}`$ (crystal structure) have been defined. We employ the same notation as and quote the set $`B_l`$ in Appendix A. In particular for $`\widehat{𝔤}_nA_n^{(1)}`$, we will always assume the convention
$$x_a=\overline{x}_i,\overline{x}_a=x_i,\overline{a}=i\text{if}a=\overline{i},\mathrm{\hspace{0.33em}1}in.$$
Let us recall some other notations and the tensor product rule.
$$\phi _i(b)=\mathrm{max}\{j\stackrel{~}{f}_i^jb0\},\epsilon _i(b)=\mathrm{max}\{j\stackrel{~}{e}_i^jb0\}.$$
For two crystals $`B`$ and $`B^{}`$, the tensor product $`BB^{}`$ is defined as the set $`BB^{}=\{b_1b_2b_1B,b_2B^{}\}`$ with the actions of $`\stackrel{~}{e}_i`$ and $`\stackrel{~}{f}_i`$ specified by
$`\stackrel{~}{e}_i(b_1b_2)`$ $`=`$ $`\{\begin{array}{cc}\stackrel{~}{e}_ib_1b_2\hfill & \text{ if }\phi _i(b_1)\epsilon _i(b_2)\hfill \\ b_1\stackrel{~}{e}_ib_2\hfill & \text{ if }\phi _i(b_1)<\epsilon _i(b_2),\hfill \end{array}`$ (3)
$`\stackrel{~}{f}_i(b_1b_2)`$ $`=`$ $`\{\begin{array}{cc}\stackrel{~}{f}_ib_1b_2\hfill & \text{ if }\phi _i(b_1)>\epsilon _i(b_2)\hfill \\ b_1\stackrel{~}{f}_ib_2\hfill & \text{ if }\phi _i(b_1)\epsilon _i(b_2).\hfill \end{array}`$ (6)
Here $`0b`$ and $`b0`$ are understood to be $`0`$. Consequently one has
$$\begin{array}{cc}\phi _i(b_1b_2)\hfill & =\phi _i(b_2)+(\phi _i(b_1)\epsilon _i(b_2))_+,\hfill \\ \epsilon _i(b_1b_2)\hfill & =\epsilon _i(b_1)+(\epsilon _i(b_2)\phi _i(b_1))_+,\hfill \end{array}$$
(7)
where the symbol $`(x)_+`$ is defined by $`(x)_+=\mathrm{max}(x,0)`$. The Weyl group operator $`S_i(iI)`$ is defined by
$$S_ib=\{\begin{array}{cc}\stackrel{~}{f}_i^{\phi _i(b)\epsilon _i(b)}b\hfill & \text{ if }\phi _i(b)\epsilon _i(b),\hfill \\ \stackrel{~}{e}_i^{\epsilon _i(b)\phi _i(b)}b\hfill & \text{ if }\phi _i(b)\epsilon _i(b).\hfill \end{array}$$
(8)
$`S_i`$ satisfies the Coxeter relations . For two crystals $`B`$ and $`B^{}`$, we let $`P:BB^{}B^{}B`$ denote the transposition $`P(uv)=vu`$. It is easy to check
$$S_iP=PS_i\text{ for any }iI.$$
(9)
For two crystals of the form $`B=B_{l_1}\mathrm{}B_{l_N}`$ and $`B^{}=B_{l_1^{}}\mathrm{}B_{l_N^{}^{}}`$, the tensor products $`BB^{}`$ and $`B^{}B`$ are isomorphic, i.e., they have the same crystal structure. The isomorphism $`R:BB^{}\stackrel{}{}B^{}B`$ is called the combinatorial $`R`$ matrix . (In this paper we do not consider the energy associated with $`R`$.) It is obtained by a successive application of the elementary ones $`B_{l_i}B_{l_j^{}}\stackrel{}{}B_{l_j^{}}B_{l_i}`$. We will use the same symbols $`\stackrel{~}{e}_i,\stackrel{~}{f}_i,\epsilon _i,\phi _i,S_i,P`$ and $`R`$ irrespective of the crystals that they act.
Let $`\mathrm{\Lambda }_i`$ denote a fundamental weight and let $`P_{cl}=_{iI}\mathrm{\Lambda }_i`$ be the classical weight lattice. (See Section 3.1 of the paper for a precise treatment.) We define a linear map $`\sigma :P_{cl}P_{cl}`$ as in the rightmost column of Table I. It is a Dynkin diagram automorphism. When $`\widehat{𝔤}_nC_n^{(1)}`$ (resp. $`\widehat{𝔤}_n=C_n^{(1)}`$), it agrees with the one introduced after Corollary 4.6.3 of the paper with $`B=B_l`$ (resp. $`B=B_{2l}`$) for any $`l`$.
We also let $`\sigma `$ act on the index set $`I`$ by the rule $`i^{}=\sigma (i)\sigma (\mathrm{\Lambda }_i)=\mathrm{\Lambda }_i^{}`$. For $`A_n^{(1)}`$, the letters $`a,a1`$ should be interpreted mod $`n+1`$.
We also introduce the data $`d_{>0}`$, and the sequences $`i_d,\mathrm{},i_1I`$ and $`a_d,\mathrm{},a_0`$ for each algebra as in Table II.
Here $`a_k`$’s are taken from the letters appearing in the description of the crystals $`B_l`$ in Appendix A. In the third and the fourth columns of Table II, the symbols $`\left\{\genfrac{}{}{0pt}{}{0,1}{1,0}\right\}`$ and $`\left\{\genfrac{}{}{0pt}{}{1}{\overline{1}}\right\}`$ mean that either the simultaneous upper choice or the simultaneous lower choice are allowed.
###### Proposition 1.
For any $`l_{>0}`$ the operator $`\sigma ^{}:=S_{i_1}\mathrm{}S_{i_d}:B_lB_l`$ is a bijection having the properties:
$$\sigma ^{}\stackrel{~}{f}_i=\stackrel{~}{f}_{\sigma (i)}\sigma ^{},\sigma ^{}\stackrel{~}{e}_i=\stackrel{~}{e}_{\sigma (i)}\sigma ^{}.$$
Moreover, the action of $`\sigma ^{}`$ is explicitly given by the second column of Table I.
The second column of Table I specifies the transformation of those letters labeling the elements of $`B_l`$. See Appendix A. For example in $`A_n^{(1)}`$ case, $`\sigma ^{}((x_1,\mathrm{},x_{n+1}))=(x_2,\mathrm{},x_{n+1},x_1)`$ in terms of the coordinates. Similarly for $`A_{2n1}^{(2)}`$, $`\sigma ^{}((x_1,\mathrm{},x_n,\overline{x}_n,\mathrm{},\overline{x}_1))=(\overline{x}_1,x_2,\mathrm{},x_n,\overline{x}_n,\mathrm{},\overline{x}_2,x_1)`$. The proposition implies $`\sigma ^{}S_i=S_{\sigma (i)}\sigma ^{}`$, from which the alternative expression $`\sigma ^{}=S_{i_{k+1}}\mathrm{}S_{i_d}S_{\sigma ^1(i_1)}\mathrm{}S_{\sigma ^1(i_k)}`$ is available for any $`0kd`$ due to $`S_i^2=id`$. We identify $`\sigma `$ with $`\sigma ^{}`$, thereby extend the definition of its domain. Namely, we also let $`\sigma `$ act on $`B_l`$ for any $`l`$ via $`\sigma =S_{i_{k+1}}\mathrm{}S_{i_d}S_{\sigma ^1(i_1)}\mathrm{}S_{\sigma ^1(i_k)}`$. (We do not exhibit $`l`$.) The result is independent of $`0kd`$ and enjoys the properties $`\sigma \stackrel{~}{f}_i=\stackrel{~}{f}_{\sigma (i)}\sigma `$ and $`\sigma \stackrel{~}{e}_i=\stackrel{~}{e}_{\sigma (i)}\sigma `$. Proposition 1 was known for some $`k`$. For the tensor product crystal $`B=B_{l_1}\mathrm{}B_{l_N}`$, we write $`\sigma _B=\sigma \mathrm{}\sigma :BB`$, where $`\sigma `$ on the right side acts on each component $`B_{l_j}`$ of the tensor product according to the above rule $`\sigma =S_{i_{k+1}}\mathrm{}S_{i_d}S_{\sigma ^1(i_1)}\mathrm{}S_{\sigma ^1(i_k)}`$. Obviously one has $`\sigma _B\stackrel{~}{f}_i=\stackrel{~}{f}_{\sigma (i)}\sigma _B`$ and $`\sigma _B\stackrel{~}{e}_i=\stackrel{~}{e}_{\sigma (i)}\sigma _B`$ on $`B`$, therefore
$$\sigma _BS_i=S_{\sigma (i)}\sigma _B\text{ on }B=B_{l_1}\mathrm{}B_{l_N}.$$
(10)
The combinatorial $`R`$ matrix $`R:BB^{}\stackrel{}{}B^{}B`$ satisfies
$$R(\sigma _B\sigma _B^{})=(\sigma _B^{}\sigma _B)R.$$
(11)
When acting on crystals, $`\sigma `$ without an index shall always act on a single crystal $`B_l`$ with some $`l`$.
For each $`a\{1,\mathrm{},n+1\}(\widehat{𝔤}_n=A_n^{(1)})`$, $`\{1,\mathrm{},n,\overline{n},\mathrm{},\overline{1}\}(\widehat{𝔤}_nA_n^{(1)})`$, we set
$$\delta _l[a]=(x_a=l,x_a^{}=0\text{ for }a^{}a)B_l$$
(12)
with the notation in Appendix A. Using the crystal structure explicitly one easily finds
$`\delta _l[a_k]=S_{i_k}\left(\delta _l[a_{k1}]\right)=\stackrel{~}{e}_{i_k}^{\mathrm{max}}\left(\delta _l[a_{k1}]\right)1kd,`$ (13)
$`\stackrel{~}{e}_{i_d}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}u=\delta _l[a_d]\text{for any }uB_l,`$ (14)
$`\phi _{i_k}(\delta _l[a_{k1}])=0,`$
for any $`\widehat{𝔤}_n`$, where $`\stackrel{~}{e}_i^{\mathrm{max}}b=\stackrel{~}{e}_i^{\epsilon _i(b)}b`$. In particular (13) implies $`\delta _l[a_d]=\sigma ^1(\delta _l[a_0])`$ due to Proposition 1. Thus it is natural to extend the definition of $`i_kI`$ and $`a_k`$ to all $`k`$ by
$$i_{k+d}=\sigma ^1(i_k),\delta _l[a_{k+d}]=\sigma ^1\left(\delta _l[a_k]\right).$$
(15)
This way of extending the index $`k`$ of $`a_k`$ is independent of $`l`$, and (13) also persists for all $`k`$. We have $`\{a_kk\}=\{1,\mathrm{},n+1\}`$ for $`\widehat{𝔤}_n=A_n^{(1)}`$ and $`\{1,\mathrm{},n,\overline{n},\mathrm{},\overline{1}\}`$ for $`\widehat{𝔤}_nA_n^{(1)}`$.
In this paper we will concern some asymptotic domain of the crystal $`B_M`$ when $`M`$ gets large. For $`a\{a_kk\}`$ and $`M1`$, we introduce the “domain” $`B_M[a]B_M`$ by
$$\begin{array}{cc}& \widehat{𝔤}_n=A_n^{(1)}:\hfill \\ & B_M[a]=\{(u_1,\mathrm{},u_{n+1})B_Mu_au_b\text{ for any }b\{1,\mathrm{},n+1\}\{a\}\},\hfill \\ & \widehat{𝔤}_nA_n^{(1)}:\hfill \\ & B_M[a]=\{(u_1,\mathrm{},\overline{u}_1)B_M|\begin{array}{c}u_a\overline{u}_a|u_b\overline{u}_b|\text{ for any }\hfill \\ b\{1,\mathrm{},n\},ba,\overline{a}\hfill \end{array}\}.\hfill \end{array}$$
(16)
In the rest of the paper we will have assertions under the conditions like $`cuc^{}BB_M[a]B^{}`$ with $`B`$ and $`B^{}`$ of the form $`B=B_{l_1}\mathrm{}B_{l_N}`$, $`B^{}=B_{l_1^{}}\mathrm{}B_{l_N^{}^{}}`$ ($`N,N^{}0`$). The mathematically invalid “definition” (16) should be understood that the associated assertions are valid on condition that $`M1`$ and the inequalities in (16) are satisfied. Thus for example it amounts to assuming the following:
$`\epsilon _{i_k}\phi _{i_k}\text{on }BB_M[a_{k1}]B^{},\phi _{i_k}\epsilon _{i_k}\text{on }BB_M[a_k]B^{},`$ (17)
$`\epsilon _{i_k}(ux)=\epsilon _{i_k}(u)\text{ for }uxB_M[a_k]B,`$ (18)
$`\phi _{i_{k+1}}(xu)=\phi _{i_{k+1}}(u)\text{ for }xuBB_M[a_k],`$ (19)
$`\text{If }ux\stackrel{}{}yv\text{ then }uxB_M[a]ByvBB_M[a].`$ (20)
In the above (17) can be checked by using the explicit formula for $`\phi _i`$ and $`\epsilon _i`$. (18) and (19) follow directly from (17) and (7). By the weight reason (20) is obvious. Moreover we may effectively treat as
$$S_{i_k}(BB_M[a_{k1}]B^{})=BB_M[a_k]B^{},$$
(21)
for any $`k`$.
## 3 Combinatorial $`R`$ matrices
Let $`B`$ be any crystal of the form $`B=B_{l_1}\mathrm{}B_{l_N}`$. Our goal in this section is to prove
###### Theorem 2.
When $`M`$ is sufficiently large, the combinatorial $`R`$ matrix giving the isomorphism $`R:B_MB\stackrel{}{}BB_M`$ is expressed as
$$R=(\sigma _B\sigma )PS_{i_{k+d}}\mathrm{}S_{i_{k+2}}S_{i_{k+1}}$$
on $`B_M[a_k]BB_MB`$ for any $`k`$.
###### Definition 3.
Let $`B`$ be any crystal of the form $`B=B_{l_1}\mathrm{}B_{l_N}`$. Given $`uxB_M[a_0]B`$ and $`yvBB_M[a_0]`$ we set
$`u^{|k}x^{|k}`$ $`=S_{i_k}\mathrm{}S_{i_1}(ux)B_M[a_k]B,`$
$`y^{k|}v^{k|}`$ $`=S_{i_k}\mathrm{}S_{i_1}(yv)BB_M[a_k],`$
for $`0kd`$.
It should be noted that $`u^{|k}`$ for example is not defined solely from $`u`$ but only with the other element $`x`$.
It suffices to show Theorem 2 for $`k=0`$. To see this, let $`uxB[a_0]B,\mathrm{\hspace{0.17em}1}kd`$ and assume $`ux\stackrel{}{}(\sigma _B\sigma )PS_{i_d}\mathrm{}S_{i_1}(ux)`$ under the isomorphism $`B_MBBB_M`$ with $`M`$ sufficiently large. Multiplying $`S_{i_k}\mathrm{}S_{i_1}`$ on the both sides and using (10), one gets
$$u^{|k}x^{|k}\stackrel{}{}(\sigma _B\sigma )PS_{\sigma ^1(i_k)}\mathrm{}S_{\sigma ^1(i_1)}S_{i_d}\mathrm{}S_{i_{k+1}}(u^{|k}x^{|k}).$$
(22)
In view of (15) and (21) this proves $`1kd`$ case. To see it for the other $`k`$, multiply $`(\sigma ^m\sigma _B^m)`$ on the left side and $`(\sigma _B^m\sigma ^m)`$ on the right side of (22) for an integer $`m`$ and use (10) and (11).
Henceforth we shall concentrate on the $`k=0`$ case of Theorem 2 in the rest of this section. Suppose $`B_M[a_0]Bux\stackrel{}{}yvBB_M[a_0]`$ under the isomorphism $`B_MBBB_M`$ with $`M`$ sufficiently large. Then the assertion of Theorem 2 with $`k=0`$ is equivalent to
$`v`$ $`=\sigma u^{|d},`$ (23)
$`y`$ $`=\sigma _Bx^{|d}.`$ (24)
We shall prove these relations separately. In our approach, (23) can be verified directly for any choice $`B=B_{l_1}\mathrm{}B_{l_N}`$. On the other hand, as for (24), we first deal with $`N=1`$ case and derive $`N`$ general case based on it.
Let us first treat (23).
###### Definition 4.
For any crystal $`B`$ we introduce
$`t:B`$ $`_0^d`$
$`b`$ $`(t_1(b),\mathrm{},t_d(b))`$
$`t_k(b)`$ $`=\phi _{i_k}(\stackrel{~}{e}_{i_{k1}}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}b).`$
Here $`i_k`$’s are those in Table II. It is easy to calculate $`t`$ explicitly for $`B=B_l`$. For the element $`b=(x_1,\mathrm{},x_{n+1})`$ for $`A_n^{(1)}`$ ($`b=(x_1,\mathrm{},\overline{x}_1)`$ for $`\widehat{𝔤}_nA_n^{(1)}`$), the result is summarized in
###### Lemma 5.
The map $`t:B_l_0^d`$ has the form:
$`t(b)`$ $`=(x_n,\mathrm{},x_1)\text{for }A_n^{(1)},`$
$`=(x_n,\mathrm{},x_3,x_2,\left\{{\displaystyle \genfrac{}{}{0pt}{}{x_1,\overline{x}_1}{\overline{x}_1,x_1}}\right\},\overline{x}_2,\mathrm{},\overline{x}_{n1})\text{for }A_{2n1}^{(2)},`$
$`=(x_n,\mathrm{},x_1,x_0,\overline{x}_1,\mathrm{},\overline{x}_{n1})\text{for }A_{2n}^{(2)},`$
$`=(2x_n+x_0,x_{n1},\mathrm{},x_2,\left\{{\displaystyle \genfrac{}{}{0pt}{}{x_1,\overline{x}_1}{\overline{x}_1,x_1}}\right\},\overline{x}_2,\mathrm{},\overline{x}_{n1})\text{for }B_n^{(1)},`$
$`=(x_n,\mathrm{},x_1,x_0,\overline{x}_1,\mathrm{},\overline{x}_{n1})\text{for }C_n^{(1)},`$
$`=(x_n+x_{n1},x_{n2},\mathrm{},x_2,\left\{{\displaystyle \genfrac{}{}{0pt}{}{x_1,\overline{x}_1}{\overline{x}_1,x_1}}\right\},\overline{x}_2,\mathrm{},\overline{x}_{n2},\overline{x}_{n1}+x_n)\text{for }D_n^{(1)},`$
$`=(2x_n+x_0,x_{n1},\mathrm{},x_1,x_{\mathrm{}},\overline{x}_1,\mathrm{},\overline{x}_{n1})\text{for }D_{n+1}^{(2)}.`$
See Appendix A for the notation $`x_0`$ in $`A_{2n}^{(2)},C_n^{(1)}`$ and $`x_{\mathrm{}}`$ in $`D_{n+1}^{(2)}`$. The upper and lower choices correspond to those in Table II. From Lemma 5 we derive a useful fact.
###### Proposition 6.
The map $`t:B_l_0^d`$ is injective for any $`l_1`$.
###### Lemma 7.
Under Definition 3 one has
$$\epsilon _{i_k}(u^{|k})=\phi _{i_k}(u^{|k1}x^{|k1}),1kd.$$
(25)
Proof. Definition 3 tells $`u^{|k}x^{|k}=S_{i_k}(u^{|k1}x^{|k1})`$. Since $`u^{|k1}x^{|k1}B_M[a_{k1}]B`$, this $`S_{i_k}`$ acts as
$`u^{|k}x^{|k}`$ $`=\stackrel{~}{e}_{i_k}^qu^{|k1}\stackrel{~}{e}_{i_k}^q^{}x^{|k1},`$ (26)
$`q`$ $`=\epsilon _{i_k}(u^{|k1})\phi _{i_k}(x^{|k1})\left(\phi _{i_k}(u^{|k1})\epsilon _{i_k}(x^{|k1})\right)_+,`$
$`q^{}`$ $`=\left(\epsilon _{i_k}(x^{|k1})\phi _{i_k}(u^{|k1})\right)_+.`$
Thus we have $`\epsilon _{i_k}(u^{|k})=\epsilon _{i_k}(u^{|k1})q\stackrel{(\text{7})}{=}\phi _{i_k}(u^{|k1}x^{|k1})`$. $`\mathrm{}`$
For integers $`p,q`$ depending on $`M`$ in general, we write $`pq`$ to mean $`p<q`$ or $`pq=M`$-independent for $`M1`$.
###### Lemma 8.
Let $`uB_M[a_0]B_M`$ for sufficiently large $`M`$. For $`(c_1,\mathrm{},c_d)_0^d`$, define $`u^k(0kd)`$ by $`u^k=\stackrel{~}{e}_{i_k}^{\epsilon _{i_k}(u^{k1})c_k}u^{k1}(1kd)`$ and $`u^0=u`$. Suppose $`c_jt_j(u)`$ for all $`1jd`$. Then the following hold for $`1kd`$:
$`\phi _{i_k}(u^{k1})=t_k(u),`$ (27)
$`\epsilon _{i_k}(u^k)=c_k,`$ (28)
$`t_k(\sigma u^d)=c_k.`$ (29)
Although $`u^k`$ depends on $`c_j`$’s, the right side of (27) is independent of them. Actually (28) is trivial. The other relations in the lemma can be verified with a direct calculation by using the crystal structure of $`B_M`$. Under Definition 3, Lemma 8 immediately leads to ($`1kd`$)
$`\phi _{i_k}(u^{|k1})=t_k(u),`$ (30)
$`\phi _{i_k}(v^{k1|})=t_k(v),`$ (31)
$`\epsilon _{i_k}(u^{|k})=t_k(\sigma u^{|d}).`$ (32)
The right sides of (30) and (31) are independent of $`x`$ and $`y`$ in Definition 3, respectively.
Proof of (23). Suppose $`B_M[a_0]Bux\stackrel{}{}yvBB_M[a_0]`$ under the isomorphism $`B_MBBB_M`$ with $`M`$ sufficiently large. We employ the notation in Definition 3. Applying $`S_{i_{k1}}\mathrm{}S_{i_1}`$ to the both sides of $`ux\stackrel{}{}yv`$, one gets $`u^{|k1}x^{|k1}\stackrel{}{}y^{k1|}v^{k1|}`$, therefore $`\phi _{i_k}(u^{|k1}x^{|k1})=\phi _{i_k}(y^{k1|}v^{k1|})`$. But
$`\phi _{i_k}(u^{|k1}x^{|k1})`$ $`\stackrel{(\text{25})}{=}\epsilon _{i_k}(u^{|k})\stackrel{(\text{32})}{=}t_k(\sigma u^{|d}),`$
$`\phi _{i_k}(y^{k1|}v^{k1|})`$ $`\stackrel{(\text{19})}{=}\phi _{i_k}(v^{k1|})\stackrel{(\text{31})}{=}t_k(v),`$
are valid for $`1kd`$, telling that $`t(v)=t(\sigma u^{|d})`$. Thus (23) follows from Proposition 6. $`\mathrm{}`$
Now we proceed to the proof of (24) with the simple choice $`B=B_l`$.
###### Lemma 9.
Suppose $`B_M[a_0]B_l\delta _M[a_0]z\stackrel{}{}\stackrel{~}{z}\stackrel{~}{u}B_lB_M[a_0]`$ under the isomorphism $`B_MB_lB_lB_M`$ with $`M`$ sufficiently large. Then we have $`t(\stackrel{~}{u})=t(z)`$.
Proof. Define $`u^{|k}z^{|k}`$ by Definition 3 starting from $`uz=u^{|0}z^{|0}=\delta _M[a_0]z`$. From (23) we already know that $`\stackrel{~}{u}=\sigma u^{|d}`$. Thus we have
$`t_k(\stackrel{~}{u})`$ $`=t_k(\sigma u^{|d})\stackrel{(\text{32})}{=}\epsilon _{i_k}(u^{|k})\stackrel{(\text{25})}{=}\phi _{i_k}(u^{|k1}z^{|k1})`$
$`\stackrel{(\text{7})}{=}\phi _{i_k}(z^{|k1})+\left(\phi _{i_k}(u^{|k1})\epsilon _{i_k}(z^{|k1})\right)_+`$
$`\stackrel{(\text{30})}{=}\phi _{i_k}(z^{|k1})+\left(t_k(u)\epsilon _{i_k}(z^{|k1})\right)_+.`$
Note from the explicit forms in Lemma 5 that $`t_k(u)=t_k(\delta _M[a_0])=0`$, from which $`t_k(\stackrel{~}{u})=\phi _{i_k}(z^{|k1})`$ follows. In view of $`\phi _{i_k}(u^{|k1})\stackrel{(\text{30})}{=}t_k(u)=0`$ and (17), it also follows that $`z^{|k}=\stackrel{~}{e}_{i_k}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}z`$. Therefore we conclude $`\phi _{i_k}(z^{|k1})=t_k(z)`$ from Definition 4. $`\mathrm{}`$
###### Lemma 10.
Given $`yvB_lB_M[a_0]`$, set
$$y^{(k)}v^{(k)}=\stackrel{~}{e}_{i_k}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}(yv)$$
for $`0kd`$. Then we have $`t(\sigma v^{(d)})=t(y)`$.
Proof. Since $`y^{(k)}=\stackrel{~}{e}_{i_k}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}y`$, we have
$$v^{(k)}=\stackrel{~}{e}_{i_k}^{\epsilon _{i_k}(v^{(k1)})\phi _{i_k}(y^{(k1)})}v^{(k1)}=\stackrel{~}{e}_{i_k}^{\epsilon _{i_k}(v^{(k1)})t_k(y)}v^{(k1)}.$$
Applying (29) we obtain $`t(\sigma v^{(d)})=t(y)`$. $`\mathrm{}`$
Proof of (24) for $`B=B_l`$. Suppose $`B_M[a_0]B_lux\stackrel{}{}yvB_lB_M[a_0]`$ under the isomorphism $`B_MB_lB_lB_M`$ with $`M`$ sufficiently large. We employ the notation in Definition 3. Applying $`(\sigma \sigma )\stackrel{~}{e}_{i_d}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}`$ to the both sides and using (14), one finds
$$B_M[a_0]B_l\delta _M[a_0]\sigma x^{|d}\stackrel{}{}(\sigma \sigma )\stackrel{~}{e}_{i_d}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}(yv)B_lB_M[a_0].$$
Setting
$`\delta _M[a_0]\sigma x^{|d}`$ $`\stackrel{}{}\stackrel{~}{z}\stackrel{~}{u}B_lB_M[a_0],`$
$`(\sigma \sigma )\stackrel{~}{e}_{i_d}^{\mathrm{max}}\mathrm{}\stackrel{~}{e}_{i_1}^{\mathrm{max}}(yv)`$ $`=(\sigma y^{(d)})(\sigma v^{(d)}),`$
we have $`\stackrel{~}{z}=\sigma y^{(d)}`$ and $`\stackrel{~}{u}=\sigma v^{(d)}`$. Thus $`t(\stackrel{~}{u})=t(\sigma v^{(d)})`$. But we know $`t(\stackrel{~}{u})=t(\sigma x^{|d})`$ from Lemma 9 and $`t(\sigma v^{(d)})=t(y)`$ from Lemma 10. Therefore $`t(\sigma x^{|d})=t(y)`$, compelling $`y=\sigma x^{|d}`$ by Proposition 6. $`\mathrm{}`$
To show (24) for the general choice $`B=B_{l_1}\mathrm{}B_{l_N}`$, we prepare
###### Definition 11.
For $`s,s^{}_0,b,b^{}B_l`$ and $`iI`$, we let the vertex diagram
denote the relations
$$b^{}=\stackrel{~}{e}_i^{(\epsilon _i(b)s)_+}b,s^{}=\phi _i(b)+(s\epsilon _i(b))_+.$$
Here $`l`$ can be any positive integer but we do not exhibit it in the diagram. The color $`iI`$ is attached to the horizontal line. The diagram should not be confused with the one representing the combinatorial $`R`$ matrix . Given $`i`$, $`(b^{},s^{})`$ is uniquely fixed from $`(b,s)`$. Thus for example the diagram
is uniquely determined from $`s_0,i`$ and $`b_1\mathrm{}b_N`$. From Definition 11 it implies
$$b_1^{}\mathrm{}b_N^{}=\stackrel{~}{e}_i^{\left(\epsilon _i(b_1\mathrm{}b_N)s_0\right)_+}(b_1\mathrm{}b_N).$$
(33)
Having established Theorem 2 for $`B=B_l`$, we already know that
$$B_M[a_0]B_lub\stackrel{}{}\sigma b^{|d}\sigma u^{|d}B_lB_M[a_0]$$
under the isomorphism $`B_MB_lB_lB_M`$ for $`M`$ sufficiently large, where $`u^{|k}b^{|k}=S_{i_k}\mathrm{}S_{i_1}(ub)(0kd)`$.
###### Proposition 12.
Under the above stated setting, the following diagram holds.
Proof. By Definition 11 we are to check
$`b^{|k}`$ $`=\stackrel{~}{e}_{i_k}^{\left(\epsilon _{i_k}(b^{|k1})t_k(u)\right)_+}b^{|k1},`$
$`t_k(\sigma u^{|d})`$ $`=\phi _{i_k}(b^{|k1})+\left(t_k(u)\epsilon _{i_k}(b^{|k1})\right)_+,`$
for $`1kd`$. To show the former, set $`x=b`$ and apply (30) in $`q^{}`$ appearing in (26). The left side of the latter reads
$`t_k(\sigma u^{|d})`$ $`\stackrel{(\text{32})}{=}\epsilon _{i_k}(u^{|k})\stackrel{(\text{25})}{=}\phi _{i_k}(u^{|k1}b^{|k1})`$
$`\stackrel{(\text{7})}{=}\phi _{i_k}(b^{|k1})+\left(\phi _{i_k}(u^{|k1})\epsilon _{i_k}(b^{|k1})\right)_+,`$
which is equal to the right side owing to (30). $`\mathrm{}`$
Proof of (24) for $`B=B_{l_1}\mathrm{}B_{l_N}`$. Given any $`x=b_1\mathrm{}b_NB`$ and $`uB_M[a_0]`$, set $`s_{0,k}=t_k(u)(1kd)`$. Let $`b_j^{|k}B_{l_j},s_{j,k}_0(1kd,1jN)`$ be the ones uniquely determined from the diagram:
By a repeated use of Proposition 12, one has
$`ub_1\mathrm{}b_N`$ $`\stackrel{}{}\sigma b_1^{|d}\sigma b_2^{|d}\mathrm{}\sigma b_N^{|d}v`$
$`=\left(\sigma _B\left(b_1^{|d}\mathrm{}b_N^{|d}\right)\right)v`$
under the isomorphism $`B_MBBB_M`$ with sufficiently large $`M`$. On the other hand we introduce $`u^{|k}x^{|k}:=S_{i_k}\mathrm{}S_{i_1}(ux)B_M[a_k]B`$. (Although not necessary here, $`vB_M[a_0]`$ is uniquely determined from $`t(v)=(s_{N,1},\mathrm{},s_{N,d})`$ by Proposition 6, and we already know that the result coincides with $`v=\sigma u^{|d}`$ from (23).) We are to show $`x^{|d}=b_1^{|d}\mathrm{}b_N^{|d}`$. In fact, $`x^{|k}=b_1^{|k}\mathrm{}b_N^{|k}`$ can be proved by induction on $`k`$ as follows. (We set $`b_1^{|0}\mathrm{}b_N^{|0}=b_1\mathrm{}b_N`$.) It is obvious for $`k=0`$. From (33) we have
$$b_1^{|k}\mathrm{}b_N^{|k}=\stackrel{~}{e}_{i_k}^{(ms_{0,k})_+}(b_1^{|k1}\mathrm{}b_N^{|k1}),$$
(34)
where $`m=\epsilon _{i_k}(b_1^{|k1}\mathrm{}b_N^{|k1})`$. On the other hand, $`x^{|k}`$ is determined from the recursion relation:
$$x^{|k}=\stackrel{~}{e}_{i_k}^{\left(\epsilon _{i_k}(x^{|k1})\phi _{i_k}(u^{|k1})\right)_+}x^{|k1},$$
(35)
because of $`\epsilon _{i_k}(u^{|k1})1`$. Note that $`\phi _{i_k}(u^{|k1})\stackrel{(\text{30})}{=}t_k(u)=s_{0,k}`$. Thus the two recursion relations (34) and (35) lead to $`x^{|k}=b_1^{|k}\mathrm{}b_N^{|k}`$ under the induction assumption $`x^{|k1}=b_1^{|k1}\mathrm{}b_N^{|k1}`$. $`\mathrm{}`$
We have finished the proof of (24) for any $`B=B_{l_1}\mathrm{}B_{l_N}`$, and thereby the proof of Theorem 2.
###### Example 13.
Consider $`\widehat{𝔤}_n=A_3^{(1)}(d=3)`$. The data in Table II reads $`i_ja_j4j`$ mod 4 with $`i_j\{0,1,2,3\}`$ and $`a_j\{1,2,3,4\}`$. To save the space the element $`(x_1,x_2,x_3,x_4)=(3,2,1,0)B_6`$ is denoted by $`111223`$ for example. Then one has $`\sigma (111223)=112444`$ according to Table I.
Let us take $`bc=111223344B_6B_3`$ and seek its image under the isomorphism $`R:B_6B_3\stackrel{}{}B_3B_6`$. It is known that $`R(bc)=223111344`$. This can be derived by taking $`N=1,M=6`$ and $`k=3`$ in Theorem 2, which reads $`R=(\sigma \sigma )PS_2S_3S_0`$. Namely we may regard $`bcB_M[1]B_3`$, where $`M=6`$ and $`x_1=3`$ for $`b`$ are already sufficiently large so that
$`bc=111223344`$ $`\stackrel{S_0}{}112234344`$
$`\stackrel{S_3}{}112234334`$
$`\stackrel{S_2}{}112224334`$
$`\stackrel{P}{}334112224\stackrel{\sigma \sigma }{}223111344=R(bc).`$
For comparison, take a smaller example $`R(11223344)=22311344B_3B_5`$. Theorem 2 is not applicable in this case under any choice of $`k`$ because $`S_i(11223344)=11223344`$ for any $`i\{0,1,2,3\}`$ and $`(\sigma \sigma )P(11223344)=2331124422311344`$.
## 4 Cellular Automata
The factorization of the combinatorial $`R`$ matrix shown in Theorem 2 induces the factorization of the time evolution of the associated cellular automaton. Consider the isomorphism
$$B_M(\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{})\stackrel{}{}(\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{})B_M$$
induced by the successive application of combinatorial $`R`$ matrix $`B_MB_{l_j}B_{j_j}B_M`$. We impose the boundary condition on $`b_jB_{l_j}`$ as $`b_j=\delta _{l_j}[a_k]`$ for $`|j|1`$, where the choice of $`k`$ is arbitrary. Assume the following properties:
(i) $`\delta _M[a_k]\delta _l[a_k]\stackrel{}{}\delta _l[a_k]\delta _M[a_k]\text{ for any }M,l,`$
(ii) $`u\delta _{l_j}[a_k]\delta _{l_{j+1}}[a_k]\mathrm{}\delta _{l_{j+N}}[a_k]\stackrel{}{}\stackrel{~}{b}_j\mathrm{}\stackrel{~}{b}_{j+N}\delta _M[a_k]`$
$`\text{for any }uB_M\text{ if }N\text{ is sufficiently large},`$
where $`\stackrel{~}{b}_{l_j}\mathrm{}\stackrel{~}{b}_{l_{j+N}}`$ is some element in $`B_{l_j}\mathrm{}B_{l_{j+N}}`$. ((i) is indeed valid by the weight reason.) Then under the isomorphism $`B_M(\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{})(\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{})B_M`$,
$$\delta _M[a_k](\mathrm{}b_jb_{j+1}\mathrm{})\stackrel{}{}(\mathrm{}b_j^{}b_{j+1}^{}\mathrm{})\delta _M[a_k]$$
(36)
is valid for some $`b_j^{}`$’s. One may regard the system as an automaton which undergoes the time evolution $`p=\mathrm{}b_jb_{j+1}\mathrm{}p^{}=\mathrm{}b_j^{}b_{j+1}^{}\mathrm{}`$. When $`M`$ gets large, the transformation $`pp^{}`$ stabilizes to a certain map, which we denote by $`T(p)=p^{}`$. By taking $`B=\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}`$ in Theorem 2 and using (10), (15) we obtain
###### Corollary 14.
Under the boundary condition $`b_j=\delta _{l_j}[a_k]B_{l_j}`$ for $`|j|1`$, the time evolution $`T`$ acts on $`B=\mathrm{}B_{l_j}B_{l_{j+1}}\mathrm{}`$ as
$$T^t=\{\begin{array}{cc}\sigma _B^tS_{i_{k+td}}\mathrm{}S_{i_{k+2}}S_{i_{k+1}}& \text{ if }t_0,\hfill \\ \sigma _B^tS_{i_{k+td+1}}\mathrm{}S_{i_{k1}}S_{i_k}& \text{ if }t_{<0}.\hfill \end{array}$$
Actually all the Weyl group operators $`S_{i_m}`$ for $`t>0`$ (resp. $`t<0`$) in the above act as $`S_{i_m}=\stackrel{~}{e}_{i_m}^{\mathrm{max}}`$ (resp. $`S_{i_m}=\stackrel{~}{f}_{i_m}^{\mathrm{max}}`$) on $`B`$ since they always hit such states $`pB`$ that $`\epsilon _{i_m}(p)1`$ (resp. $`\phi _{i_m}(p)1`$). Corollary 14 exhibits a factorization of the time evolution of the automaton having the background (vacuum) configuration $`\mathrm{}\delta _{l_j}[a_k]\delta _{l_{j+1}}[a_k]\mathrm{}`$ specified by $`k`$. When the partial factor $`S_{i_m}S_{i_{m1}}\mathrm{}S_{i_{k+1}}(k+dmk+1)`$ in $`T`$ is applied, the background, hence the boundary condition, changes into $`\mathrm{}\delta _{l_j}[a_m]\delta _{l_{j+1}}[a_m]\mathrm{}`$ according to (13). A generalization of $`T`$ that does not change the background in every intermediate step is constructed as follows:
$$\begin{array}{cc}𝒯_m\hfill & =\sigma _{B,k,m}S_{i_m}\mathrm{}S_{i_{k+2}}S_{i_{k+1}},mk+1,\hfill \\ \sigma _{B,k,m}\hfill & =\mathrm{}\sigma _{k,m}\sigma _{k,m}\mathrm{},\hfill \\ \sigma _{k,m}\hfill & =S_{i_{k+1}}S_{i_{k+2}}\mathrm{}S_{i_m},\hfill \end{array}$$
(37)
where $`\sigma _{k,m}`$ acts on each component $`B_{l_j}`$ of the tensor product. We understand $`𝒯_k=id`$. For any $`mk`$, the operator $`𝒯_m`$ retains the background in the original form $`\mathrm{}\delta _{l_j}[a_k]\delta _{l_{j+1}}[a_k]\mathrm{}`$ due to (13). Moreover from Corollary 14 we find that the original evolution under $`t`$-time step $`T^t`$ is included in $`\{𝒯_m\}`$ as $`T^t=𝒯_{k+td}(t0)`$. For $`A_n^{(1)}`$ the evolution of the state $`p`$ according to $`p=𝒯_k(p),𝒯_{k+1}(p),𝒯_{k+2}(p),\mathrm{}`$ agrees with that obtained by the ball-moving algorithm in the box-ball systems under a convention adjustment. In particular for $`A_1^{(1)}`$, the evolution rule in terms of the Dyck language essentially agrees with the crystal theoretic signature rule in applying $`S_i=\stackrel{~}{e}_i^{\mathrm{max}}`$ or $`\stackrel{~}{f}_i^{\mathrm{max}}(i=0,1)`$ explained in section 1.3 of the reference .
###### Example 15.
Consider $`\widehat{𝔤}_n=A_5^{(2)}(d=5)`$. To save the space the element $`(x_1,x_2,x_3,\overline{x}_3,\overline{x}_2,\overline{x}_1)=(2,1,2,1,0,1)B_7`$ is denoted by $`11233\overline{3}\overline{1}`$ for example. Let us take $`k=0`$, hence the initial background is $`\mathrm{}\delta _{l_j}[\overline{3}]\delta _{l_{j+1}}[\overline{3}]\mathrm{}`$. We employ $`i_5,\mathrm{},i_1=2,0,1,2,3`$ and $`a_5,\mathrm{},a_0=\overline{3},\overline{2},1,2,3,\overline{3}`$ correspondingly. Take
$$p=\mathrm{}\overline{3}\overline{3}1\overline{2}3\overline{3}\overline{1}2\overline{3}\overline{3}\overline{3}\overline{3}\mathrm{}\mathrm{}B_2B_2B_1B_2B_1B_2B_2\mathrm{},$$
where $``$ stands for $``$ and $`\mathrm{}`$ parts on the both ends represent the configuration identical with the background. Then the time evolution (36) is given by
$$T(p)=\mathrm{}\overline{3}\overline{3}\overline{3}\overline{3}1\overline{3}\overline{2}\overline{1}23\overline{3}\overline{3}\mathrm{}.$$
According to Corollary 14, the evolution is decomposed into the following steps.
$`p`$ $`=\mathrm{}\overline{3}\overline{3}1\overline{2}3\overline{3}\overline{1}2\overline{3}\overline{3}\overline{3}\overline{3}\mathrm{}`$
$`S_3(p)`$ $`=\mathrm{}331\overline{2}3\overline{3}\overline{1}23333\mathrm{}`$
$`S_2S_3(p)`$ $`=\mathrm{}221\overline{3}2\overline{3}\overline{1}23322\mathrm{}`$
$`S_1S_2S_3(p)`$ $`=\mathrm{}111\overline{3}2\overline{3}\overline{2}13311\mathrm{}`$
$`S_0S_1S_2S_3(p)`$ $`=\mathrm{}\overline{2}\overline{2}\overline{3}\overline{2}\overline{1}\overline{3}\overline{2}133\overline{2}\overline{2}\mathrm{}`$
$`S_2S_0S_1S_2S_3(p)`$ $`=\mathrm{}\overline{3}\overline{3}\overline{3}\overline{3}\overline{1}\overline{3}\overline{2}123\overline{3}\overline{3}\mathrm{}`$
$`\sigma _BS_2S_0S_1S_2S_3(p)`$ $`=\mathrm{}\overline{3}\overline{3}\overline{3}\overline{3}1\overline{3}\overline{2}\overline{1}23\overline{3}\overline{3}\mathrm{}=T(p),`$
where in the last step we used the fact that $`\sigma _B`$ interchanges the letters $`1`$ and $`\overline{1}`$ in each component as specified in Table I. Here the background configuration is changing except the last step.
Alternatively the evolution may also be decomposed according to (37) into the following steps, in which the original background $`\mathrm{}\delta _{l_j}[\overline{3}]\delta _{l_{j+1}}[\overline{3}]\mathrm{}`$ is kept unchanged.
$`p=𝒯_0(p)`$ $`=\mathrm{}\overline{3}\overline{3}1\overline{2}3\overline{3}\overline{1}2\overline{3}\overline{3}\overline{3}\overline{3}\mathrm{}`$
$`𝒯_1(p)`$ $`=\mathrm{}\overline{3}\overline{3}1\overline{2}\overline{3}3\overline{1}2\overline{3}\overline{3}\overline{3}\overline{3}\mathrm{}`$
$`𝒯_2(p)`$ $`=\mathrm{}\overline{3}\overline{3}1\overline{2}\overline{3}\overline{2}\overline{1}\overline{3}22\overline{3}\overline{3}\mathrm{}`$
$`𝒯_3(p)`$ $`=\mathrm{}\overline{3}\overline{3}\overline{3}\overline{2}1\overline{2}\overline{1}\overline{3}22\overline{3}\overline{3}\mathrm{}`$
$`𝒯_4(p)`$ $`=\mathrm{}\overline{3}\overline{3}\overline{3}\overline{2}1\overline{3}\overline{2}\overline{1}22\overline{3}\overline{3}\mathrm{}`$
$`𝒯_5(p)`$ $`=\mathrm{}\overline{3}\overline{3}\overline{3}\overline{3}1\overline{3}\overline{2}\overline{1}23\overline{3}\overline{3}\mathrm{}=T(p).`$
Regarding $`\overline{3}`$ as empty space in a box, one can interpret the above patterns as a motion of particles and anti-particles which can form a neutral (weight 0) bound state. We hope to report on the explicit algorithm for general $`\widehat{𝔤}_n`$ elsewhere.
Acknowledgements The authors thank Masato Okado and Yasuhiko Yamada for useful discussions. One of the authors (A.K.) appreciates the warm hospitality of the organizers of The Baxter Revolution in Mathematical Physics, held at Australian National University, Canberra during February 13–19, 2000, where a part of this work was presented.
## Appendix A Parameterization of $`B_l`$
We list the parameterization of the crystal $`B_l`$. In $`A_n^{(1)}`$ case, it may be identified with the set of semistandard Young tableaux of length $`l`$ one row shape on letters $`\{1,\mathrm{},n+1\}`$. For the other $`\widehat{𝔤}_n`$, $`B_l`$ may be viewed as a similar set with some additional constraints. The relevant letters are $`\{1,\mathrm{},n,\overline{n},\mathrm{},\overline{1}\}`$ as well as $`0`$ and/or $`\mathrm{}`$ depending on $`\widehat{𝔤}_n`$. The crystal structure (actions of $`\stackrel{~}{e}_i,\stackrel{~}{f}_i`$) can be found in the article for $`C_n^{(1)}`$ and the ones for the other $`\widehat{𝔤}_n`$.
$$A_n^{(1)}:B_l=\{(x_1,\mathrm{},x_{n+1})^{n+1}|x_i0,\underset{i=1}{\overset{n+1}{}}x_i=l\}.$$
$$A_{2n1}^{(2)}:B_l=\{(x_1,\mathrm{},x_n,\overline{x}_n,\mathrm{},\overline{x}_1)^{2n}|x_i,\overline{x}_i0,\underset{i=1}{\overset{n}{}}(x_i+\overline{x}_i)=l\}.$$
$`A_{2n}^{(2)}:`$ $`B_l=\{(x_1,\mathrm{},x_n,\overline{x}_n,\mathrm{},\overline{x}_1)^{2n}|x_i,\overline{x}_i0,{\displaystyle \underset{i=1}{\overset{n}{}}}(x_i+\overline{x}_i)l\}.`$
$`\text{We set }x_0=l{\displaystyle \underset{i=1}{\overset{n}{}}}(x_i+\overline{x}_i).`$
$$B_n^{(1)}:B_l=\{(x_1,\mathrm{},x_n,x_0,\overline{x}_n,\mathrm{},\overline{x}_1)^{2n}\times \{0,1\}|\begin{array}{c}x_0=0\text{or}\mathrm{\hspace{0.33em}1},x_i,\overline{x}_i0,\\ x_0+_{i=1}^n(x_i+\overline{x}_i)=l\end{array}\}.$$
$`C_n^{(1)}:`$ $`B_l=\{(x_1,\mathrm{},x_n,\overline{x}_n,\mathrm{},\overline{x}_1)^{2n}|\begin{array}{c}x_i,\overline{x}_i0,\hfill \\ l_{i=1}^n(x_i+\overline{x}_i)l2\hfill \end{array}\}.`$
$`\text{We set }x_0=(l{\displaystyle \underset{i=1}{\overset{n}{}}}(x_i+\overline{x}_i))/2.`$
$$D_n^{(1)}:B_l=\{(x_1,\mathrm{},x_n,\overline{x}_n,\mathrm{},\overline{x}_1)^{2n}|\begin{array}{c}x_n=0\text{or}\overline{x}_n=0,x_i,\overline{x}_i0,\\ _{i=1}^n(x_i+\overline{x}_i)=l\end{array}\}.$$
$`D_{n+1}^{(2)}:`$ $`B_l=\{(x_1,\mathrm{},x_n,x_0,\overline{x}_n,\mathrm{},\overline{x}_1)^{2n}\times \{0,1\}|\begin{array}{c}x_0=0\text{ or }1,x_i,\overline{x}_i0,\hfill \\ x_0+_{i=1}^n(x_i+\overline{x}_i)l\hfill \end{array}\}.`$
$`\text{We set }x_{\mathrm{}}=lx_0{\displaystyle \underset{i=1}{\overset{n}{}}}(x_i+\overline{x}_i).`$ |
warning/0003/hep-th0003150.html | ar5iv | text | # Nonperturbative approach to a simple model with ultravioletly divergent eigenenergies in perturbation theory
## I Introduction
Models in quantum field theories usually have the problem of ultraviolet divergence in the framework of perturbation theory (see, e.g., ). For renormalizable models, ultraviolet divergence can be removed by renormalization techniques, but at the cost that renormalized masses of particles can not be explained in the framework of the theory. On the other hand, non-renormalizable models are usually rejected, since for them it is not clear at the present stage how to obtain finite results for high order corrections in perturbation theory. However, rigorously to say, the break-down of perturbative treatment to a model does not necessarily mean that the model cannot give meaningful results when it is treated nonperturbatively. The point here is that nonperturbative approach to behaviors of states of ultravioletly divergent models is generally quite difficult. In particular, it is still not clear whether the Hamiltonian of a model suffering ultraviolet divergence in the framework of perturbation theory could have eigenstates with finite energies. This problem is of interest, since, if the answer is positive, then, the ultraviolet divergence may be avoidable in nonperturbative treatment and finite energies of ground states may be associated with masses of particles observed experimentally.
In this paper we will show that the answer is indeed positive. For this purpose, we study a simple quantized model, which gives ultravioletly divergent results in the framework of perturbation theory. The model has an interaction structure similar to that of QED in the number representation, so that most of the arguments given to it can be extended to the case of QED. (In this paper, by interaction structure in a representation, we mean the way in which basis states of the representation are coupled by the interaction term of the Hamiltonian, that is, the structure of the non-zero off-diagonal elements of the Hamiltonian matrix in the representation.) For this reason, whether the simple model has an appropriate form in configuration space does not matter here.
Similar to QED, the Hilbert space of the simple model studied in this paper is infinite even when the momentum space is cut off and discretized. For the infinite Hilbert space, not all the theorems for eigen-solutions of Hamiltonians in finite Hilbert spaces hold. For example, eigenstates of the Hamiltonian of the simple model do not span the whole infinite Hilbert space. In particular, we find that eigenstates and eigenenergies of the Hamiltonian can be expressed in terms of themselves, and, as a result, the eigenenergies can remain finite when the cut-off of momentum is taken off, even if second order corrections to the eigenenergies in perturbation theory is ultravioletly divergent.
Concretely, this paper is organized in the following way. In section II, we introduce the quantized model, discuss the structure of basis states and the structure of the Hamiltonian matrix in the basis states. In section III, we truncate the infinite Hilbert space of the model, so that obtain a series of finite Hilbert spaces, the limit of which gives the infinite Hilbert space. For each truncated finite Hilbert space, we construct another set of basis states by making use of energy eigenstates of another finite Hilbert space. The Hamiltonian matrix in the representation of the new set of basis states is quite simple and can be diagonalized easily. Section IV is devoted to discussions for eigenstates and eigenenergies of the Hamiltonian of the model in the infinite Hilbert space. We show that it is possible for the eigenenergies to be finite, even if perturbative treatment gives divergent results. Conclusions and discussions are given in section V.
## II A simple quantized model
Since realistic models, such as the standard model, are complicated, in this paper, as a first step to the method that is to be developed for nonperturbative approach to eigenstates and eigenenergies of Hamiltonians in quantum field theories, we choose a model as simple as possible to study. Such a model should have an interaction structure similar to that of QED and may have ultraviolet divergence, so that the method developed in this paper for the model can be extended to treat realistic models, such as QED and the standard model. Since here we are interested in properties of eigen-solutions of Hamiltonians only, we are not to start from a classical Lagrangian expressed in configuration space, but start from a Hamiltonian expressed in terms of creation and annihilation operators for free fields.
The simplest model satisfying the above requirements is composed of a quantized fermion field and a quantized boson field in 1-dimensional momentum space. Denoting the creation and annihilation operators for a free fermion field with momentum $`p`$ and for a free boson field with momentum $`k`$ as $`b^{}(p)`$, $`b(p)`$, and $`a^{}(k)`$, a(k), respectively, the Hamiltonian of the model is taken as
$$H=H_f+H_b+H_I$$
(1)
where
$`H_f={\displaystyle p_0b^{}(p)b(p)𝑑p}`$ (2)
$`H_b={\displaystyle k_0a^{}(k)a(k)𝑑k}`$ (3)
$`H_I={\displaystyle }(V(p_1,p_2,k)b^{}(p_2)b(p_1)a^{}(k)+h.c.)\delta (p_1p_2k)dp_1dp_2dk`$ (4)
with $`p_0=|p|`$ and $`k_0=|k|`$. The operators $`b^{}(p)`$ and $`b(p^{})`$ satisfy the usual anticommutation relations and $`a^{}(k)`$ and $`a(k^{})`$ satisfy the usual commutation relations. Here $`\mathrm{}`$ and $`c`$ are taken to be unit, $`\mathrm{}=c=1`$. The interaction structure of this Hamiltonian in the number representation, given by the expression of $`H_I`$ in Eq. (4), is clearly similar to (although simpler than) that of QED.
For the sake of convenience in discussing properties of the Hamiltonian $`H`$, we discretize the momentum space and take a cut-off $`\mathrm{\Lambda }`$, that is, we take
$$p=p_i=i\mathrm{\Delta }p\mathrm{\Lambda },k=k_j=j\mathrm{\Delta }p\mathrm{\Lambda },$$
(5)
where $`i,j=0,1,\mathrm{},N`$ with $`N=2\mathrm{\Lambda }/\mathrm{\Delta }p`$. The Hamiltonian of the model expressed in terms of summations over $`p=p_i`$ and $`k=k_j`$ is
$$H(\mathrm{\Lambda })=H_f(\mathrm{\Lambda })+H_b(\mathrm{\Lambda })+H_I(\mathrm{\Lambda }),$$
(6)
where
$`H_f(\mathrm{\Lambda })={\displaystyle \underset{p}{}}p_0b^{}(p)b(p)`$ (7)
$`H_b(\mathrm{\Lambda })={\displaystyle \underset{k}{}}k_0a^{}(k)a(k)`$ (8)
$`H_I(\mathrm{\Lambda })={\displaystyle \underset{p_1,p_2}{}}(V(p_1,p_2)b^{}(p_2)b(p_1)a^{}(p_1p_2)+h.c.).`$ (9)
In the limit of $`\mathrm{\Delta }p0`$ and $`\mathrm{\Lambda }\mathrm{}`$, $`H(\mathrm{\Lambda })`$ becomes the $`H`$ in Eq. (1). In this paper we assume that with finite cut-off $`\mathrm{\Lambda }`$ the model is free from divergence.
Since $`H_I(\mathrm{\Lambda })a^{}(k)|0=0`$, where $`|0`$ is the vacuum state, states $`a^{}(k)|0`$, denoted by $`|s_k`$ after normalization, are eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$. We are not interested in this kind of trivial eigenstates here. What we are interested in are fermion-type eigenstates. The simplest fermion-type eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$ are in the Hilbert space spanned by the state $`b^{}(p)|0`$, denoted by $`|f_p`$ after normalization, and all the states that can be coupled to $`|f_p`$ by $`H_I^m(\mathrm{\Lambda })`$, the product of $`m`$ $`H_I(\mathrm{\Lambda })`$, with $`m=0,1,2,\mathrm{}`$. It is this Hilbert space that we are to study in this paper, which will be denoted by $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$.
Concretely to say, basis states of the Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ can be taken as
$$|f_{pk_1\mathrm{}k_m}s_{k_1}\mathrm{}s_{k_m}N_{p,k_1,\mathrm{},k_m}b^{}(pk_1\mathrm{}k_m)a^{}(k_1)\mathrm{}a^{}(k_m)|0$$
(10)
for $`m=0,1,2,\mathrm{}`$ (the case for $`m=0`$ is just $`|f_p`$), where $`N_{p,k_1,\mathrm{},k_m}`$ are normalization coefficients. Noticing that the basis states $`|f_{pk_1\mathrm{}k_m}s_{k_1}\mathrm{}s_{k_m}`$ for all possible $`k_1,\mathrm{},k_m`$ can be coupled to the state $`|f_p`$ by $`H_I^m(\mathrm{\Lambda })`$, we will denote the set of them by $`\{H_I^m(\mathrm{\Lambda })|f_p\}`$ in what follows. Then, the basis states of the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ in Eq. (10) are elements of the following sets
$$\{H_I^0(\mathrm{\Lambda })|f_p=|f_p\},\{H_I(\mathrm{\Lambda })|f_p\},\{H_I^2(\mathrm{\Lambda })|f_p\},\{H_I^3(\mathrm{\Lambda })|f_p\},\mathrm{}.$$
(11)
In some cases in the following sections, for brevity, instead of the expression in Eq. (10), we use $`|\xi _{i_m}(m,p,\mathrm{\Lambda })`$ to denote basis states in the set $`\{H_I^m(\mathrm{\Lambda })|f_p\}`$, i.e., we use $`i_m`$ to denote $`(k_1,\mathrm{},k_m)`$. For example, $`|\xi _{i_0}(0,p,\mathrm{\Lambda })`$ indicates $`|f_p`$ with $`i_0=1`$, $`|\xi _{i_1}(1,p,\mathrm{\Lambda })`$ indicates $`|f_{pk_1}s_{k_1}`$ with $`i_1=k_1`$, and so on.
The interaction structure of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the basis states in Eq. (10), i.e., the structure of the non-zero off-diagonal elements of the Hamiltonian matrix in the basis states, has an interesting tree structure: The basis state $`|f_p`$ is coupled to basis states in the set $`\{H_I(\mathrm{\Lambda })|f_p\}`$ only; basis states in the set $`\{H_I(\mathrm{\Lambda })|f_p\}`$ are coupled to basis states in the sets $`\{|f_p\}`$ and $`\{H_I^2(\mathrm{\Lambda })|f_p\}`$ only; $`\mathrm{}`$; basis states in the set $`\{H_I^m(\mathrm{\Lambda })|f_p\}`$ are coupled to basis states in the sets $`\{H_I^{m1}(\mathrm{\Lambda })|f_p\}`$ and $`\{H_I^{m+1}(\mathrm{\Lambda })|f_p\}`$ only; $`\mathrm{}`$. It is easy to show that, with this structure of the Hamiltonian matrix, when the coupling strength $`V(p_1,p_2)`$ in Eq. (9) is strong enough, second order corrections to the eigenenergies of the Hamiltonian $`H(\mathrm{\Lambda })`$ are ultravioletly divergent in perturbation theory when the cut-off $`\mathrm{\Lambda }`$ approaches to infinity.
## III Truncated Hilbert space and $`\psi _s`$ representation
The Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ spanned by the basis states in the sets given in (11) is infinite, although the momentum space has been cut off by $`\mathrm{\Lambda }`$. Since the problem of eigen-solutions of a Hamiltonian in an infinite Hilbert space is more difficult than that in a finite Hilbert space and many theorems in the latter case are invalid in the former case, in this section we truncate the Hilbert space to a series of finite ones and discuss properties of the eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the truncated finite Hilbert spaces. Then, in the next section, we discuss what we could have when the truncated finite Hilbert spaces resume the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$.
### A Truncated finite Hilbert spaces
A truncated finite Hilbert space, denoted by $`L_n(p,\mathrm{\Lambda })`$, is spanned by basis states in a set $`A(n,p,\mathrm{\Lambda })`$ defined by
$$A(n,p,\mathrm{\Lambda })=\{|f_p\}\{H_I(\mathrm{\Lambda })|f_p\}\mathrm{}\{H_I^n(\mathrm{\Lambda })|f_p\}.$$
(12)
When $`n`$ goes to infinity, the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$ will become the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$.
Normalized eigenstates and the corresponding eigenenergies of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the truncated finite Hilbert space $`L_n(p,\mathrm{\Lambda })`$ are denoted by $`|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })`$ and $`E_{\alpha _n}(n,p,\mathrm{\Lambda })`$, respectively,
$$H(\mathrm{\Lambda })|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })=E_{\alpha _n}(n,p,\mathrm{\Lambda })|\psi _{\alpha _n}(n,p,\mathrm{\Lambda }).$$
(13)
These states $`|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })`$ also span the Hilbert space $`L_n(p,\mathrm{\Lambda })`$ and can be expanded in the basis states in the set $`A(n,p,\mathrm{\Lambda })`$,
$$|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })=\underset{m=0}{\overset{n}{}}\underset{i_m}{}C_{\alpha _n,i_m}(n,m,p,\mathrm{\Lambda })|\xi _{i_m}(m,p,\mathrm{\Lambda }),$$
(14)
where $`C_{\alpha _n,i_m}(n,m,p,\mathrm{\Lambda })`$ are expanding coefficients.
The set $`A(n,p,\mathrm{\Lambda })`$ has an interesting structure: It can be expressed by making use of the sets $`A(n1,pk,\mathrm{\Lambda })`$. To see this, first note that the product of two basis states $`|s_k`$ and $`|f_{pkk_1\mathrm{}k_m}s_{k_1}\mathrm{}s_{k_m}`$ is
$$|s_k|f_{pkk_1\mathrm{}k_m}s_{k_1}\mathrm{}s_{k_m}=|f_{pkk_1\mathrm{}k_m}s_{k_1}\mathrm{}s_{k_m}s_k.$$
(15)
From Eq. (10), it is easy to see that each basis state in the set $`\{H_I^m(\mathrm{\Lambda })|f_p\}`$ with $`m1`$ can be expressed as the product of a basis state $`|s_k`$ and a basis state in the set $`\{H_I^{m1}(\mathrm{\Lambda })|f_{pk}\}`$. Then, denoting the set of the product of the basis states $`|s_k`$ and the basis states in the set $`A(n1,pk,\mathrm{\Lambda })`$ for all possible $`k`$ as $`B(n1,p,k,\mathrm{\Lambda })`$,
$$B(n1,p,k,\mathrm{\Lambda })=\{|s_k|\xi :\mathrm{for}|\xi A(n1,pk,\mathrm{\Lambda })\mathrm{and}\mathrm{all}\mathrm{possible}k\},$$
(16)
the set $`A(n,p,\mathrm{\Lambda })`$ can be reexpressed as
$$A(n,p,\mathrm{\Lambda })=\{|f_p\}\underset{k}{}B(n1,p,k,\mathrm{\Lambda }).$$
(17)
In the next subsection, we show that this structure of the set $`A(n,p,\mathrm{\Lambda })`$ enables us to construct another set of basis states, in which the Hamiltonian matrix can be diagonalized explicitly.
### B $`\psi _s`$-representation in truncated finite Hilbert spaces
Making use of the definition of product of basis states given in Eq. (15) and the expansion of eigenstates in Eq. (14), it is easy to get $`|s_k|\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$, the product of a state $`|s_k`$ and an eigenstate $`|\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ in the truncated Hilbert space $`L_{n1}(pk,\mathrm{\Lambda })`$, which will be denoted by $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$,
$$|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })=\underset{m=0}{\overset{n1}{}}\underset{i_m}{}C_{\alpha _{n1},i_m}(n1,m,pk,\mathrm{\Lambda })|s_k|\xi _{i_m}(m,pk,\mathrm{\Lambda }).$$
(18)
These states $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ span the same Hilbert space as the states in the set $`B(n1,p,k,\mathrm{\Lambda })`$ in Eq. (16). Therefore, these states together with the state $`|f_p`$ also span the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$.
Now we show that in the limit of $`\mathrm{\Delta }p0`$, the states $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ are orthogonal to each other and together with $`|f_p`$ form another set of orthogonal basis states for the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$. To prove this, we rewrite $`|\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ in the form
$$|\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })=\underset{m=0}{\overset{n1}{}}|h_{\alpha _{n1}}(n1,m,pk,\mathrm{\Lambda }),$$
(19)
where
$$|h_{\alpha _{n1}}(n1,m,pk,\mathrm{\Lambda })=\underset{i_m}{}C_{\alpha _{n1},i_m}(n1,m,pk,\mathrm{\Lambda })|\xi _{i_m}(m,pk,\mathrm{\Lambda }),$$
(20)
for example,
$`|h_{\alpha _{n1}}(n1,0,pk,\mathrm{\Lambda })=C_{\alpha _{n1},1}(n1,0,pk,\mathrm{\Lambda })|f_{pk}`$ (21)
$`|h_{\alpha _{n1}}(n1,1,pk,\mathrm{\Lambda })={\displaystyle \underset{k_1}{}}C_{\alpha _{n1},k_1}(n1,1,pk,\mathrm{\Lambda })|f_{pkk_1}s_{k_1}.`$ (22)
The states $`|h_{\alpha _{n1}}(n1,m,pk,\mathrm{\Lambda })`$ are projections of the state $`|\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ in the subspaces spanned by states in the sets $`\{H_I^m(\mathrm{\Lambda })|f_{pk}\}`$, respectively, therefore, $`|h_{\alpha _{n1}}(n1,m,pk,\mathrm{\Lambda })`$ with different $`m`$ are orthogonal to each other.
Making use of Eq. (19), we have
$$s_k^{}\psi _{\alpha _{n1}^{}}(n1,pk^{},\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })=\underset{m=0}{\overset{n1}{}}I_m,$$
(23)
where
$$I_m=s_k^{}h_{\alpha _{n1}^{}}(n1,m,pk^{},\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })$$
(24)
with $`|s_k^{}h_{\alpha _{n1}^{}}(n1,m,pk^{},\mathrm{\Lambda })=|s_k^{}|h_{\alpha _{n1}^{}}(n1,m,pk^{},\mathrm{\Lambda })`$ defined in the same way as $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ in Eq. (18). It is easy to verify that
$`I_0=C_{\alpha _{n1}^{},1}^{}(n1,0,pk^{},\mathrm{\Lambda })C_{\alpha _{n1},1}(n1,0,pk,\mathrm{\Lambda })\delta _{kk^{}}`$ (25)
$`I_1=C_{\alpha _{n1}^{},k}^{}(n1,1,pk^{},\mathrm{\Lambda })C_{\alpha _{n1},k^{}}(n1,1,pk,\mathrm{\Lambda })`$ (26)
$`+{\displaystyle \underset{k_1}{}}C_{\alpha _{n1}^{},k_1}^{}(n1,1,pk^{},\mathrm{\Lambda })C_{\alpha _{n1},k_1}(n1,1,pk,\mathrm{\Lambda })\delta _{kk^{}}.`$ (27)
Here we assume that in the limit of $`\mathrm{\Delta }p0`$, there is no singularity in the coefficients $`C_{\alpha _{n1},i_m}(n1,m,pk,\mathrm{\Lambda })`$ for fixed $`n1`$ and $`m`$, i.e., $`C_{\alpha _{n1},i_m}(n1,m,pk,\mathrm{\Lambda })`$ is a smooth function of $`i_m`$. Then, since
$$\underset{k_1}{}|C_{\alpha _{n1},k_1}(n1,1,pk,\mathrm{\Lambda })|^2<1,$$
(28)
when the interval $`\mathrm{\Delta }p`$ for discretizing the momentum $`k_1`$ goes to zero, the value of $`C_{\alpha _{n1}^{},k}^{}(n1,1,pk^{},\mathrm{\Lambda })C_{\alpha _{n1},k^{}}(n1,1,pk,\mathrm{\Lambda })`$ should approach to zero, and Eq. (27) gives
$$\underset{\mathrm{\Delta }p0}{lim}I_1=\underset{k_1}{}C_{\alpha _{n1}^{},k_1}^{}(n1,1,pk^{},\mathrm{\Lambda })C_{\alpha _{n1},k_1}(n1,1,pk,\mathrm{\Lambda })\delta _{kk^{}}.$$
(29)
Similar results can also be obtained for the other $`I_m`$, and finally we have
$`\underset{\mathrm{\Delta }p0}{lim}s_k^{}\psi _{\alpha _{n1}^{}}(n1,pk^{},\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ (30)
$`={\displaystyle \underset{m=0}{\overset{n1}{}}}{\displaystyle \underset{i_m}{}}C_{\alpha _{n1}^{},i_m}^{}(n1,m,pk^{},\mathrm{\Lambda })C_{\alpha _{n1},i_m}(n1,m,pk,\mathrm{\Lambda })\delta _{kk^{}}=\delta _{\alpha _{n1}\alpha _{n1}^{}}\delta _{kk^{}}.`$ (31)
That is to say, the states $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ are normalized and orthogonal to each other. Therefore, these states and the state $`|f_p`$ form another set of normalized orthogonal basis states for the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$. The representation given by this set of basis states in the limit of $`\mathrm{\Delta }p0`$, will be termed $`\psi _s`$-representation in what follows.
The Hamiltonian $`H(\mathrm{\Lambda })`$ has a quite simple matrix form in the $`\psi _s`$-representation of the Hilbert space $`L_n(p,\mathrm{\Lambda })`$. In fact, similar to Eq. (31), one can verify that
$`\underset{\mathrm{\Delta }p0}{lim}s_k^{}\psi _{\alpha _{n1}^{}}(n1,pk^{},\mathrm{\Lambda })|H(\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ (32)
$`=s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })|H(\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })\delta _{\alpha _{n1}\alpha _{n1}^{}}\delta _{kk^{}}.`$ (33)
Therefore, the non-zero off-diagonal elements of the Hamiltonian matrix in the $`\psi _s`$-representation are those connecting the state $`|f_p`$ and the states $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ only. Diagonalization of the Hamiltonian matrix with such a simple structure is quite easy, which gives
$`E_{\alpha _n}(n,p,\mathrm{\Lambda })f_p|H(\mathrm{\Lambda })|f_p`$ (34)
$`={\displaystyle \underset{k,\alpha _{n1}}{}}{\displaystyle \frac{\left|f_p|H(\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })\right|^2}{E_{\alpha _n}(n,p,\mathrm{\Lambda })s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })|H(\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })}}`$ (35)
and
$$|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })=D_{\alpha _n}(n,p,\mathrm{\Lambda })|f_p+\underset{k,\alpha _{n1}}{}D_{\alpha _n,k\alpha _{n1}}(n,p,\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda }),$$
(36)
where
$$D_{\alpha _n,k\alpha _{n1}}(n,p,\mathrm{\Lambda })=\frac{s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })|H(\mathrm{\Lambda })|f_pD_{\alpha _n}(n,p,\mathrm{\Lambda })}{E_{\alpha _n}(n,p,\mathrm{\Lambda })s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })|H(\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })}$$
(37)
and
$`D_{\alpha _n}(n,p,\mathrm{\Lambda })`$ (38)
$`=\left[1+{\displaystyle \underset{k,\alpha _{n1}}{}}{\displaystyle \frac{|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })|H(\mathrm{\Lambda })|f_p|^2}{\left(E_{\alpha _n}(n,p,\mathrm{\Lambda })s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })|H(\mathrm{\Lambda })|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })\right)^2}}\right]^{\frac{1}{2}}.`$ (39)
## IV Eigenenergies and eigenstates in infinite Hilbert space
When $`n`$ goes to infinity, the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$ discussed in the previous section will become the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$. However, eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the infinite Hilbert space cannot be obtained by simply letting the index $`n`$ go to infinity in the eigenstates $`|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })`$ of the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$. In fact, when $`n`$ becomes $`n+1`$, the number of eigenstates in the truncated Hilbert space will become $`N`$ times larger ($`N`$ is the number of discretized momenta); and despite of how large $`n`$ is, there exist eigenstates $`|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })`$ for which the values of $`|h_{\alpha _n}(n,n,p,\mathrm{\Lambda })|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })|^2`$ are not small, i.e., there exist eigenstates whose projections in the subspace spanned by states in the set $`\{H_I^n(\mathrm{\Lambda })|f_p\}`$ are not small. Therefore, the definition for eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ should be treated more carefully. In subsection IV A, we give the definition and discuss some properties of the eigenstates defined by it. In subsection IV B, we show that eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ can be expressed in terms of themselves, and the corresponding eigenenergies can remain finite when the cut-off $`\mathrm{\Lambda }`$ is taken off, even if they are ultravioletly divergent in perturbation theory. Subsection IV C is devoted to a brief discussion for evolution of states with time in the infinite Hilbert space.
### A Energy eigenstates in infinite Hilbert space
A normalized eigenstate of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$, denoted by $`|\varphi _\beta (p,\mathrm{\Lambda })`$, with eigenenergy $`E_\beta (p,\mathrm{\Lambda })`$ is defined in the following way: For each positive number $`ϵ`$, there exists a number $`N(ϵ)`$ such that for each $`n`$ not smaller than $`N(ϵ)`$, there exists an eigenstate $`|\psi _\beta (n,p,\mathrm{\Lambda })`$ of the Hamiltonian $`H(\mathrm{\Lambda })`$ with eigenenergy $`E_\beta (n,p,\mathrm{\Lambda })`$ in the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$ satisfying
$$1|\psi _\beta (n,p,\mathrm{\Lambda })|\varphi _\beta (p,\mathrm{\Lambda })|^2<ϵ$$
(40)
and
$$|E_\beta (p,\mathrm{\Lambda })E_\beta (n,p,\mathrm{\Lambda })|<ϵ.$$
(41)
Then, each eigenstate $`|\varphi _\beta (p,\mathrm{\Lambda })`$ is the limit of a series of states in the truncated Hilbert spaces $`L_n(p,\mathrm{\Lambda })`$,
$$|\varphi _\beta (p,\mathrm{\Lambda })=\underset{n\mathrm{}}{lim}|\psi _\beta (n,p,\mathrm{\Lambda }).$$
(42)
A property of an eigenstate $`|\varphi _\beta (p,\mathrm{\Lambda })`$ is that its projection is less than $`ϵ`$ in the subspace of the infinite Hilbert space spanned by states in the sets $`\{H_I^{N(ϵ)+1}\}`$, $`\{H_I^{N(ϵ)+2}\}`$, $`\mathrm{}`$. In fact, substituting Eq. (40) into the normalization condition
$$\underset{\alpha _{N(ϵ)}}{}|\psi _{\alpha _{N(ϵ)}}(N(ϵ),p,\mathrm{\Lambda })|\varphi _\beta (p,\mathrm{\Lambda })|^2+\underset{m=N(ϵ)+1}{\overset{\mathrm{}}{}}\underset{i_m}{}|\xi _{i_m}(m,p,\mathrm{\Lambda })|\varphi _\beta (p,\mathrm{\Lambda })|^2=1,$$
(43)
we have
$$\underset{m=N(ϵ)+1}{\overset{\mathrm{}}{}}\underset{i_m}{}|\xi _{i_m}(m,p,\mathrm{\Lambda })|\varphi _\beta (p,\mathrm{\Lambda })|^2<ϵ.$$
(44)
Therefore, the projection of each of the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$ in the subspace spanned by states in the set $`\{H_I^n(\mathrm{\Lambda })|f_p\}`$ approaches to zero when $`n\mathrm{}`$. Due to this property of the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$, one can see that for large $`n`$ and small $`ϵ`$, the number of the states $`|\psi _\beta (n,p,\mathrm{\Lambda })`$, which are close to the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$, must be much smaller than the total number of the states $`|\psi _{\alpha _n}(n,p,\mathrm{\Lambda })`$ in the truncated Hilbert space $`L_n(p,\mathrm{\Lambda })`$. As a result, the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$ span a small part of the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ only.
### B Expressions of eigenenergies and eigenstates in $`\varphi _s`$-subspace
As discussed above, when $`n\mathrm{}`$, not all the states $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$, but a fraction of them, namely, $`|s_k\psi _\beta (n1,pk,\mathrm{\Lambda })`$ have definite limit
$$\underset{n\mathrm{}}{lim}|s_k\psi _\beta (n1,pk,\mathrm{\Lambda })=|s_k\varphi _\beta (pk,\mathrm{\Lambda }).$$
(45)
The subspace of the Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$ spanned by the states
$$|f_p\mathrm{and}|s_k\varphi _\beta (pk,\mathrm{\Lambda })$$
(46)
for all possible $`k`$ and $`\beta `$ will be termed $`\varphi _s`$-subspace. Since all the states $`|s_k\psi _{\alpha _{n1}}(n1,pk,\mathrm{\Lambda })`$ with $`\alpha _{n1}\beta `$ do not have definite limit when $`n`$ goes to $`\mathrm{}`$, the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$ must lie in the $`\varphi _s`$-subspace of the infinite Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$. Therefore, if one diagonalizes the Hamiltonian $`H(\mathrm{\Lambda })`$ in the $`\varphi _s`$-subspace, the eigenstates obtained in the $`\varphi _s`$-subspace must contain all the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$ of the whole Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$. As a result, in order to obtain the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$, one does not need to diagonalize the Hamiltonian $`H(\mathrm{\Lambda })`$ in the whole Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$, but diagonalization in the $`\varphi _s`$-subspace is enough.
Similar to Eq. (31), one can show that the states in (46) are normalized and orthogonal to each other. Furthermore, the Hamiltonian matrix in the $`\varphi _s`$-subspace has a structure similar to that in the $`\psi _s`$-representation of the Hilbert space $`L_n(p,\mathrm{\Lambda })`$, i.e., there is coupling between the state $`|f_p`$ and each of the states $`|s_k\varphi _\beta (pk,\mathrm{\Lambda })`$, but the coupling between each two of the states $`|s_k\varphi _\beta (pk,\mathrm{\Lambda })`$ approaches to zero when $`\mathrm{\Delta }p0`$. Then, similar to Eqs. (35) and (36), in the limit of $`\mathrm{\Delta }p0`$, one can obtain expressions of the eigenstates and eigenenergies of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the $`\varphi _s`$-subspace. In particular, for eigenenergies and eigenstates of the Hamiltonian $`H(\mathrm{\Lambda })`$ in the whole Hilbert space $`L_{\mathrm{}}(p,\mathrm{\Lambda })`$, after simplification we have
$`E_\beta (p,\mathrm{\Lambda })=p_0+{\displaystyle \underset{k,\beta ^{}}{}^{}}{\displaystyle \frac{\left|f_p|H(\mathrm{\Lambda })|f_{pk}s_k\right|^2|d_\beta ^{}(pk,\mathrm{\Lambda })|^2}{E_\beta (p,\mathrm{\Lambda })k_0E_\beta ^{}(pk,\mathrm{\Lambda })}}`$ (47)
$`|\varphi _\beta (p,\mathrm{\Lambda })=D_\beta (p,\mathrm{\Lambda })|f_p+{\displaystyle \underset{k,\beta ^{}}{}^{}}D_{\beta ,k\beta ^{}}(p,\mathrm{\Lambda })|s_k\varphi _\beta ^{}(pk,\mathrm{\Lambda }),`$ (48)
where
$`D_{\beta ,k\beta ^{}}(p,\mathrm{\Lambda })={\displaystyle \frac{f_{pk}s_k|H(\mathrm{\Lambda })|f_pd_\beta ^{}(pk,\mathrm{\Lambda })}{E_\beta (p,\mathrm{\Lambda })k_0E_\beta ^{}(pk,\mathrm{\Lambda })}}D_\beta (p,\mathrm{\Lambda })`$ (49)
$`D_\beta (p,\mathrm{\Lambda })=\left[1+{\displaystyle \underset{k,\beta ^{}}{}^{}}{\displaystyle \frac{\left|f_p|H(\mathrm{\Lambda })|f_{pk}s_k\right|^2|d_\beta ^{}(pk,\mathrm{\Lambda })|^2}{\left(E_\beta (p,\mathrm{\Lambda })k_0E_\beta ^{}(pk,\mathrm{\Lambda })\right)^2}}\right]^{\frac{1}{2}}`$ (50)
$$d_\beta ^{}(pk,\mathrm{\Lambda })=\varphi _\beta ^{}(pk,\mathrm{\Lambda })|f_{pk}$$
(51)
and the primes over the summations mean that when $`k_0=0`$ the index $`\beta ^{}`$ is not equal to $`\beta `$. Note that the right hand sides of Eqs. (47) and (48) contain the eigenenergies and eigenstates themselves.
Now let us compare the expression of the eigenenergy $`E_\beta (p,\mathrm{\Lambda })`$ on the right hand side of Eq. (47) and the second order correction for it in perturbation theory, which is
$$E_\beta ^{(2)}(p,\mathrm{\Lambda })=\underset{k}{}^{}\frac{|f_p|H_I(\mathrm{\Lambda })|f_{pk}s_k|^2}{p_0|pk|k_0}.$$
(52)
The main difference between the right hand side of Eq. (47) and the right hand side of Eq. (52) is that there is a term $`|d_\beta ^{}(pk,\mathrm{\Lambda })|^2`$ in the numerator of the right hand side of Eq. (47). Equation (51) shows that this term is less than one. Let us consider the case that the summation on the right hand side of Eq. (52) goes to infinity as the cut-off $`\mathrm{\Lambda }\mathrm{}`$. For the summation on the right hand side of Eq. (47), if the value of $`|d_\beta ^{}(pk,\mathrm{\Lambda })|^2`$ decreases fast enough with increasing $`k_0`$, then, it would be possible for it to be convergent in the limit of $`\mathrm{\Lambda }\mathrm{}`$. That is to say, for a model suffering ultraviolet divergence in perturbation theory, when it is treated nonperturbatively, it is possible for some of, even all of, its eigenstates to have finite eigenenergies.
At last, we would like to mention that the theory of relativity may give additional restrictions to possible physical eigenstates with finite eigenenergies in the limit of $`\mathrm{\Lambda }\mathrm{}`$. For example, the requirement that $`E_\beta (p,\mathrm{})`$ is finite does not guarantee that it satisfies the relation
$$E_\beta ^2(p,\mathrm{})=E_\beta ^2(0,\mathrm{})+p^2.$$
(53)
Furthermore, the theory of relativity requires that Lorentzian transformation can transform an eigenstate $`|\varphi _\beta (p,\mathrm{})`$ to an eigenstate $`|\varphi _\beta (p^{},\mathrm{})`$. But, it is not clear if all the eigenstates given in Eq. (48) satisfy this requirement.
### C Evolution of states in infinite Hilbert space
Finally, let us give a brief discussion for evolution of states in the infinite Hilbert space $`L_{\mathrm{}}(p,\lambda )`$. In the infinite Hilbert space, the subspace spanned by the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$ of the Hamiltonian $`H(\mathrm{\Lambda })`$, which will be called eigen-subspace , play a special role. States in the eigen-subspace can be expanded in the eigenstates $`|\varphi _\beta (p,\mathrm{\Lambda })`$, therefore, they evolve in the same way as those in a finite Hilbert space, say, for an initial state $`|\mathrm{\Phi }(t=0)=|\mathrm{\Phi }_0`$,
$$|\mathrm{\Phi }(t)=\underset{p,\beta }{}\varphi _\beta (p,\mathrm{\Lambda })|\mathrm{\Phi }_0e^{iE_\beta (p,\mathrm{\Lambda })t}|\varphi _\beta (p,\mathrm{\Lambda }).$$
(54)
In the general case, evolution of states is more complicated than that in Eq. (54). For example, consider a state $`|\mathrm{\Psi }_0`$ lie in a finite subspace $`L_n(p,\lambda )`$ of the infinite Hilbert space $`L_{\mathrm{}}(p,\lambda )`$. It can be divided into two parts: one in the eigen-subspace, denoted by $`|\mathrm{\Psi }_0^{es}`$, the other out side of the eigen-subspace, denoted by $`|\mathrm{\Psi }_0^{nes}`$. Then, the $`|\mathrm{\Psi }_0^{es}`$ part of the state will evolve as in Eq. (54). But, the evolution of $`|\mathrm{\Psi }_0^{nes}`$ is not so clear. As time increases, it is possible for it to spread to subspaces of the Hilbert space spanned by states in the sets $`\{H_I^m(\mathrm{\Lambda })|f_p\}`$ for whatever large $`m`$. In this case, when we are concerned with properties of the state in a finite subspace $`L_n(p,\mathrm{\Lambda })`$ only, the probability of the state $`|\mathrm{\Psi }^{nes}(t)`$ in the subspace $`L_n(p,\mathrm{\Lambda })`$ will become smaller and smaller as $`t`$ increases, which remind one of properties of dissipative systems. At the present stage, it is not clear if this feature of evolution of states is common to all the states out side of the eigen-subspace.
## V Conclusions and discussions
In this paper, we study a simple quantized model, which has an interaction structure similar to (but simpler than ) that of QED and gives ultravioletly divergent results in the framework of perturbation theory. We show that when the eigen-problem of the Hamiltonian of the model is treated nonperturbatively, it is in fact possible for eigenenergies of the Hamiltonian to be finite. The eigenstates of the Hamiltonian are found to span part of the whole infinite Hilbert space only. Evolution of states in the infinite Hilbert space shows features quite different from that in finite Hilbert spaces.
We expect that the method introduced in this paper is also useful in studying more realistic models, such as the standard model and models including the gravity. In the application of this method to realistic models, masses of particles are to be explained as energies of ground states of interacting fields, therefore, there would be no need to resort to Higgs mechanism to get masses for particles. For the standard model, without introducing masses to free fields, the Hamiltonian of the model can not have any non-zero finite eigenenergy due to the lack of the dimension of mass. But, making use of the method discussed in this paper, it would be possible for one to get information on ratios of the eigenenergies and properties of the eigenstates of the Hamiltonian of the standard model. In order to obtain non-zero finite eigenenergies without introducing masses to free fields, one must include the gravity. One of the most serious problems one meets when including the gravity in quantum field theories is that the extended theories are generally non-renormalizable. However, as discussed in this paper, it is not impossible for non-renormalizable models to have energy eigenstates with finite energies when treated rigorously. Therefore, besides the string theory (see, e.g., ), the method introduced in this paper may supply another possible way of overcoming the difficulty of ultraviolet divergence. |
warning/0003/math0003005.html | ar5iv | text | # Ensembles d’unicité pour les polynômes 11footnote 1Classification mathématique: 30D05, 58F23. Mots clés: ensemble d’unicité, ensemble de Julia, mesure invariante.
## 1 Introduction
On note $`^1=\{\mathrm{}\}`$ la droite projective complexe. Un compact $`E`$ de $``$ est appelé ensemble d’unicité si pour tous polynômes non constants $`f`$ et $`g`$ vérifiant $`f^1(E)=g^1(E)`$ on a $`f=g`$. S’il existe un polynôme $`P\mathrm{id}`$ tel que $`P^1(E)=E`$, alors $`E`$ n’est pas un ensemble d’unicité.
###### Question 1
Supposons que $`P^1(E)E`$ pour tout polynôme $`P\mathrm{id}`$. $`E`$ est-t-il un ensemble d’unicité?
Les ensembles d’unicité pour les polynômes de même degré sont déterminés par Ostrovskii, Pakovitch et Zaidenberg . Les ensembles d’unicité pour les fonctions entières ou méromorphes avec un nombre minimal d’éléments sont étudiés par Nevanlinna et également par plusieurs autres auteurs (voir par exemple ).
###### Question 2
Soit $`\mu `$ une mesure de probabilité à support compact dans $``$. Pour quels polynômes $`f`$ et $`g`$ de degrés $`d1`$ et $`d^{}1`$ on a $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$?
Les deux questions posées ci-dessus sont étroitement liées. En effet, si la capacité logarithmique $`E`$ est positive et si $`f^1(E)=g^1(E)`$, alors $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$ pour la mesure d’équilibre $`\mu `$ de $`E`$. Réciproquement, si $`E`$ est le support de $`\mu `$ et si $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$, on a $`f^1(E)=g^1(E)`$. De plus, on peut trouver un compact $`E_0`$ de capacité logarithmique positive tel que $`f^1(E_0)=g^1(E_0)`$ (voir le paragraphe 3). Lorsque $`fg`$, $`E`$ et $`E_0`$ ne sont donc pas ensembles d’unicité. Les deux cas particuliers de ces problèmes sont le problème de détermination des fonctions ayant le même ensemble de Julia ou la même mesure totalement invariante et le problème de détermination des fonctions permutables (voir également ). Notre résultat principal est le théorème suivant:
###### Théorème 1
Soient $`\mu `$ une mesure de probabilité à support compact dans $``$ et $`f`$, $`g`$ deux polynômes de degrés $`d1`$ et $`d^{}1`$. Soit $`m`$ le plus grand diviseur commun de $`d`$ et $`d^{}`$. Supposons que $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$. Alors il existe un polynôme $`Q`$ de degré $`m`$ et des polynômes $`f_0`$, $`g_0`$ tels que $`f=f_0Q`$, $`g=g_0Q`$ et tels que l’une des conditions suivantes soit vraie:
1. $`f_0=\mathrm{id}`$ ou $`g_0=\mathrm{id}`$.
2. $`d>m`$, $`d^{}>m`$ et pour une certaine coordonnée $`z`$ de $``$ on a $`f_0(z)=z^{d/m}`$, $`g_0(z)=az^{d^{}/m}`$$`a0`$ est une constante.
3. $`d>m`$, $`d^{}>m`$ et pour une certaine coordonnée $`z`$ de $``$ on a $`f_0=\pm \mathrm{T}_{d/m}`$, $`g_0=\pm \mathrm{T}_{d^{}/m}`$$`\mathrm{T}_k`$ est le polynôme de Tchebychev de degré $`k`$.
###### Corollaire 1
Soit $`E`$ un compact de capacité logarithmique positive de $``$. Alors $`E`$ est un ensemble d’unicité si et seulement si pour tout polynôme $`P\mathrm{id}`$ on a $`P^1(E)E`$.
Soit $`P`$ un polynôme de degré au moins deux. L’ensemble de Julia rempli $`K_P`$ de $`P`$ est l’ensemble des points d’orbite bornée, i.e. les points $`z`$ tels que les suites $`\{P^n(z)\}_n`$ soient bornées. Ici on note $`P^n:=P\mathrm{}P`$ le $`n`$-ième itéré de $`P`$. L’ensemble de Julia $`J_P`$ est le bord topologique de l’ensemble $`K_P`$. Alors $`K_P`$ est le plus grand compact totalement invariant par $`P`$, i.e. $`P^1(K_P)=K_P`$. L’ensemble $`J_P`$ est le plus petit compact totalement invariant par $`P`$ qui contient plus qu’un élément (voir par exemple \[2, 4.2.2\]).
###### Corollaire 2
Soit $`E`$ un compact de capacité logarithmique positive de $``$. Supposons que $`E`$ n’est pas un ensemble d’unicité et que $`E`$ n’est pas invariant par aucune rotation de $``$. Alors il existe un polynôme $`P`$ de degré au moins deux tel que $`P^1(E)=E`$ et $`J_PEK_P`$.
Si dans un ouvert $`U`$ l’intersection $`J_PU`$ est un ensemble non vide, inclus dans une courbe réelle lisse, alors $`P(z)=z^d`$ ou $`P(z)=\pm \mathrm{T}_d(z)`$ pour une certaine coordonnée $`z`$ de $``$ \[15, p.127\]. La notation $`\mathrm{T}_d`$ signifie le polynôme de Tchebychev de degré $`d`$ défini par $`\mathrm{T}_d(\mathrm{cos}t):=\mathrm{cos}dt`$. Si $`P(z)=z^d`$, $`K_P`$ est le disque unité, $`J_P`$ est le cercle unité; si $`P(z)=\pm \mathrm{T}_d(z)`$, $`K_p`$ et $`J_p`$ sont égaux au segment $`[1,1]`$. On en déduit facilement que, toute courbe réelle lisse par morceaux, qui n’est invariante par aucune rotation, est un ensemble d’unicité.
L’ensemble $`J_P`$ est le support d’une mesure de probabilité $`\mu _P`$ qui est totalement invariante par $`P`$, i.e. $`(\mathrm{deg}P)^1P^{}(\mu _P)=\mu _P`$. C’est la seule mesure de probabilité à support compact qui est totalement invariante par $`P`$ sauf dans le cas où $`P(z)=z^d`$ pour une certaine coordonnée $`z`$ de $``$. Dans ce cas exceptionnel, toute mesure totalement invariante est une combinaison linéaire de la mesure de Lebesgue sur le cercle unité et de la masse de Dirac en $`0`$. Dans le théorème 1, si $`g_0=\mathrm{id}`$, $`\mathrm{deg}f_02`$ et si $`f_0`$ n’est pas conjugué à $`z^{d/m}`$ on a $`\mu =\mu _{f_0}`$. Dans la deuxième condition du théorème 1, $`\mu `$ est une combinaison linéaire de la masse de Dirac en $`0`$ et de la mesure de Lebesgue sur le cercle $`\{z:|z|=|a|^{d/(dd^{})}\}`$. Dans la troisième condition de ce théorème, $`\mu `$ est la mesure totalement invariante des polynômes de Tchebychev.
Dans la preuve du théorème 1, on se ramène à des systèmes dynamiques holomorphes en une et en plusieurs variables. D’abord, on peut choisir une fonction subharmonique $`\phi _0`$ telle que $`i\overline{}\phi _0=\mu `$ et $`d^1\phi _0f=d_{}^{}{}_{}{}^{1}\phi _0g=:\phi _1`$. Notons $`\delta _f(z):=\mathrm{exp}(2i\pi /d)z+a_0+a_1z^1+\mathrm{}`$ le germe d’application holomorphe défini au voisinage de $`\mathrm{}`$ qui préserve les fibres de $`f`$. Notons $`\delta _g`$ le germe analogue pour $`g`$. Alors $`\delta _f`$ et $`\delta _g`$ préservent les lignes de niveau de $`\phi _1`$. Dans une coordonnée locale convenable, les lignes de niveau de $`\phi _1`$ au voisinage de $`\mathrm{}`$ sont les cercles de centre $`\mathrm{}`$. Par conséquent, pour cette coordonnée locale, $`\delta _f`$ et $`\delta _g`$ sont des rotations. D’où $`\delta _f\delta _g=\delta _g\delta _f`$ et $`\delta _f^{d/m}=\delta _g^{d^{}/m}=:\delta `$. Notons $`_f(z)`$ la fibre de $`f`$ qui contient $`z`$, i.e. $`_f(z)=f^1f(z)`$. Le polynôme $`Q`$ sera défini comme un polynôme dont toute fibre générique $`_Q(z)`$ est égale à l’intersection $`_f(z)_g(z)`$. Au voisinage de $`\mathrm{}`$, la fibre $`_Q(z)`$ est l’orbite de $`z`$ par $`\delta `$. Ceci entraîne que $`Q`$ est bien défini et qu’il est de degré $`m`$ (proposition 2). Si $`m=d`$ ou $`m=d^{}`$, on peut choisir $`P=f`$ ou $`P=g`$; la condition 1 du théorème 1 est alors vraie. Pour la suite de la preuve, on peut supposer que $`d>1`$ et $`d^{}>1`$ sont premiers entre eux.
Notons $`\mathrm{\Phi }`$ un polynôme de degré $`dd^{}`$ vérifiant $`_\mathrm{\Phi }(z)=g^1g(_f(z))`$ pour tout $`z`$. Au voisinage de $`\mathrm{}`$, $`_\mathrm{\Phi }(z)`$ est l’ensemble des points $`\delta _g^n\delta _f^m(z)`$ car $`\delta _f`$ et $`\delta _g`$ commutent. Ceci entraîne que $`\mathrm{\Phi }`$ existe et que $`_\mathrm{\Phi }(z)=f^1f(_g(z))`$ (proposition 2). Alors on peut décomposer $`\mathrm{\Phi }=f_1g=g_1f`$$`f_1`$ (resp. $`g_1`$) est un polynôme de degré $`d`$ (resp. $`d^{}`$). En suite, on peut montrer que $`𝒞_{f_1}=g(𝒞_f)`$ et $`𝒞_{g_1}=f(𝒞_g)`$$`𝒞`$ signifie l’ensemble critique. Ceci, sous certaines conditions posées sur les coefficients dominants et sur les valeurs de polynômes en $`0`$, permet de construire un endomorphisme polynomial $`𝒟_{d,d^{}}`$ de l’ensemble $`\mathrm{\Sigma }(d,d^{})`$ des couples $`(f,g)`$. L’ensemble des couples $`(f,g)`$ qui vérifient les hypothèses du théorème 1, est un sous-ensemble algébrique $`𝒩(d,d^{})`$ invariant par $`𝒟_{d,d^{}}`$. L’ensemble des couples $`(f,g)`$ qui vérifient la condition 2 ou 3 du théorème 1, décrit deux courbes $`𝒞_1(d,d^{})`$ et $`𝒞_2(d,d^{})`$. Il reste à prouver que $`𝒩(d,d^{})=𝒞_1(d,d^{})𝒞_2(d,d^{})`$. On montrera que les points périodiques de $`𝒩(d,d^{})`$ appartiennent à $`𝒞_1(d,d^{})𝒞_2(d,d^{})`$. Ceci est dû à la solution d’une équation bien connue: $`fg=gf`$ (en particulier, pour les points fixes, on obtient directement $`\mathrm{\Phi }=fg=gf`$). Finalement, l’invariance de $`𝒩(d,d^{})`$ par $`𝒟_{d,d^{}}`$ implique que $`𝒩(d,d^{})=𝒞_1(d,d^{})𝒞_2(d,d^{})`$.
Remerciement.— Je tiens à remercier Charles Favre et Nessim Sibony pour leurs aides pendant la préparation de cet article.
## 2 Factorisation de polynômes
Soit $`P`$ un polynôme de degré $`d1`$. Il existe une fonction holomorphe unique $`B(z)=z+a_0+a_1z^1+a_2z^2+\mathrm{}`$ définie au voisinage de $`\mathrm{}`$ telle que $`BPB^1=\alpha z^d`$$`\alpha `$ est le coefficient dominant de $`P`$ (voir par exemple \[2, 6.10.1\]). On définit la fonction $`\delta _P`$ par la formule $`\delta _P(z):=B^1(\theta _dB(z))`$$`\theta _d:=\mathrm{exp}(2\pi i/d)`$. Cette fonction est définie au voisinage de $`\mathrm{}`$, permute les éléments de chaque fibre de $`P`$ et vérifie $`P\delta _P=P`$, $`\delta _P^d=\mathrm{id}`$. Soit $`\mathrm{\Phi }`$ une application définie dans un voisinage suffisamment petit de $`\mathrm{}`$ à l’image dans un espace $`X`$. Alors $`\mathrm{\Phi }`$ s’écrit sous la forme $`GP`$ si et seulement si $`\delta _P`$ préserve les fibres de $`\mathrm{\Phi }`$, i.e. $`\mathrm{\Phi }\delta _P=\mathrm{\Phi }`$.
###### Proposition 1
Soient $`X`$ un espace métrique, $`a`$ un point de $`X`$ et $`\mathrm{\Phi }`$ une application définie dans un voisinage suffisamment petit de $`\mathrm{}`$ à l’image dans $`X\{a\}`$ vérifiant $`lim_z\mathrm{}\mathrm{\Phi }(z)=a`$. Soient $`P_1`$ et $`P_2`$ deux polynômes de degrés $`d_11`$ et $`d_21`$ tels que $`\mathrm{\Phi }`$ s’écrive sous les formes $`\mathrm{\Phi }=G_1P_1=G_2P_2`$. Soit $`m`$ le plus grand diviseur commun de $`d_1`$ et $`d_2`$. Alors $`\delta _{P_1}\delta _{P_2}=\delta _{P_2}\delta _{P_1}`$ et $`\delta _{P_1}^{d_1/m}=\delta _{P_2}^{d_2/m}`$. En particulier, si $`d_1=d_2`$ il existe un polynôme linéaire $`\sigma `$ tel que $`P_1=\sigma P_2`$.
* Preuve— On montre que $`\delta _{P_1}\delta _{P_2}=\delta _{P_2}\delta _{P_1}`$. Supposons que $`\delta :=\delta _{P_1}^1\delta _{P_2}\delta _{P_1}\delta _{P_2}^1\mathrm{id}`$. Remarquons que $`\delta `$ s’écrit sous la forme $`\delta (z)=z+a_0+a_1z^1+a_2z^2+\mathrm{}`$. La dynamique d’une telle application est bien connue (voir par exemple \[2, 6.5\]). Il existe un point $`z`$ tel que $`\delta ^n(z)`$ tende vers $`\mathrm{}`$ quand $`n+\mathrm{}`$. Par conséquent, $`\mathrm{\Phi }(\delta ^n(z))`$ tend vers $`a`$. C’est une contradiction car $`\mathrm{\Phi }\delta =\mathrm{\Phi }`$. De même manière, on montre que $`\delta _{P_1}^{d_1/m}\delta _{P_2}^{d_2/m}=\mathrm{id}`$ et $`\delta _{P_1}^{d_1/m}=\delta _{P_2}^{d_2/m}`$.
Si $`d_1=d_2`$, on a $`\delta _{P_1}=\delta _{P_2}`$. Par conséquent, $`_{P_1}(z)=_{P_2}(z)`$ au voisinage de $`\mathrm{}`$. Par analyticité, ceci est vrai pour tout $`z`$. Alors on peut définir la fonction $`\sigma `$ holomorphe dans $``$ par $`\sigma :=P_1P_2^1`$. Il est clair que $`\sigma `$ est bijective. Par conséquent, $`\sigma `$ est linéaire et $`P_1=\sigma P_2`$.
$`\mathrm{}`$
Pour tout polynôme $`P`$, on note $`𝒞_P`$ l’ensemble critique de $`P`$. Un point de $`𝒞_P`$ sera compté $`k`$ fois s’il est de multiplicité $`k`$.
###### Proposition 2
Soient $`P_1`$ et $`P_2`$ deux polynômes de degrés $`d_11`$ et $`d_21`$. Soit $`m`$ le plus grand diviseur commun de $`d_1`$ et $`d_2`$.
1. Si $`\mathrm{\Phi }`$ est un polynôme vérifiant $`\mathrm{\Phi }\delta _{P_1}=\mathrm{\Phi }`$ au voisinage de $`\mathrm{}`$, alors il existe un polynôme $`R`$ tel que $`\mathrm{\Phi }=RP_1`$.
2. Si $`\delta _{P_1}^{d_1/m}=\delta _{P_2}^{d_2/m}`$ alors il existe un polynôme $`Q`$ de degré $`m`$ et des polynômes $`R_1`$, $`R_2`$ tels que $`P_1=R_1Q`$ et $`P_2=R_2Q`$. En particulier, si $`d^{}=m`$ il existe un polynôme $`R`$ tel que $`P_1=RP_2`$.
3. Si $`m=1`$ et si $`\delta _{P_1}\delta _{P_2}=\delta _{P_2}\delta _{P_1}`$, il existe un polynôme $`\mathrm{\Phi }`$ de degré $`d_1d_2`$ et des polynômes $`P_1^{}`$, $`P_2^{}`$ tels que $`\mathrm{\Phi }=P_1^{}P_2=P_2^{}P_1`$. De plus, on a $`𝒞_{P_1^{}}=P_2(𝒞_{P_1})`$ et $`𝒞_{P_2^{}}=P_1(𝒞_{P_2})`$.
* Preuve— 1. Le fait que $`\mathrm{\Phi }\delta _{P_1}`$ implique que $`_{P_1}(z)_\mathrm{\Phi }(z)`$ pour $`z`$ dans un voisinage de $`\mathrm{}`$. Par analyticité, ceci est vrai pour tout $`z`$. Par conséquent, on peut définir la fonction $`R`$ holomorphe dans $``$ par $`R(z):=\mathrm{\Phi }P^1(z)`$. On a $`\mathrm{\Phi }=RP_1`$. Comme $`\mathrm{\Phi }`$ et $`P`$ sont des polynômes, $`lim_z\mathrm{}R(z)=\mathrm{}`$. Donc $`R`$ est un polynôme.
2. On note $`w:=(w_1,w_2)`$ les coordonnées de $`^2`$. Soient $`\mathrm{\Pi }:^2`$ défini par $`\mathrm{\Pi }(z):=(P_1(z),P_2(z))`$ et $`𝒞:=\mathrm{\Pi }()`$. Alors $`𝒞`$ est une courbe algébrique de $`^2`$. De plus, elle est parabolique car elle est l’image de $``$ par une application holomorphe. D’où $`𝒞`$ est un $``$ ou un $`^{}`$ immergé dans $`^2`$. La courbe compactifiée $`\overline{𝒞}^2`$ est donc rationnelle. Comme $`P(z)/Q(z)`$ a une limite finie ou infinie quand $`z\mathrm{}`$, $`\overline{𝒞}`$ coupe la droite infinie $`L`$ en un seul point $`a`$. De plus, au voisinage de $`a`$, $`\overline{𝒞}`$ est irréductible car c’est l’image d’un voisinage de $`\mathrm{}`$ par l’application $`\mathrm{\Pi }`$. On en déduit que $`𝒞=\overline{𝒞}\{a\}`$ est un $``$ immergé. Soit $`\phi :𝒞`$ une application holomorphe, injective en dehors d’un nombre fini de points. On pose $`Q:=\phi ^1\mathrm{\Pi }`$. Alors $`Q`$ est une fonction holomorphe définie en dehors d’un nombre fini de points au voisinage desquels elle est bornée. Par conséquent, $`Q`$ se prolonge en une fonction holomorphe sur $``$. On vérifie facilement que $`lim_z\mathrm{}Q(z)=\mathrm{}`$. Donc $`Q`$ est un polynôme. On pose $`\delta :=\delta _{P_1}^{d_1/m}=\delta _{P_2}^{d_2/m}`$.
Au voisinage de $`\mathrm{}`$, on a
$$Q\delta =\phi ^1(P_1\delta ,P_2\delta )=\phi ^1(P_1,P_2)=Q.$$
Soient $`z_1`$ et $`z_2`$ suffisamment proches de $`\mathrm{}`$ tels que $`Q(z_1)=Q(z_2)`$. Alors $`P_1(z_1)=P_1(z_1)`$ et $`P_2(z_1)=P_2(z_2)`$. Il existe donc les entiers $`0n_1d_11`$ et $`0n_2d_21`$ tels que $`z_1=\delta _{P_1}^{n_1}(z_2)=\delta _{P_2}^{n_2}(z_2)`$. De plus, on sait que
$$\underset{z\mathrm{}}{lim}\frac{\delta _{P_j}^{n_j}(z)}{z}=\mathrm{exp}(2n_j\pi i/d_j).$$
Par conséquent, l’égalité $`\delta _{P_1}^{n_1}(z_2)=\delta _{P_2}^{n_2}(z_2)`$ pour $`z_2`$ suffisamment proche de $`\mathrm{}`$ implique que $`n_jm`$ est divisible par $`d_j`$ pour $`j=1`$ ou $`2`$. On a alors $`z_1=\delta ^{n_1m/d_1}(z_2)`$.
Les deux arguments ci-dessus montrent que $`\delta _Q=\delta `$. Par conséquent, $`\mathrm{deg}Q=m`$. D’après la première partie, il existe des polynômes $`R_1`$ et $`R_2`$ tels que $`P_1=R_1Q`$ et $`P_2=R_2Q`$.
On remarque que $`_Q(z)=_{P_1}(z)_{P_2}(z)`$ pour un $`z`$ générique car ceci est vrai au voisinage de $`\mathrm{}`$. On peut également prouver cette partie par la même méthode que l’on utilisera dans la troisième partie.
3. Comme $`m=1`$, il existe des entiers relatifs $`n_1`$ et $`n_2`$ tels que $`n_1d_2+n_2d_1=1`$. Posons $`\delta :=\delta _{P_1}^{n_1}\delta _{P_2}^{n_2}`$. Alors $`\delta `$ s’écrit sous la forme $`\delta (z)=\mathrm{exp}(2\pi i/dd^{})z+a_0+a_1z^1+\mathrm{}`$. Notons $`(z)=P_2^1P_2(_{P_1}(z))`$. Un point de $`(z)`$ est compté $`k`$ fois s’il est de multiplicité $`k`$. Alors $`(z)`$ est de cardinal $`dd^{}`$ et au voisinage de $`\mathrm{}`$, $`(z)`$ est l’orbite de $`z`$ par $`\delta `$ car $`\delta _{P_1}\delta _{P_2}=\delta _{P_2}\delta _{P_1}`$. On en déduit que $`(z)=P_1^1P_1(_{P_2}(z))`$ au voisinage de $`\mathrm{}`$. Par analyticité, ceci est vrai pour tout $`z`$.
Pour tout $`1ndd^{}1`$, on note $`S_n(z)`$ la somme symétrique des termes du type $`z_1\mathrm{}z_n`$ avec $`\{z_1,\mathrm{},z_n\}(z)`$. Alors $`S_n(z)`$ est une fonction holomorphe sur $``$. Il est clair que $`|\delta _{P_1}^i\delta _{P_2}^j(z)/z|`$ tend vers $`1`$ quand $`z\mathrm{}`$ pour tous $`i`$ et $`j`$. Par conséquent, $`S_n(z)=\text{O}(|z|^n)`$ quand $`z\mathrm{}`$. On en déduit que $`\mathrm{deg}S_nndd^{}1`$. Comme au voisinage de $`\mathrm{}`$, $`(z)`$ est l’orbite de $`z`$ par $`\delta `$, on a $`(z)=(z_1)`$ pour tout $`z_1(z)`$. Par conséquent, $`S_n`$ est constant sur $`(z)`$. Le fait que $`(z)`$ est de cardinal $`dd^{}>\mathrm{deg}S_n`$ implique que $`S_n`$ est un polynôme constant. Soit $`\mathrm{\Phi }`$ le polynôme de degré $`dd^{}`$ défini par:
$$\mathrm{\Phi }(z):=z^{dd^{}}S_1z^{dd^{}1}+\mathrm{}+(1)^{dd^{}1}S_{dd^{}1}z.$$
Alors au voisinage de $`\mathrm{}`$, $`_\mathrm{\Phi }(z)=(z)`$. Par analyticité, ceci est vrai pour tout $`z`$. Par conséquent, $`\mathrm{\Phi }\delta _{P_1}=\mathrm{\Phi }`$ et $`\mathrm{\Phi }\delta _{P_2}=\mathrm{\Phi }`$. D’après la première partie, il existe des polynômes $`P_1^{}`$ et $`P_2^{}`$ tels que $`\mathrm{\Phi }=P_1^{}P_2=P_2^{}P_1`$.
On a $`𝒞_\mathrm{\Phi }=𝒞_{P_2}P_2^1(𝒞_{P_1^{}})`$. Ici la notation $`𝒞`$ signifie l’ensemble critique et un point critique sera compté $`k`$ fois s’il est de multiplicité $`k`$. Remarquons qu’un point critique $`z`$ de $`\mathrm{\Phi }`$ est de multiplicité $`k`$ si et seulement si $`z`$ est un point de multiplicité $`k+1`$ de $`(z)`$. Comme $`(z)=P_2^1P_2(_{P_1}(z))`$, on a $`𝒞_\mathrm{\Phi }=𝒞_{P_2}P_2^1P_2(𝒞_{P_1})`$. Alors $`P_2^1(𝒞_{P_1^{}})=P_2^1P_2(𝒞_{P_1})`$ car $`𝒞_\mathrm{\Phi }=𝒞_{P_2}P_2^1(𝒞_{P_1^{}})`$. D’où $`𝒞_{P_1^{}}=P_2(𝒞_{P_1})`$. De même, on a $`𝒞_{P_2^{}}=P_1(𝒞_{P_2})`$.
$`\mathrm{}`$
Soient $`\mu `$ une mesure de probabilité à support compact et $`f`$, $`g`$ deux polynômes de degrés $`d`$ et $`d^{}`$ vérifiant $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$. Soit $`\phi `$ une fonction subharmonique vérifiant $`i\overline{}\phi =\mu `$. Cette fonction est harmonique sur la composante non bornée de $`\text{supp}(\mu )`$; elle est unique à une constante près. De plus, $`\phi \mathrm{ln}|z|`$ est harmonique et bornée au voisinage de $`\mathrm{}`$. Posons $`\psi :=d^1\phi f`$. Alors $`i\overline{}\psi =d^1f^{}(\mu )`$. Comme $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$, on a $`d_{}^{}{}_{}{}^{1}\phi g=\psi +c`$$`c`$ est une constante. Lorsque $`dd^{}`$, quitte à remplacer $`\phi `$ par $`\phi dd^{}c/(dd^{})`$, on peut supposer que $`c=0`$. Dans tous les cas, on peut appliquer la proposition 1 pour la fonction $`\psi `$. On obtient $`\delta _f\delta _g=\delta _g\delta _f`$ et $`\delta _f^{d/m}=\delta _g^{d^{}/m}`$$`m`$ est le plus grand diviseur commun de $`d`$ et $`d^{}`$. D’après la proposition 2, on a:
###### Corollaire 3
Il existe un polynôme $`Q`$ de degré $`m`$ et des polynômes $`f_0`$, $`g_0`$ tels que $`f=f_0Q`$, $`g=g_0Q`$ et $`(d/m)^1f_0^{}(\mu )=(d^{}/m)^1g_0^{}(\mu )`$. En particulier, si $`d^{}`$ divise $`d`$, il existe un polynôme $`P`$ de degré $`d/d^{}`$ tel que $`f=Pg`$ et $`d^1d^{}P^{}(\mu )=\mu `$.
## 3 Endomorphisme polynomial $`𝒟_{d,d^{}}`$
Soient $`\mu _0`$ une mesure de probabilité à support compact de $``$ et $`f_0`$, $`g_0`$ deux polynômes de degrés $`d>1`$, $`d^{}>1`$ vérifiant $`d^1f_0^{}(\mu _0)=d_{}^{}{}_{}{}^{1}g_0^{}(\mu _0)`$ . On suppose que $`d`$ et $`d^{}`$ sont premiers entre eux (voir le corollaire 3) et que $`d>d^{}`$. Soient $`\beta 0`$ et $`\alpha 0`$ les coefficients dominants de $`f_0`$ et $`g_0`$. On choisit un point $`a`$ tel que $`f_0(a)=g_0(a)`$. Soit $`b:=f_0(a)=g_0(a)`$. Quitte à remplacer $`f_0`$ par $`\sigma _1f_0\sigma _2`$, $`g_0`$ par $`\sigma _1g_0\sigma _2`$ et $`\mu _0`$ par $`(\sigma _1)_{}(\mu _0)`$ on peut supposer que $`a=b=0`$ et $`\beta =1`$$`\sigma _2(z):=Az+a`$, $`\sigma _1(z):=A^d\beta ^1(zb)`$ et $`A^{}`$. Soit $`\phi `$ la fonction subharmonique vérifiant $`i\overline{}\phi =\mu _0`$ et $`d^1\phi f_0=d_{}^{}{}_{}{}^{1}\phi g_0`$. Posons $`\phi _0(z):=\mathrm{max}(0,\phi (z))`$, $`\phi _1:=d^1\phi _0f_0`$, $`E_0:=\phi _0^1(0)`$ et $`E_1:=f_0^1(E_0)`$. Alors $`\phi _0`$ est subharmonique; $`\phi _1=d_{}^{}{}_{}{}^{1}\phi _0g_0`$ et $`E_1=g_0^1(E_0)`$. Comme $`\phi `$ tend vers l’infini quand $`z\mathrm{}`$, $`E_0`$ est compact. Alors $`\phi _0`$ est la fonction de Green de $`^1E_0`$ avec un seul pôle en $`\mathrm{}`$. On en déduit que $`E_0`$ est de capacité logarithmique positive (voir par exemple, \[16, III.8\]).
Notons $`\mathrm{\Sigma }(d,d^{},\alpha )`$ l’ensemble des couples $`(f,g)`$$`f`$ (resp. $`g`$) est un polynôme de degré $`d`$ (resp. $`d^{}`$) à coefficient dominant $`1`$ (resp. $`\alpha `$) qui s’annule en $`0`$.
###### Lemme 1
Il existe un couple unique $`(f_1,g_1)\mathrm{\Sigma }(d,d^{},\alpha ^d)`$ et un compact $`E_1`$ de capacité logarithmique positive tels que $`f_1g_0=g_1f_0`$ et tels que $`E_0=f_1^1(E_1)=g_1^1(E_1)`$. De plus, on a $`𝒞_{f_1}=g_0(𝒞_{f_0})`$ et $`𝒞_{g_1}=f_0(𝒞_{g_0})`$.
* Preuve— D’après le corollaire 3 et la proposition 2, il existe un polynôme $`\mathrm{\Phi }`$ de degré $`dd^{}`$ et des polynômes $`f_1`$ et $`f_2`$ tels que $`\mathrm{\Phi }=f_1g_0=g_1f_0`$. Quitte à remplacer $`\mathrm{\Phi }`$, $`f_1`$ et $`g_1`$ par $`\sigma \mathrm{\Phi }`$, $`\sigma f_1`$ et $`\sigma g_1`$, on peut supposer que $`\mathrm{\Phi }(0)=0`$ et que le coefficient dominant de $`\mathrm{\Phi }`$ soit $`\alpha ^d`$$`\sigma `$ est un certain polynôme linéaire. On a alors $`(f_1,g_1)\mathrm{\Sigma }(d,d^{},\alpha ^d)`$.
Montrons qu’au voisinage de $`\mathrm{}`$, $`\delta _{f_1}`$ préserve les lignes de niveau de $`\phi _0`$. Soient $`a_1`$ et $`a_2`$ suffisamment proches de $`\mathrm{}`$ tels que $`f_1(a_1)=f_1(a_2)`$. Il faut prouver que $`\phi _0(a_1)=\phi _0(a_2)`$. Il existe $`b_1`$ et $`b_2`$ tels que $`g_0(b_1)=a_1`$ et $`g_0(b_2)=(a_2)`$. Alors $`\mathrm{\Phi }(b_1)=\mathrm{\Phi }(b_2)`$. Par construction de $`\mathrm{\Phi }`$ (voir la preuve de la proposition 2), il existe $`m`$ et $`n`$ tels que $`b_1=\delta _{f_0}^m\delta _{g_0}^n(b_2)`$. Comme $`\phi _1=d^1\phi _0f_0=d_{}^{}{}_{}{}^{1}\phi _0g_0`$, les applications $`\delta _{f_0}`$ et $`\delta _{g_0}`$ préservent les lignes de niveau de $`\phi _1`$. D’où $`\phi _1(b_1)=\phi _1(b_2)`$. On obtient
$$\phi _0(a_1)=\phi _0g_0(b_1)=d^{}\phi _1(b_1)=d^{}\phi _1(b_2)=\phi _0g_0(b_2)=\phi _0(a_1).$$
Alors au voisinage de $`\mathrm{}`$, $`\delta _{f_1}`$ préserve les lignes de niveau de $`\phi _0`$, i.e. les lignes de niveau de $`\phi _0`$ sont réunions de fibres de $`f_1`$. Comme $`\phi _0`$ est harmonique dans $`E_0=\phi _0^1(0)`$, elle est réelle analytique dans $`E_0`$. Par analyticité, $`E_0`$ est une réunion de fibres de $`f_1`$. Par conséquent, $`E_0`$ est une réunion de fibres de $`f_1`$. Posons $`E_1:=f_1(E_0)`$. Alors $`E_0=f_1^1(E_1)`$. Il est clair que $`E_1`$ est de capacité logarithmique positive. Les relations $`f_1g_0=g_1f_0`$ et $`f_0^1(E_0)=g_0^1(E_0)`$ entraînent $`g_1^1(E_1)=E_0`$. Les polynômes $`f_1`$ et $`g_1`$ sont uniques car la fonction $`\mathrm{\Phi }`$ est unique (voir la preuve de la proposition 2). D’après la proposition 2, on a $`𝒞_{f_1}=g_0(𝒞_{f_0})`$ et $`𝒞_{g_1}=f_0(𝒞_{g_0})`$.
$`\mathrm{}`$
###### Remarque 1
1. On peut construire les couples $`(f_k,g_k)\mathrm{\Sigma }(d,d^{},\alpha ^{d^k})`$ et les compacts $`E_k`$ tels que $`f_kg_{k1}=g_kf_{k1}`$ et $`E_{k1}=f_k^1(E_k)=g_k^1(E_k)`$.
2. On fixe un $`k1`$ et un $`m0`$. On pose $`\stackrel{~}{d}:=d^k`$, $`\stackrel{~}{d}^{}:=d_{}^{}{}_{}{}^{k}`$, $`\stackrel{~}{\alpha }:=\alpha ^{d^m(d^{k1}+d^{k2}d^{}+\mathrm{}+d_{}^{}{}_{}{}^{k1})}`$, $`\stackrel{~}{f}_i:=f_{ik+k+m1}\mathrm{}f_{ik+m}`$, $`\stackrel{~}{g}_i:=g_{ik+k+m1}\mathrm{}g_{ik+m}`$, $`\stackrel{~}{E}_i:=E_{ik+k+m1}`$ pour $`i=0`$ ou 1. Alors $`(\stackrel{~}{f}_0,\stackrel{~}{g}_0)\mathrm{\Sigma }(\stackrel{~}{d},\stackrel{~}{d}^{},\stackrel{~}{\alpha })`$ et $`(\stackrel{~}{f}_1,\stackrel{~}{g}_1)\mathrm{\Sigma }(\stackrel{~}{d},\stackrel{~}{d}^{},\stackrel{~}{\alpha }^{\stackrel{~}{d}})`$. On vérifie facilement que $`\stackrel{~}{f}_1\stackrel{~}{g}_0=\stackrel{~}{g}_1\stackrel{~}{f}_0`$ et que $`\stackrel{~}{E}_0=\stackrel{~}{f}_1^1(\stackrel{~}{E}_1)=\stackrel{~}{g}_1^1(\stackrel{~}{E}_1)`$. Par l’unicité, $`(\stackrel{~}{f}_1,\stackrel{~}{g}_1)`$ est le couple que l’on peut construire comme dans le lemme 1 mais pour les polynômes $`\stackrel{~}{f}_0`$ et $`\stackrel{~}{g}_0`$.
On remarque qu’un couple $`(f,g)\mathrm{\Sigma }(d,d^{},\alpha )`$ est déterminé uniquement par les points critiques de $`f`$ et de $`g`$. Notons $`\mathrm{\Pi }_{d,d^{},\alpha }:^{d1}\times ^{d^{}1}\mathrm{\Sigma }(d,d^{},\alpha )`$ l’application qui associe un point $`(x,y)=(x_1,\mathrm{},x_{d1},y_1,\mathrm{},y_{d^{}1})`$ le couple $`(f,g)\mathrm{\Sigma }(d,d^{},\alpha )`$ vérifiant $`𝒞_f=\{x_1,\mathrm{},x_{d1}\}`$ et $`𝒞_g=\{y_1,\mathrm{},y_{d^{}1}\}`$. Cette application définit un revêtement ramifié au-dessus de $`\mathrm{\Sigma }(d,d^{},\alpha )`$. On définit l’application $`𝒟_{d,d^{}}:^{d1}\times ^{d^{}1}\times ^{}^{d1}\times ^{d^{}1}\times ^{}`$ par:
$$𝒟_{d,d^{}}(x,y,\alpha ):=(g(x_1),\mathrm{},g(x_{d1}),f(y_1),\mathrm{},f(y_{d^{}1}),\alpha ^d)$$
$`(f,g):=\mathrm{\Pi }_{d,d^{},\alpha }(x,y)`$. Les polynômes $`f`$ et $`g`$ sont déterminés par les formules explicites suivantes:
$$f(z)=\frac{1}{d}_0^z(tx_1)\mathrm{}(tx_{d1})𝑑t$$
et
$$g(z)=\frac{\alpha ^d}{d^{}}_0^z(ty_1)\mathrm{}(ty_{d^{}1})𝑑t.$$
Il est clair que $`𝒟_{d,d^{}}`$ est un endomorphisme polynomial.
###### Remarque 2
1. D’après le lemme précédent, si $`\mathrm{\Pi }_{d,d^{},\alpha }(x,y)=(f_0,g_0)`$ on a $`\mathrm{\Pi }_{d,d^{},\alpha ^d}(x^{},y^{})=(f_1,g_1)`$$`(x^{},y^{},\alpha ^d):=𝒟_{d,d^{}}(x,y,\alpha )`$.
2. D’après la remarque 1, si $`\mathrm{\Pi }_{\stackrel{~}{d},\stackrel{~}{d}^{},\stackrel{~}{\alpha }}(\stackrel{~}{x},\stackrel{~}{y})=(\stackrel{~}{f}_0,\stackrel{~}{g}_0)`$ on a $`\mathrm{\Pi }_{\stackrel{~}{d},\stackrel{~}{d}^{},\stackrel{~}{\alpha }}(\stackrel{~}{x}^{},\stackrel{~}{y}^{})=(\stackrel{~}{f}_1,\stackrel{~}{g}_1)`$$`(\stackrel{~}{x}^{},\stackrel{~}{y}^{},\stackrel{~}{\alpha }^{\stackrel{~}{d}}):=𝒟_{\stackrel{~}{d},\stackrel{~}{d}^{}}(\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{\alpha })`$.
## 4 Ensemble invariant $`𝒩(d,d^{})`$
Notons $`(d,d^{})`$ l’ensemble des points $`(x,y,\alpha )^{d1}\times ^{d^{}1}\times ^{}`$ vérifiant $`f^{}g=g^{}f`$$`(f,g):=\mathrm{\Pi }_{d,d^{},\alpha }(x,y)`$, $`(x^{},y^{},\alpha ^d):=𝒟_{d,d^{}}(x,y,\alpha )`$ et $`(f^{},g^{}):=\mathrm{\Pi }_{d,d^{},\alpha ^d}(x^{},y^{})`$. Notons $`𝒩(d,d^{})`$ l’ensemble des $`(x,y,\alpha )(d,d^{})`$ vérifiant les deux propriétés suivantes:
* 1. $`𝒫(d,d^{})`$: pour tout $`n0`$, on a $`𝒟_{d,d^{}}^n(x,y,\alpha )(d,d^{})`$.
* 2. Pour tous $`k1`$ et $`m0`$, si $`\mathrm{\Pi }_{\stackrel{~}{d},\stackrel{~}{d}^{},\stackrel{~}{\alpha }}(\stackrel{~}{x},\stackrel{~}{y})=(\stackrel{~}{f},\stackrel{~}{g})`$ alors $`(\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{\alpha })`$ vérifie la condition $`𝒫(\stackrel{~}{d},\stackrel{~}{d}^{})`$$`(x_n,y_n,\alpha _n):=𝒟_{d,d^{}}^n(x,y,\alpha )`$, $`(f_n,g_n):=\mathrm{\Pi }_{d,d^{},\alpha _n}(x_n,y_n)`$, $`\stackrel{~}{f}:=f_{k+m1}\mathrm{}f_m`$, $`\stackrel{~}{g}:=g_{k+m1}\mathrm{}g_m`$, $`\stackrel{~}{d}:=d^k`$, $`\stackrel{~}{d}^{}:=d_{}^{}{}_{}{}^{k}`$ et $`\stackrel{~}{\alpha }:=\alpha ^{d^m(d^{k1}+d^{k2}d^{}+\mathrm{}+d_{}^{}{}_{}{}^{k1})}`$ (voir la remarque 2).
Alors $`𝒩(d,d^{})`$ est un sous-ensemble algébrique faiblement invariant par $`𝒟_{d,d^{}}`$ i.e. $`𝒟_{d,d^{}}(𝒩(d,d^{}))𝒩(d,d^{})`$. De plus, $`𝒟_{d,d^{}}^n(𝒩(d,d^{}))`$ est faiblement invariant par $`𝒟_{d,d^{}}`$ pour tout $`n0`$.
Soient $`f_0`$, $`g_0`$, $`\alpha `$ et $`E_0`$ vérifiant les hypothèses du paragraphe précédent. D’après le lemme 1 et les remarques 1, 2, on a $`(x,y,\alpha )𝒩(d,d^{})`$ pour tout $`(x,y)`$ vérifiant $`\mathrm{\Pi }_{d,d^{},\alpha }(x,y)=(f_0,g_0)`$.
Nous construisons maintenant deux sous-ensembles $`𝒞_1(d,d^{})`$ et $`𝒞_2(d,d^{})`$ de $`𝒩(d,d^{})`$ grâce à des exemples précis sur $`f_0`$, $`g_0`$, $`\alpha `$ et $`E_0`$. Par suite, on montre que $`𝒩(d,d^{})=𝒞_1(d,d^{})𝒞_2(d,d^{})`$.
Soient $`\sigma _1`$, $`\sigma _2`$ deux polynômes linéaires, $`a0`$ et $`\alpha 0`$ tels que $`(f,g)\mathrm{\Sigma }(d,d^{},\alpha )`$$`f(z):=\sigma _1(z^d)\sigma _2`$ et $`g(z):=\sigma _1(az^d^{})\sigma _2`$. On a $`f^1(E)=g^1(E)`$ pour $`E:=\sigma _1(\{z:|z|=|a|^{d/(dd^{})}\})`$. Par conséquent, $`(x,y,\alpha )𝒩(d,d^{})`$ pour tout $`(x,y)`$ vérifiant $`\mathrm{\Pi }_{d,d^{},\alpha }(x,y)=(f,g)`$. On note $`𝒞_1(d,d^{})`$ l’ensemble de ces points $`(x,y,\alpha )`$.
Notons $`\mathrm{T}_k`$ le polynôme de Tchebychev de degré $`k`$ défini par $`\mathrm{T}(\mathrm{cos}z):=\mathrm{cos}kz`$. On sait que l’ensemble de Julia de $`\mathrm{T}_k`$ est l’intervalle $`[1,1]`$, que le coefficient dominant de $`\mathrm{T}_k`$ est égal à $`2^{k1}`$ et que les points critiques de $`\mathrm{T}_k`$ sont les points $`\mathrm{cos}t\pm 1`$ avec $`t`$ vérifiant $`\mathrm{sin}kt=0`$. Soient $`\sigma _1`$, $`\sigma _2`$ deux polynômes linéaires et $`\alpha ^{}`$ tels que $`(f,g)\mathrm{\Sigma }(d,d^{},\alpha )`$$`f:=\sigma _1(\pm \mathrm{T}_d)\sigma _2`$, $`g:=\sigma _1(\pm \mathrm{T}_d^{})\sigma _2`$. Posons $`E:=\sigma _1([1,1])`$. Alors $`f^1(E)=g^1(E)`$. On en déduit que $`(x,y,\alpha )𝒩(d,d^{})`$. Notons $`𝒞_2(d,d^{})`$ l’ensemble de tels points $`(x,y,\alpha )`$.
###### Lemme 2
$`𝒞_1(d,d^{})`$ et $`𝒞_2(d,d^{})`$ sont des courbes algébriques réductibles dont aucune composante n’est incluse dans un hyperplan du type $`\{\alpha =\text{constante}\}`$.
* Preuve— Soient $`\sigma _1(z)=a_1z+b_1`$, $`\sigma _2(z)=a_2z+b_2`$ et $`f`$, $`g`$, $`\alpha `$, $`x`$, $`y`$ définis ci-dessus.
Pour la courbe $`𝒞_1(d,d^{})`$, comme $`(f,g)\mathrm{\Sigma }(d,d^{},\alpha )`$, on obtient les relations suivantes: $`a_1a_2^d=1`$, $`a_1aa_2^d^{}=\alpha `$ et $`a_1b_2^d+b_1=a_1ab_2^d^{}+b_1=0`$. On a également, $`x=(b_2/a_2,\mathrm{},b_2/a_2)`$ et $`y=(b_2/a_2,\mathrm{},b_2/a_2)`$. On obtient facilement que $`b_2=0`$ ou $`\alpha =(b_2/a_2)^{dd^{}}`$. Ceci montre que $`𝒞_1`$ est une courbe algébrique réductible et qu’aucune de ses composantes n’est incluse dans $`\{\alpha =\text{constante}\}`$.
Pour la courbe $`𝒞_2(d,d^{})`$, on obtient $`a_12^{d1}a_2^d=1`$, $`a_12^{d^{}1}a_2^d^{}=\alpha `$ et $`\pm a_1\mathrm{T}_d(b_2)+b_1=\pm a_1\mathrm{T}_d^{}(b_2)+b_1=0`$. On a $`\{x_1,\mathrm{},x_{d1}\}=\sigma _2^1(𝒞_{\mathrm{T}_{d_2}})`$, $`\{y_1,\mathrm{},y_{d^{}1}\}=\sigma _2^1(𝒞_{\mathrm{T}_{d_2}})`$. On remarque que $`b_2`$ est une solution de l’équation $`\pm \mathrm{T}_d(z)=\pm \mathrm{T}_d^{}(z)`$. Cette équation n’a qu’un nombre fini de solution. Les autres nombres $`a_1`$ et $`a_2`$ (resp. $`b_1`$) s’écrivent en fonction de $`\alpha `$ (resp. de $`\alpha `$ et de $`b_2`$). Par conséquent, $`𝒞_2(d,d^{})`$ est une courbe algébrique dont aucune composante n’est incluse dans un hyperplan du type $`\{\alpha =\text{constante}\}`$.
$`\mathrm{}`$
###### Lemme 3
* 1. $`𝒟_{d,d^{}}^1(𝒞_1(d,d^{}))𝒩(d,d^{})𝒞_1(d,d^{})`$.
* 2. $`𝒟_{d,d^{}}^1(𝒞_2(d,d^{}))𝒩(d,d^{})𝒞_2(d,d^{})`$.
* Preuve— 1. Soit $`(x,y,\alpha )𝒟_{d,d^{}}^1(𝒞_1(d,d^{}))𝒩(d,d^{})`$. Posons $`(f,g):=\mathrm{\Pi }_{d,d^{},\alpha }(x,y)`$, $`(x_1,y_1,\alpha _1):=𝒟_{d,d^{}}(x,y,\alpha )`$ et $`(f_1,g_1):=\mathrm{\Pi }_{d,d^{},\alpha _1}(x_1,y_1)`$. Comme $`(x_1,y_1,\alpha _1)𝒞_1(d,d^{})`$, il existe des polynômes linéaires $`\sigma _1`$, $`\sigma _2`$ et une constante non nulle $`a`$ tels que $`f_1(z)=\sigma _1([\sigma _2(z)]^d)`$ et $`g_1=\sigma _1(a[\sigma _2(z)^d^{}])`$. Comme $`(x,y,\alpha )𝒩(d,d^{})`$, on a $`f_1g=g_1f`$. Posons $`f^{}:=\sigma _2f`$ et $`g^{}:=\sigma _2g`$. Alors $`a[f^{}(z)]^d^{}=[g^{}(z)]^d`$. Posons $`\mathrm{\Phi }(z):=a[f^{}(z)]^d^{}=[g^{}(z)]^d`$. Soit $`\lambda `$ une racine de $`\mathrm{\Phi }`$. Alors la multiplicité de $`\lambda `$ est divisible par $`d`$ et par $`d^{}`$. Comme $`d`$ et $`d^{}`$ sont premiers entre eux, la multiplicité de $`\lambda `$ est divisible par $`dd^{}`$. D’autre part, $`\mathrm{deg}\mathrm{\Phi }=dd^{}`$. On déduit que $`\lambda `$ est la seule racine de $`\mathrm{\Phi }`$. Il est également la seule racine de $`f^{}`$ et de $`g^{}`$. Alors il existe un polynôme linéaire $`\sigma `$ et un $`b`$ tels que $`f^{}(z)=[\sigma (z)]^d`$ et $`g^{}(z)=b[\sigma (z)]^d^{}`$. D’où $`(x,y,\alpha )𝒞_1(d,d^{})`$.
2. De même manière, on se ramène à une équation du type $`\mathrm{T}_dg^{}=\pm \mathrm{T}_d^{}f^{}`$. Il faut montrer qu’il existe un polynôme linéaire $`\sigma `$ tel que $`f^{}=\pm \mathrm{T}_d\sigma `$ et $`g^{}=\pm \mathrm{T}_d^{}\sigma `$. Il est clair que $`f_{}^{}{}_{}{}^{1}([1,1])=g_{}^{}{}_{}{}^{1}([1,1])`$. On déduit de la définition de $`𝒟_{d,d^{}}`$ que $`g^{}(𝒞_f^{})=𝒞_{\mathrm{T}_d}`$. Par conséquent, les points critiques de $`f^{}`$ sont tous de multiplicité 1. De même pour $`g^{}`$.
Pour $`|z|`$ suffisamment grand on a $`_f^{}(z)_g^{}(z)=\{z\}`$. En effet, utilisant les développements asymptotiques de $`\delta _f^{}`$ et $`\delta _g^{}`$, on obtient pour tous $`1nd1`$ et $`1md^{}1`$:
$$\underset{z\mathrm{}}{lim}\frac{\delta _f^{}^m(z)}{\delta _g^{}^n(z)}=\mathrm{exp}(2m\pi i/d2n\pi i/d^{})0$$
car $`d`$ et $`d^{}`$ sont premiers entre eux. Par analyticité, pour un $`z`$ générique $`_f^{}(z)_g^{}(z)=\{z\}`$. Soit $`p𝒞_f^{}`$. Montrons que $`f^{}(p)=\pm 1`$. Supposons que $`f^{}(p)=a\pm 1`$. On sait que $`g^{}(𝒞_f^{})=𝒞_{\mathrm{T}_d}[1,1]`$ et $`f_{}^{}{}_{}{}^{1}([1,1])=g_{}^{}{}_{}{}^{1}([1,1])`$. D’où $`a]1,1[`$. Alors au voisinage de $`p`$, $`f_{}^{}{}_{}{}^{1}([1,1])`$ est la réunion de deux courbes réelles analytiques qui se coupent en $`p`$. Par conséquent, $`p`$ est un point critique de $`g^{}`$. Comme $`\delta _f^{}`$ et $`\delta _g^{}`$ commutent, leurs prolongements analytiques commutent aussi au voisinage de $`p`$. On en déduit que $`_f^{}(z)_g^{}(z)`$ contient au moins deux éléments pour tout $`z`$ suffisamment proche de $`p`$. C’est une contradiction. Donc $`f^{}(p)=\pm 1`$. De même $`g^{}(q)=\pm 1`$ pour $`q𝒞_g^{}`$.
Le fait que $`f^{}`$ est de degré $`d`$ implique que pour $`d`$ impair $`f_{}^{}{}_{}{}^{1}(1)`$ est une réunion de $`(d1)/2`$ points critiques et d’un point non critique; pour $`d`$ pair $`f_{}^{}{}_{}{}^{1}(1)`$ est la réunion de $`d/2`$ points critiques ou la réunion de $`d/21`$ points critiques avec deux points non critiques. De même pour $`f_{}^{}{}_{}{}^{1}(1)`$. Quitte à remplacer $`f^{}`$, $`g^{}`$ par $`\pm f^{}\sigma `$ et $`g^{}\sigma `$ pour un certain polynôme linéaire $`\sigma `$, on peut supposer que $`\pm 1`$ ne sont pas critiques pour $`f`$ et que pour $`d`$ impair $`f^{}(1)=1`$, $`f^{}(1)=1`$ et pour $`d`$ pair $`f^{}(1)=f^{}(1)=1`$. On remarque qu’au voisinage de $`\pm 1`$, $`f_{}^{}{}_{}{}^{1}([1,1])`$ est un arc réel analytique. Par conséquent, $`g^{}(\pm 1)=\pm 1`$ et $`\pm 1`$ ne sont pas critiques pour $`g^{}`$. Quitte à remplacer $`g^{}`$ par $`\pm g^{}`$, on peut supposer que pour $`d^{}`$ impair $`g^{}(1)=1`$, $`g^{}(1)=1`$ et pour $`d^{}`$ pair $`g^{}(1)=g^{}(1)=1`$.
Alors il existe des polynômes $`P`$, $`Q`$ tels que pour $`d`$ impair $`f^{}(z)+1=(z+1)P^2(z)`$, $`f^{}(z)1=(z1)Q^2(z)`$ et pour $`d`$ pair $`f^{}(z)+1=P^2(z)`$, $`f^{}(z)1=(z1)(z+1)Q^2(z)`$. Posons $`\psi (z):=(z+z^1)/2`$. On vérifie facilement qu’il existe une fonction rationnelle $`R(z)`$ telle que:
$$\frac{f^{}+1}{f^{}1}\psi (z)=R^2(z).$$
On en déduit que $`f^{}(z)\psi =(F+F^1)/2`$$`F:=(R+1)/(R1)`$. On a donc $`\psi ^1f^{}\psi =F^{\pm 1}`$. Ceci implique que $`\mathrm{deg}F=d`$. De même, il existe une fonction rationnelle $`G`$ de degré $`d^{}`$ telle que $`\psi ^1g^{}\psi =G^{\pm 1}`$. D’autre part, $`\psi ^1\mathrm{T}_d\psi (z)=z^{\pm d}`$ et $`\psi ^1\mathrm{T}_d^{}\psi (z)=z^{\pm d^{}}`$. On déduit de la relation $`\mathrm{T}_dg^{}=\pm \mathrm{T}_d^{}f^{}`$ que $`F^{\pm d^{}}=\pm G^{\pm d}`$. Alors comme dans la partie précédente, les multiplicités des zéros et des pôles de $`F^{\pm d^{}}=\pm G^{\pm d}`$ sont divisibles par $`dd^{}`$. Or c’est une fonction de degré $`dd^{}`$. D’où $`F(z)=Az^{\pm d}`$ et $`G(z)=Bz^{\pm d^{}}`$. Le fait que $`f^{}(1)=g^{}(1)=1`$ entraîne $`A=B=1`$. D’où $`f^{}=\mathrm{T}_d`$ et $`g^{}=\mathrm{T}_d^{}`$.
$`\mathrm{}`$
###### Lemme 4
Soit $`(x,y,\alpha )𝒩(d,d^{})`$ un point prépériodique de $`𝒟_{d,d^{}}`$, i.e. $`𝒟_{d,d^{}}^{k+m}(x,y,\alpha )=𝒟_{d,d^{}}^m(x,y,\alpha )`$ pour certains $`k1`$ et $`m0`$. Alors $`(x,y,\alpha )`$ appartient à $`𝒞_1(d,d^{})𝒞_2(d,d^{})`$.
* Preuve— On utilise les notations de la définition de l’ensemble $`𝒩(d,d^{})`$. Soient $`(\stackrel{~}{x}^{},\stackrel{~}{y}^{},\stackrel{~}{\alpha }^{}):=𝒟_{\stackrel{~}{d},\stackrel{~}{d}^{}}(\stackrel{~}{x},\stackrel{~}{y},\stackrel{~}{\alpha })`$ et $`(\stackrel{~}{f}^{},\stackrel{~}{g}^{}):=\mathrm{\Pi }_{\stackrel{~}{d},\stackrel{~}{d}^{},\stackrel{~}{\alpha }}(\stackrel{~}{x}^{},\stackrel{~}{y}^{})`$. Comme $`𝒟_{d,d^{}}^{k+m}(x,y,\alpha )=𝒟_{d,d^{}}^m(x,y,\alpha )`$, on a $`\stackrel{~}{f}^{}=\stackrel{~}{f}`$ et $`\stackrel{~}{g}^{}=\stackrel{~}{g}`$ (voir la remarque 2). Par définition de $`𝒩(\stackrel{~}{d},\stackrel{~}{d}^{})`$, on a $`\stackrel{~}{f}^{}\stackrel{~}{g}=\stackrel{~}{g}^{}\stackrel{~}{f}`$. D’où $`\stackrel{~}{f}\stackrel{~}{g}=\stackrel{~}{g}\stackrel{~}{f}`$. Cette équation a été résolue par Fatou et Julia . Dans notre cas, $`\stackrel{~}{d}>1`$ et $`\stackrel{~}{d}^{}>1`$ sont premiers entre eux. D’après le théorème de Fatou-Julia, il existe un polynôme linéaire $`\sigma _1`$ tel que l’une des conditions suivantes soit vraie:
1. $`\sigma _1\stackrel{~}{f}\sigma _1^1=z^{\stackrel{~}{d}}`$ et $`\sigma _1\stackrel{~}{g}\sigma _1^1=az^{\stackrel{~}{d}^{}}`$$`a0`$ est une constante.
2. $`\sigma _1\stackrel{~}{f}\sigma _1^1=\pm \mathrm{T}_{\stackrel{~}{d}}`$ et $`\sigma _1\stackrel{~}{g}\sigma _1^1=\pm \mathrm{T}_{\stackrel{~}{d}^{}}`$.
Considérons le second cas, le premier cas sera traité de même manière. On remarque que $`\mathrm{T}_{rs}=\mathrm{T}_r\mathrm{T}_s`$. En particulier, $`\mathrm{T}_{\stackrel{~}{d}}=\mathrm{T}_{\stackrel{~}{d}/d}\mathrm{T}_d`$. Le fait que $`\stackrel{~}{f}=f_{k+m1}\mathrm{}f_m`$ implique
$$(\sigma _1f_{k+m1}\sigma _1^1)\mathrm{}(\sigma _1f_m\sigma _1^1)=\sigma _1\stackrel{~}{f}\sigma ^1=\pm \mathrm{T}_{\stackrel{~}{d}}.$$
D’après la proposition 1, il existe un polynôme linéaire $`\sigma _2`$ tel que $`\sigma _1f_m\sigma _1^1=\sigma _2\mathrm{T}_d`$. De même, il existe $`\sigma _2^{}`$ tel que $`\sigma _1g_m\sigma _1^1=\sigma _2^{}\mathrm{T}_d^{}`$. Alors $`f_m=\sigma _3\mathrm{T}_d\sigma _1`$ et $`g_m=\sigma _3^{}\mathrm{T}_d^{}\sigma _1`$$`\sigma _3:=\sigma _1^1\sigma _2`$ et $`\sigma _3^{}:=\sigma _1^1\sigma _2`$. On sait que pour $`d`$ et $`d^{}`$ premiers entre eux l’ensemble critique de $`\mathrm{T}_d`$ (resp. de $`\mathrm{T}_d^{}`$) est invariant par $`\mathrm{T}_d^{}`$ (resp. par $`\mathrm{T}_d`$). Par construction de $`𝒟_{d,d^{}}`$, on a
$$𝒞_{f_{m+1}}=g_m(𝒞_{f_m})=\sigma _3^{}\mathrm{T}_d^{}(𝒞_{\mathrm{T}_d})=\sigma _3^{}(𝒞_{\mathrm{T}_d}).$$
Alors il existe un polynôme linéaire $`\sigma _4`$ tel que $`f_{m+1}=\sigma _4\mathrm{T}_d\sigma _{3}^{}{}_{}{}^{1}`$. De même, il existe un polynôme linéaire $`\sigma _4^{}`$ tel que $`g_{m+1}=\sigma _4^{}\mathrm{T}_d^{}\sigma _3^1`$.
Comme on a montré ci–dessus pour $`f_m`$ et $`g_m`$, il suffit de remplacer $`m`$ par $`m+1`$ afin d’obtenir $`f_{m+1}=\sigma _6\mathrm{T}_d\sigma _5`$ et $`g_{m+1}=\sigma _6^{}\mathrm{T}_d\sigma _5`$$`\sigma _5`$, $`\sigma _6`$ et $`\sigma _6^{}`$ sont linéaires. On déduit des quatres dernières égalités que l’ensemble critique de $`\mathrm{T}_d`$ (resp. de $`\mathrm{T}_d^{}`$) est invariant par $`\sigma _5\sigma _3^{}`$ (resp. par $`\sigma _5\sigma _3`$). D’où $`\sigma _5\sigma _3^{}(z)=\pm z`$ et $`\sigma _5\sigma _3(z)=\pm z`$. Par conséquent, $`\sigma _3^{}(z)=\sigma _3(\pm z)`$ et donc $`f_m=\sigma _3\mathrm{T}_d\sigma _1`$ et $`g_m=\sigma _3(\pm \mathrm{T}_d^{})\sigma _1`$. Ceci signifie que $`(f_m,g_m,\alpha _m)𝒞_2(d,d^{})`$. D’après le lemme 3, $`(f,g,\alpha )𝒞_2(d,d^{})`$ car $`(f,g,\alpha )=(f_0,g_0,\alpha _0)`$.
$`\mathrm{}`$
###### Proposition 3
On a $`𝒩(d,d^{})=𝒞_1(d,d^{})𝒞_2(d,d^{})`$.
Soit $`S`$ un sous-ensemble algébrique périodique de $`𝒩(d,d^{})`$, i.e. $`𝒟_{d,d^{}}^n(S)=S`$ pour un certain $`n1`$. On montre que $`S𝒞_1(d,d^{})𝒞_2(d,d^{})`$. Soit $`a`$ une racine d’ordre $`d^{m1}1`$ de l’unité. On pose $`K_a`$ l’ensemble des points $`(x,y,a)`$. Alors $`K_a`$ est périodique de période $`m`$.
###### Lemme 5
Pour tout $`a`$, l’ensemble $`SK_a`$ est fini.
* Preuve— Soit $`V`$ une composante irréductible, périodique de $`SK_a`$. Il faut montrer que $`dimV=0`$. Supposons par l’absurde que $`dimV1`$. Pour simplifier les notations, on suppose par la suite que $`a=1`$ et on pose $`𝒟:=𝒟_{d,d^{}}`$. Notons $`s:=(x,y)=(x_1,\mathrm{},x_{d1},y_1,\mathrm{},y_{d^{}1})`$ les coordonnées de $`K_a^{d+d^{}2}`$ . Comme $`𝒟`$ est polynomial, elle se prolonge en une application méromorphe de $`^{d+d^{}2}`$ dans lui-même. Notons encore $`𝒟`$ ce prolongement. Soit $`L:=^{d+d^{}2}K_a`$ l’hyperplan à l’infini muni des coordonnées homogènes $`w:=[x_1:\mathrm{}:x_{d1}:y_1:\mathrm{}:y_{d^{}1}]`$. On pose $`(f,g):=\mathrm{\Pi }_{d,d^{},a}(s)`$. Les formules explicites des polynômes $`f`$ et $`g`$ sont données dans le paragraphe précédent. On remarque que $`f(y_i)`$ (resp. $`g(x_i)`$) est un polynôme homogène de degré $`d`$ (resp. $`d^{}`$) en variables $`x`$ et $`y`$. Par conséquent, il existe une constante $`c>0`$ telle que $`|f(y_j)|c\lambda ^d`$ et $`|g(x_i)|c\lambda ^d^{}`$
$$\lambda :=\mathrm{max}(\underset{1\nu d1}{\mathrm{max}}|x_\nu |,\underset{1\nu d^{}1}{\mathrm{max}}|y_\nu |).$$
Comme $`d>d^{}`$, l’ensemble d’indétermination $`I`$ de $`𝒟`$ est égal à
$$I=\{wL:f(y_1)=\mathrm{}=f(y_{d^{}1})=0\}$$
et l’ensemble $`X:=𝒟(LI)`$ vérifie
$$X\{wL:x_1=\mathrm{}=x_{d1}=0\}.$$
Comme $`dimV1`$, l’intersection $`\overline{V}L\mathrm{}`$. Comme $`V`$ est prépériodique, $`\overline{V}(IX)\mathrm{}`$. Soient $`s^{(n)}=(x^{(n)},y^{(n)})V`$ tendant vers un point $`w_0\overline{V}(IX)`$ quand $`n+\mathrm{}`$. On pose $`(f_n,g_n):=\mathrm{\Pi }_{d,d^{},a}(s^{(n)})`$, $`\overline{s}^{(n)}:=𝒟(s^{(n)})`$ et $`(\overline{f}_n,\overline{g}_n):=\mathrm{\Pi }_{d,d^{},a}(\overline{s}^{(n)})`$. Par définition de $`𝒟_{d,d^{}}`$, on a $`𝒞_{\overline{f}_n}=g_n(𝒞_{f_n})`$ et $`𝒞_{\overline{g}_n}=f_n(𝒞_{g_n})`$. Par définition de $`𝒩(d,d^{})`$, on a $`\overline{f}_ng_n=\overline{g}_nf_n`$. Soient
$$\lambda _n=\mathrm{max}(\underset{1\nu d1}{\mathrm{max}}|x_\nu ^{(n)}|,\underset{1\nu d^{}1}{\mathrm{max}}|y_\nu ^{(n)}|).$$
Alors $`|\overline{x}_\nu ^{(n)}|c\lambda _n^d^{}`$ et $`|\overline{y}_\nu ^{(n)}|c\lambda _n^d`$. On pose $`\sigma _1(z):=\lambda z`$, $`\sigma _2(z):=\lambda ^{d1}z`$, $`\sigma _3(z):=\lambda ^{d^{}1}z`$ et $`\sigma _4(z):=\lambda ^{(d1)(d^{}1)}z`$. On pose également $`f_n^{}:=\sigma _2^1f_n\sigma _1`$, $`g_n^{}:=\sigma _3^1g_n\sigma _1`$, $`\overline{f}_n^{}:=\sigma _4^1\overline{f}_n\sigma _3`$ et $`\overline{g}_n^{}:=\sigma _4^1\overline{g}_1\sigma _2`$. Alors $`\overline{f}_n^{}g_n^{}=\overline{g}_n^{}f_n^{}`$, $`(f_n^{},g_n^{})\mathrm{\Sigma }(d,d^{},1)`$ et $`(\overline{f}_n^{},\overline{g}_n^{})\mathrm{\Sigma }(d,d^{},1)`$. On a aussi $`𝒞_{f_n^{}}=\sigma _1^1(𝒞_{f_n})=\sigma _1^1\{x_1^{(n)},\mathrm{},x_{d1}^{(n)}\}`$, $`𝒞_{g_n^{}}=\sigma ^1\{y_1^{(n)},\mathrm{},y_{d^{}1}^{(n)}\}`$, $`𝒞_{\overline{f}^{}}=\sigma _3^1\{\overline{x}_1^{(n)},\mathrm{},\overline{x}_{d1}^{(n)}\}`$ et $`𝒞_{\overline{g}^{}}=\sigma _2^1\{\overline{y}_1^{(n)},\mathrm{},\overline{y}_{d^{}1}^{(n)}\}`$. Par définition de $`\lambda _n`$ et des $`\sigma _i`$, les points critiques de $`f^{}`$ et $`g^{}`$ (resp. de $`\overline{f}^{}`$ et $`\overline{g}^{}`$) sont de modules majorés par $`1`$ (resp. par $`c`$). De plus, au moins l’un des points critiques de $`f^{}`$ ou de $`g^{}`$ est de module $`1`$. Le fait que $`(f^{},g^{})\mathrm{\Sigma }(d,d^{},1)`$ et $`(\overline{f}^{},\overline{g}^{})\mathrm{\Sigma }(d,d^{},1)`$ entraîne que les coefficients des polynômes $`f_n^{}`$, $`g_n^{}`$, $`\overline{f}_n^{}`$ et $`\overline{g}_n^{}`$ sont bornés. On vérifie facilement que $`𝒞_{\overline{f}_n^{}}=g_n^{}(𝒞_{f_n^{}})`$ et $`𝒞_{\overline{g}_n^{}}=f_n^{}(𝒞_{g_n^{}})`$. Soient $`F`$, $`G`$, $`\overline{F}`$, $`\overline{G}`$ quatre polynômes tels que $`(F,G,\overline{F},\overline{G})`$ soit adhérent à la suite $`(f_n^{},g_n^{},\overline{f}_n^{},\overline{g}_n^{})`$. Par continuité, on a $`\overline{F}G=\overline{G}F`$, $`𝒞_{\overline{F}}=G(𝒞_F)`$ et $`𝒞_{\overline{G}}=F(𝒞_G)`$. De plus, au moins un point critique de $`F`$ ou de $`G`$ est de module 1.
Cas 1.– Supposons que $`w_0I`$. On déduit de la description de $`I`$ que $`\lambda _n^{d+1}\overline{y}_\nu ^{(n)}`$ tend vers $`0`$ quand $`n+\mathrm{}`$. Ceci implique que les points critiques de $`\overline{G}`$ sont tous nuls. D’où $`\overline{G}(z)=z^d^{}`$ et $`[F(z)]^d^{}=\overline{F}G(z)`$ car $`\overline{F}G=\overline{G}F`$. Alors les multiplicités des zéros de $`\overline{F}G`$ sont divisibles par $`d^{}`$. Comme $`d=\mathrm{deg}\overline{F}`$ n’est pas divisible par $`d^{}`$, il existe au moins une racine $`a_1`$ de $`\overline{F}`$ telle que sa multiplicité $`\alpha _1`$ ne soit pas divisible par $`d^{}`$.
Supposons d’abord qu’il existe une autre racine $`a_2`$ de $`\overline{F}`$ dont la multiplicité $`\alpha _2`$ n’est pas divisible par $`d^{}`$. Soit $`b_j`$ un point arbitraire de $`G^1(a_j)`$ à multiplicité $`\beta _j`$. Alors $`\alpha _j\beta _j`$ est divisible par $`d^{}`$. Notons $`\alpha _j^{}`$ le plus grand diviseur commun de $`\alpha _j`$ et $`d^{}`$. Notons également $`\nu _j=d^{}/\alpha _j^{}`$. Alors $`\nu _j`$ divise $`\beta _j`$. On en déduit qu’il existe un polynôme non constant $`K_j`$ tel que $`G(z)a_j=[K_j(z)]^{\nu _j}`$. Comme $`\alpha _j`$ ne divise pas $`d^{}`$, on a $`\nu _j2`$. On obtient donc $`[K_1(z)]^{\nu _1}=[K_2(z)]^{\nu _2}b^{\nu _2}`$$`b^{\nu _2}=a_2a_10`$. Ceci implique
$$[K_1(z)]^{\nu _1}=\underset{j=0}{\overset{\nu _21}{}}[K_2(z)\theta _jb].$$
$`\theta _j:=\mathrm{exp}(2j\pi i/\nu _2)`$. Les facteurs du membre à droite sont deux à deux premiers entre eux. Par conséquent, il existe des polynômes $`P_j`$ tels que $`K_2(z)\theta _jb=[P_j(z)]^{\nu _1}`$. On a
$$[P_1(z)]^{\nu _1}[P_0(z)]^{\nu _1}=(\theta _1\theta _0)b0.$$
C’est une contradiction car le membre à gauche se factorise en $`\nu _1`$ facteurs qui ne sont pas tous constants.
Il reste le cas où $`a_1`$ est la seule racine de $`\overline{F}`$ dont la multiplicité n’est pas divivible par $`d^{}`$. Comme $`d=\mathrm{deg}\overline{F}`$ et $`d^{}`$ sont premiers entre eux, $`\alpha _1`$ et $`d^{}`$ sont premiers entre eux. Par conséquent, tout point de $`G^1(a_1)`$ est de multiplicité divisible par $`d^{}`$. Comme $`\mathrm{deg}G=d^{}`$ et comme $`G\mathrm{\Sigma }(d,d^{},1)`$, on a $`G(z)=z^d^{}`$. On en déduit que $`a_1=0`$. On peut donc écrire $`\overline{F}(z)=z^{\alpha _1}[P(z)]^d^{}`$$`P`$ est un polynôme unitaire. L’équation $`[F(z)]^d^{}=\overline{F}G(z)`$ entraîne $`F(z)=z^{\alpha _1}P(z^d^{})`$. Les égalités suivantes sont obtenues par les calculs de dérivées:
$$F^{}(z)=z^{\alpha _11}[\alpha _1P(z^d^{})+d^{}z^d^{}P^{}(z^d^{})]$$
et
$$\overline{F}^{}(z)=z^{\alpha _11}[P(z)]^{d^{}1}[\alpha _1P(z)+d^{}zP^{}(z)].$$
Soit $`a`$ une racine non nulle de $`F^{}`$, i.e. une racine de $`\alpha _1P(z^d^{})+d^{}z^d^{}P^{}(z^d^{})`$. Alors $`\mathrm{exp}(2k\pi i/d^{})a`$ est également une racine de $`F^{}`$ pour tout $`0kd^{}1`$. Comme $`𝒞_{\overline{F}}=G(𝒞_F)`$ et comme $`G(z)=z^d^{}`$, toute racine de $`\overline{F}^{}`$ est du type $`a^d^{}`$, i.e. une racine de $`\alpha _1P(z)+d^{}zP^{}(z)`$. De plus, la multiplicité de cette racine est divisible par $`d^{}`$. Soit $`b`$ une racine de multiplicité $`nd^{}+m`$ de $`P(z)`$ avec $`0md^{}1`$. Alors $`b`$ est une racine de multiplicité $`nd^{}+m1`$ de $`P^{}(z)`$ et donc de $`\alpha _1P(z)+d^{}zP^{}(z)`$. Par conséquent, $`b`$ est une racine de multiplicité $`(n^{}d^{}+m)(d^{}1)+(nd^{}+m1)`$ de $`\overline{F}^{}`$. Cette multiplicité n’est pas divisible par $`d^{}`$. C’est impossible. D’où $`P(z)=1`$ et $`F(z)=z^d`$. C’est aussi une contradiction car au moins l’un des points critiques de $`F`$ ou de $`G`$ est de module 1.
Cas 2.– Supposons maintenant que $`w_0X`$. Par la description de $`X`$, $`\lambda _n^1x_\nu ^{(n)}`$ tend vers $`0`$ quand $`n+\mathrm{}`$. Par conséquent, les points critiques de $`f_n^{}`$ tendent vers $`0`$. On en déduit que $`F(z)=z^d`$ et que $`𝒞_{\overline{F}}=G(𝒞_F)=\{0\}`$. On a donc $`\overline{F}(z)=z^d`$. On obtient alors $`\overline{G}(z^d)=[G(z)]^d`$. Ceci montre que les racines de $`\overline{G}(z^d)`$ sont toutes de multiplicité divisible par $`d`$. En particulier, toute racine non nulle de $`\overline{G}`$ est de multiplicité divisible par $`d`$. Mais $`\mathrm{deg}\overline{G}=d^{}<d`$. Donc $`\overline{G}`$ n’a pas de racine non nulle. Alors $`\overline{G}(z)=z^d^{}`$ et donc $`G(z)=z^d^{}`$. C’est une contradiction car au moins un point critique de $`F`$ ou de $`G`$ est de module 1.
$`\mathrm{}`$
Fin de la preuve de la proposition 2.– Si $`dimS=0`$, alors $`S`$ est simplement un point périodique. D’après le lemme 4, $`S𝒞_1(d,d^{})𝒞_2(d,d^{})`$.
Si $`dimS1`$, d’après le lemme précédent, $`SK_a`$ est un ensemble fini pour tout $`a`$. Comme $`S`$ et $`K_a`$ sont périodiques, $`SK_a`$ est périodique. Par conséquent, tout point de $`SK_a`$ est prépériodique. D’après le lemme 4, $`SK_a𝒞_1(d,d^{})𝒞_2(d,d^{})`$. Comme $`S`$ est périodique et comme $`𝒟_{d,d^{}}`$ envoie l’hyperplan $`\{\alpha =c\}`$ dans l’hyperplan $`\{\alpha =c^d\}`$, $`SK_a`$ est non vide sauf peut-être pour un nombre fini de $`a`$. On déduit que $`dimS=1`$ et que $`S`$ coupe $`𝒞_1(d,d^{})𝒞_2(d,d^{})`$ en une infinité de points. D’où $`S𝒞_1(d,d^{})𝒞_2(d,d^{})`$. Ceci est vrai pour tout sous-ensemble algébrique périodique de $`𝒩(d,d^{})`$.
Comme $`𝒟_{d,d^{}}^n(𝒩(d,d^{}))`$ est faiblement invariant pour tout $`n0`$, toute composante de $`𝒩(d,d^{})`$ s’envoie par un $`𝒟_{d,d^{}}^n`$ dans une composante périodique de $`𝒩(d,d^{})`$. Donc elle s’envoie par $`𝒟_{d,d^{}}^n`$ dans $`𝒞_1(d,d^{})𝒞_2(d,d^{})`$. D’après le lemme 3, elle est incluse dans $`𝒞_1(d,d^{})𝒞_2(d,d^{})`$.
## 5 Preuves des théorèmes et remarques
Preuve du théorème 1— Soient $`\mu `$, $`f`$ et $`g`$ vérifiant les hypothèses du théorème 1. Si $`d`$ est divisible par $`d^{}`$ ou si $`d^{}`$ est divisible par $`d`$, d’après le corollaire 3, la condition 1 du théorème 1 est satisfaisante.
Dans le cas contraire, d’après le corollaire 3, on peut supposer que $`d`$ et $`d^{}`$ sont premiers entre eux et que $`d>1`$, $`d^{}>1`$. Sans perdre en généralité, on peut supposer que $`d>d^{}`$. Alors d’après les paragraphes 3 et 4, il existe des polynômes linéaires $`\sigma _1`$, $`\sigma _2`$ et un nombre $`\alpha 0`$ tels que
$$(\sigma _1f\sigma _2,\sigma _1g\sigma _2)\mathrm{\Pi }_{d,d^{},\alpha }(𝒩(d,d^{})).$$
D’après la proposition 3, $`𝒩(d,d^{})=𝒞_1(d,d^{})𝒞_2(d,d^{})`$. Par définition de $`𝒞_1(d,d^{})`$ et $`𝒞_2(d,d^{})`$, il existe des polynômes linéaires $`\sigma _3`$ et $`\sigma _4`$ tels que l’une des conditions suivantes soit vraie:
1. $`\sigma _3f\sigma _4(z)=z^d`$ et $`\sigma _3g\sigma _4(z)=az^d^{}`$$`a0`$ est une constante.
2. $`\sigma _3f\sigma _4=\pm \mathrm{T}_d`$ et $`\sigma _3g\sigma _4=\pm \mathrm{T}_d^{}`$.
Posons $`Q:=\sigma _3^1\sigma _4^1`$, $`f_0:=fQ^1`$ et $`g_0:=gQ^1`$. On a $`f=f_0Q`$ et $`g=g_0Q`$. On a aussi $`f_0=\sigma _3^1z^d\sigma _3`$, $`g_0=\sigma _3^1(az^d^{})\sigma _3`$ ou $`f_0=\sigma _3^1(\pm \mathrm{T}_d)\sigma _3`$, $`g_0=\sigma _3^1(\pm \mathrm{T}_d^{})\sigma _3`$. Alors pour la nouvelle coordonnée $`z^{}:=\sigma _3^1(z)`$, $`f_0`$ et $`g_0`$ vérifient la condition 2 ou la condition 3 du théorème 1.
$`\mathrm{}`$
###### Lemme 6
Soient $`E`$ un compact, $`f`$ et $`g`$ deux polynômes de degrés $`d>1`$ et $`d^{}>1`$. Supposons que $`f^1(E)=g^1(E)`$ et que $`d`$, $`d^{}`$ sont premiers entre eux.
1. Si $`f(z)=az^d`$ et $`g(z)=bz^d^{}`$ avec $`a0`$ et $`b0`$, alors $`E`$ est une réunion de cercles centrés en $`0`$
2. Si $`f(z)=\pm \mathrm{T}_d`$ et $`g(z)=\pm \mathrm{T}_d^{}`$, alors $`E=[1,1]`$.
* Preuve— 1. On note $`S_r`$ le cercle de centre $`0`$ et de rayon $`r0`$. Pour tout compact non vide $`K`$, on pose $`A_K(r)`$ le maximum des longueurs des composantes connexes de $`S_rK`$. On pose
$$A_K:=sup\{A_K(r)\text{ pour tout }r>0\text{ tel que }KS_r\mathrm{}\}.$$
Posons $`F:=f^1(E)`$. On a $`A_F=d^1A_E`$. D’autre part, $`F=g^1(E)`$. D’où $`A_F=d_{}^{}{}_{}{}^{1}A_E`$. On en déduit que $`A_E=A_F=0`$. Par conséquent, $`E`$ est une réunion de cercles centrés en $`0`$.
2. Notons $`\phi `$ la fonction de Green de $`^1[1,1]`$ avec un seul pôle en $`\mathrm{}`$. On a $`d^1\phi f=d_{}^{}{}_{}{}^{1}\phi g=\phi `$. Notons $`E_1=f^1(E)=g^1(E)`$. On a
$$\underset{E_1}{\mathrm{max}}\phi (z)=d^1\underset{E}{\mathrm{max}}\phi (z)=d_{}^{}{}_{}{}^{1}\underset{E}{\mathrm{max}}\phi (z).$$
Par conséquent, $`\phi (z)=0`$ pour $`zE_1`$. D’où $`E_1[1,1]`$ et $`E[1,1]`$. Notons $`\psi (z):=(z+z^1)/2`$. On a $`\psi ^1f\psi (z)=\pm z^{\pm d}`$, $`\psi ^1g\psi (z)=\pm z^{\pm d^{}}`$. Posons $`\stackrel{~}{E}:=\psi ^1(E)`$. Alors $`\stackrel{~}{E}\psi ^1([1,1])=\{z:|z|=1\}`$. On a $`\stackrel{~}{f}^1(\stackrel{~}{E})=\stackrel{~}{g}^1(\stackrel{~}{E})`$$`\stackrel{~}{f}(z):=\pm z^d`$ et $`\stackrel{~}{g}(z):=\pm z^d^{}`$. D’après la partie précédente, $`\stackrel{~}{E}`$ est le cercle unité. D’où $`E=[1,1]`$.
$`\mathrm{}`$
Preuve du corollaire 1— Dans le corollaire 1, la condition nécessaire est évidente. Pour la condition suffisante, supposons par l’absurde qu’il existe deux polynômes distints $`f`$ et $`g`$ tels que $`f^1(E)=g^1(E)`$. Notons $`\mathrm{\Omega }`$ la composante connexe de $`^1E`$ qui contient $`\mathrm{}`$. Alors $`f^1(\mathrm{\Omega })=g^1(\mathrm{\Omega })`$. Comme $`E`$ est de capacité logarithmique positive, il existe une fonction de Green $`\phi `$ de $`\mathrm{\Omega }`$ avec un seul pôle en $`\mathrm{}`$ \[16, III.8\]. Alors $`d^1\phi f`$, $`d_{}^{}{}_{}{}^{1}\phi g`$ sont les fonctions de Green de $`f^1(\mathrm{\Omega })=g^1(\mathrm{\Omega })`$ avec un seul pôle en $`\mathrm{}`$. Comme la fonction de Green est unique, on a $`d^1\phi f=d_{}^{}{}_{}{}^{1}\phi g`$. On pose $`\phi _0(z)=0`$ si $`z\mathrm{\Omega }`$ et $`\phi _0(z)=\phi (z)`$ si $`z\mathrm{\Omega }`$. C’est une fonction subharmonique et $`\mu :=i\overline{}\phi _0`$ est la mesure d’équilibre de $`E`$ \[16, III\]. On obtient par les relations précédentes que $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$. D’après le théorème 1, on a $`f=f_0Q`$ et $`g=g_0Q`$. D’où $`f_0^1(E)=g_0^1(E)`$. Si la condition 1 du théorème 1 est vraie, on a $`f_0^1(E)=g_0^1(E)=E`$.
Si la condition 2 du théorème 1 est vraie, d’après le lemme précédent, $`E`$ est une réunion de cercles centrés en $`0`$. On a $`P^1(E)=E`$ pour toute rotation $`P`$ de centre $`0`$.
Si la condition 3 est vraie, d’après le lemme précédent, $`E=[1,1]`$. Par conséquent, $`P^1(E)=E`$ pour $`P:=\mathrm{T}_k`$.
Dans les trois cas, on obtient une contradiction avec l’hypothèse du corollaire 1.
$`\mathrm{}`$
Preuve du corollaire 2.— Comme $`E`$ est de capacité logarithmique positive, $`E`$ est un ensemble infini. D’après le corollaire 1, il existe un polynôme $`P\mathrm{id}`$ tel que $`P^1(E)=E`$. Si $`P(z)=az+b`$, on a $`|a|=1`$ et $`a1`$ car $`E`$ est compact. Alors $`P`$ est une rotation de centre $`b/(1a)`$. C’est impossible. On a donc $`\mathrm{deg}P2`$. On sait que $`J_P`$ est le plus petit compact totalement invariant par $`P`$ qui contient plus qu’un élément. D’où $`J_PE`$. Comme $`K_P`$ est le plus grand compact totalement invariant par $`P`$, on a $`EK_P`$.
$`\mathrm{}`$
Dans le cas général, si $`E`$ est un compact et si $`f`$, $`g`$ sont deux polynômes vérifiant $`f^1(E)=g^1(E)`$, il n’existe pas de mesure $`\mu `$ à support dans $`E`$ telle que $`d^1f^{}(\mu )=d_{}^{}{}_{}{}^{1}g^{}(\mu )`$. Par exemple pour $`E=\{0\}`$, $`f(z)=z(z1)`$ et $`g(z)=z^2(z1)`$, la seule mesure de probabilité $`\mu `$ supportée par $`E`$ est la masse de Dirac en $`0`$. On a $`f^1(E)=g^1(E)`$ mais $`d^1f^{}(\mu )d_{}^{}{}_{}{}^{1}g^{}(\mu )`$.
###### Proposition 4
Soient $`E`$ un compact, $`f`$ et $`g`$ deux polynômes tels que $`f^1(E)=g^1(E)`$. Alors il existe deux mesures de probabilité $`\mu _1`$ et $`\mu _2`$ à support dans $`E`$ telles que $`g_{}(d^1f^{}(\mu _1))=\mu _1`$ et $`f_{}(d_{}^{}{}_{}{}^{1}g^{}(\mu _2))=\mu _2`$.
* Preuve— Soit $`\delta _0`$ une mesure de probabilité à support dans $`E`$. Posons $`\delta _n:=g_{}(d^1f^{}(\delta _{n1}))`$ pour tout $`n1`$. Alors $`\delta _n`$ est une mesure de probabilité à support dans $`E`$. Posons $`S_n:=(\delta _0+\mathrm{}+\delta _{n1})/n`$. Alors $`S_n`$ est également une mesure de probabilité à support dans $`E`$. Il existe une suite croissante $`\{n_i\}_i`$ telle que $`S_{n_i}`$ tende faiblement vers une mesure $`\mu _1`$ quand $`i+\mathrm{}`$ car l’ensemble des mesures de probabilité à support dans $`E`$ est compact. On en déduit que $`g_{}(d^1f^{}(S_{n_i}))`$ tend faiblement vers $`g_{}(d^1f^{}(\mu _1))`$. D’autre part, $`g_{}(d^1f^{}(S_{n_i}))S_{n_i}=(\delta _{n_i}\delta _0)/n_i`$ tend vers $`0`$. On obtient finalement $`g_{}(d^1f^{}(\mu _1))=\mu _1`$. De même pour $`\mu _2`$.
$`\mathrm{}`$ |
warning/0003/cs0003078.html | ar5iv | text | # About the finding of independent vertices of a graph
## 1 Statement of the problem
Consider the class $`L`$ of undirected graphs without loops and multiple edges with weighted vertices.
Assume that there is a graph $`G=(X,\mathrm{\Gamma },M)L`$, where $`X=\{x_1,\mathrm{},x_n\}`$ be the set of the graph vertices, $`\mathrm{\Gamma }`$ is the mapping $`X`$ into $`X`$, and $`M=\{\mu (x_1),\mathrm{},\mu (x_n)\}`$ is the set of the non-negative integers – weights of the graph vertices. If $`X_1=\{x_{i_1},\mathrm{},x_{i_m}\}X`$ then $`\mathrm{\Gamma }X_1=\mathrm{\Gamma }x_{i_1}\mathrm{}\mathrm{\Gamma }x_{i_m}`$.
A graph $`G=(X,\mathrm{\Gamma },M)L`$ is called isometric if $`\mu (x_i)=\mu (x_j)`$ ($`ij`$) for all $`x_i,x_jX`$.
For any $`AX`$ we shall designate
$$\mu (A)=\underset{x_iA}{}\mu (x_i).$$
As a problem $`Z`$, given on a graph $`G=(X,\mathrm{\Gamma },M)L`$, we shall call the problem of finding of vertex set $`UX`$ such that satisfies conditions
$$U\mathrm{\Gamma }U,$$
(1)
$$U\mathrm{\Gamma }U=X$$
(2)
and supplies the maximum of a function
$$\mu (U).$$
(3)
Any vertex set $`UX`$, satisfying the condition (1), is called independent. An independent set $`UX`$, satisfying (2), is called the maximal independent set (MIS) of the graph $`G`$.
A MIS $`\widehat{U}X`$, supplying the maximum of the function (3), is called the maximum independent set (MMIS) of the graph $`G`$ (the optimum solution of the problem $`Z`$).
The problem $`Z`$ has the different applications . It has the special significance in Computation Complexity Theory, as it is NP-complete . From the point of view of applications, the significance of any NP-complete problem is that it can be considered as a mathematical model of all discrete problems.
The existing methods for solving of the problem $`Z`$ (as a rule, in an isometric graph) consist in basic in finding all MIS of the graph $`GL`$ and selection of them the maximum independent set .
The inefficiency of such approach to solving the problem $`Z`$ is proved by that the maximum number of the MISs $`\sigma _G(n)`$, a graph $`GL`$ can has, is equal to $`\sigma _G(n)=\gamma (s)\dot{3}^{r1}`$ , where $`n=3r+s`$, $`\gamma (0)=3`$, $`\gamma (1)=4`$, $`\gamma (2)=6`$. Hence, complexity of any algorithm, based on searching of all MIS of a graph $`G`$, can not have an evaluation better than $`O(3^{n/3})`$.
With the problem $`Z`$, given on a graph $`G`$, it is usually connected a problem of finding the maximum complete subgraph (the maximal clique) in the additional graph $`\overline{G}=(X,\overline{\mathrm{\Gamma }},M)L`$ as the subset of vertices $`\widehat{U}X`$, inducing the maximum clique of a graph $`\overline{G}`$, is MMIS of a graph $`G`$ .
Notice that the Maximum Clique Problem is a maximize problem, and from the point of view of the approach, accepted in Operations Research, is not dual to the problem $`Z`$.
Unfortunately, the difficulties, connected with finding MMIS of a graph $`GL`$, can not are overcome by development of a polynomial algorithm, enabling to find approximate solution of the problem $`Z`$ with a guaranteed deviation from the optimum solution . Therefore, for development of the solution methods of the problem $`Z`$, it is necessary either to try to create an algorithm, discovering its exact solution (in this case it will be proven that P= NP), or to find the exact solution of the problem for separate subclasses of graphs of $`L`$ (the majority authors go to the last way).
The main result of the given work is that the problem $`Z`$, given on an arbitrary graph $`GL`$, can be considered as the solving the same problem in subclass of the normal conjugate-orthogonal graphs. A problem is formulated that is dual to the problem $`Z`$. It is shown that, for trivial conjugate-orthogonal graphs, any of its MIS is also a MMIS.
## 2 A normal graph
Divide a set of all vertices of a graph $`G=(X,\mathrm{\Gamma },)L`$ into classes <sub>j</sub> ($`j=\overline{1,s}`$) such that if $`_{i_1}`$, $`{}_{i_2}{}^{}K_j`$, then $`\mathrm{\Gamma }_{i_1}=\mathrm{\Gamma }_{i_2}`$. The set of all such classes of the graph is designated by $`H_G`$.
###### Theorem 2.1
If $`{}_{i_1}{}^{}U`$ and $`{}_{i_1}{}^{}_{j}^{}`$ then $`{}_{j}{}^{}U`$, where $`U`$ be a MIS of a graph $`G=(X,\mathrm{\Gamma },)`$.
Assume the conditions of Theorem 2.1 are satisfied and we allow that there exists a vertex $`{}_{i_2}{}^{}_{j}^{}`$ such that $`{}_{i_2}{}^{}U`$ ($`i_1i_2`$).
As $`\mathrm{\Gamma }x_{i_2}=\mathrm{\Gamma }x_{i_1}`$ then $`{}_{i_2}{}^{}\mathrm{\Gamma }U`$ owing to (1). But, it takes into considering (2), we have $`{}_{i_2}{}^{}U`$. The contradiction have obtained. Q.E.D.
###### Theorem 2.2
For any vertex $`{}_{i}{}^{}`$ ($`i=\overline{1,n}`$) of a graph $`G=(X,\mathrm{\Gamma },)`$
$$\mathrm{\Gamma }x_i=\underset{r}{}K_{j_r}(K_{j_r}H_G)$$
It is clear that the vertex set $`\mathrm{\Gamma }x_i`$ can be divided into the classes $`K_{j_1}^{}`$, …, $`{}_{}{}^{}{}_{j_t}{}^{}`$ as it is mentioned above. Assume that these classes of vertices are distinct from similar vertex classes of the graph $`G`$, that is, $`{}_{j_r}{}^{}_{j_r}^{}`$ and $`{}_{j_r}{}^{}_{j_r}^{}`$ ($`=\overline{1,t}`$).
It follows from here that there are the vertices $`{}_{k_1}{}^{}_{j_r}^{}K_{j_r}`$ and $`x_{k_2}_{j_r}_{j_r}^{}`$ such that $`{}_{k_1}{}^{}\mathrm{\Gamma }_i\mathrm{\Gamma }x_i`$ and $`{}_{k_2}{}^{}\mathrm{\Gamma }_i`$.
As vertices $`_{k_1}`$, $`{}_{k_2}{}^{}_{j_r}^{}`$ then $`\mathrm{\Gamma }_{k_1}=\mathrm{\Gamma }_{k_2}`$ by the definition. If $`{}_{k_1}{}^{}\mathrm{\Gamma }x_i`$ then $`{}_{i}{}^{}\mathrm{\Gamma }_{k_1}`$, it signifies, $`{}_{i}{}^{}\mathrm{\Gamma }_{k_2}`$. Then we have $`x_{k_2}\mathrm{\Gamma }x_i`$. The contradiction have obtained. Q.E.D.
Thus, it is established that $`\mathrm{\Gamma }x_i`$, for any vertex $`{}_{i}{}^{}`$ ($`i=\overline{1,}`$) of a graph $`G=(X,\mathrm{\Gamma },)`$, be an union of some classes $`{}_{j_r}{}^{}H_G`$.
###### Corollary 2.1
For any class $`{}_{j}{}^{}H_G`$ of a graph $`G=(X,\mathrm{\Gamma },)`$
$$\mathrm{\Gamma }K_j=\underset{r}{}K_{j_r}(K_{j_r}H_G)$$
A graph $`G_1=(X_1,\mathrm{\Gamma }_1,_1)L`$ is called normal if for any two vertices $`y_{j_1}`$, $`{}_{j_2}{}^{}X_1`$ ($`j_1j_2`$) the relation $`\mathrm{\Gamma }y_{j_1}\mathrm{\Gamma }_{j_2}`$ takes place.
Obviously, that for any graph $`G=(X,\mathrm{\Gamma },M)L`$ can be found a mapping $`\varphi `$: $`G_1=\varphi (G)`$, where $`G_1=(X_1,\mathrm{\Gamma }_1,M_1)L`$ be the normal graph. Thus, $`G_1=\varphi (G)`$ if $`y_j=\varphi (K_j)`$ ($`K_jH_G`$) and $`\mathrm{\Gamma }_1y_j=\varphi (\mathrm{\Gamma }K_j)`$, $`\mu (y_j)=\mu (K_j)`$ for all $`y_jX_1`$ ($`j=\overline{1,s}`$).
Fig. 1 (a) shows the graph $`GL`$ with the unit weights of its vertices, and Fig. 1 (B) shows the normal graph $`G_1`$ that corresponds it (weights of its vertices are put in brackets).
Designate by $`L_H`$ the set of all normal graphs with the weighted vertices that correspond graphs of the class $`L`$.
Further, speaking about a graph $`G=(X,\mathrm{\Gamma },)`$, we shall mean that $`GL_H`$. Besides, we assume that $`CardX=n`$.
## 3 A twin-orthogonal graph
Let $`G=(X,\mathrm{\Gamma },)L_H`$.
The adjacent vertices <sub>1</sub>, $`x_2`$ of the graph $`G`$ is called orthogonal if for all vertices $`{}_{i_1}{}^{}\mathrm{\Gamma }x_1\{x_2\}`$ and $`x_{i_2}\mathrm{\Gamma }_2\{x_1\}`$, when they exist, the relation are fulfilled:
$$\mathrm{\Gamma }x_1\mathrm{\Gamma }x_{i_2},\mathrm{\Gamma }x_2\mathrm{\Gamma }x_{i_1}.$$
(4)
###### Theorem 3.1
If at least one of adjacent vertices <sub>1</sub>, $`x_2X`$ of a graph $`G`$ is dangling then the vertices $`x_1`$ and $`x_2`$ are orthogonal.
Really, suppose, for example, a vertex $`x_1`$ of a graph $`G`$, adjacent with a vertex $`{}_{2}{}^{}X`$, is dangling. Hence, $`\mathrm{\Gamma }_1=\{_2\}`$.
Then we shall have $`\mathrm{\Gamma }_1\{x_2\}`$ and $`\mathrm{\Gamma }_1\mathrm{\Gamma }x_{i_2}`$ for all $`x_{i_2}\mathrm{\Gamma }_2\{_1\}`$ when $`\mathrm{\Gamma }x_2\{x_1\}`$ (as $`x_2\mathrm{\Gamma }x_{i_2}`$). Q.E.D.
###### Theorem 3.2
Let $`U`$ be an arbitrary MIS of a graph $`G=(X,\mathrm{\Gamma },M)`$. If $`x_1`$, $`x_2`$ is the orthogonal vertices of $`G`$ either $`{}_{1}{}^{}U`$ or $`{}_{2}{}^{}U`$.
Assume that the conditions of Theorem 3.2 are satisfied, and suppose that $`x_1`$, $`x_2\mathrm{\Gamma }U`$. Then there exists at least vertex $`{}_{3}{}^{}\mathrm{\Gamma }x_1`$ such that $`x_3U`$, and at least vertex $`{}_{4}{}^{}\mathrm{\Gamma }x_2`$ such that $`x_4U`$.
By the condition (4), for the orthogonal vertices $`x_1`$, $`x_2X`$, we shall have $`x_3\mathrm{\Gamma }_4`$ and $`x_4\mathrm{\Gamma }x_3`$, that is, the vertices $`x_3`$, $`x_4U`$ are adjacent. We have received the contradiction. Q.E.D.
A graph $`\stackrel{~}{G}=(\stackrel{~}{X},\stackrel{~}{\mathrm{\Gamma }},\stackrel{~}{M})L_H`$ is called twin-orthogonal if graph vertices can divide into pairs of the orthogonal vertices. It is clear that $`rd()=n=2k`$, where $`k`$ is a non-negative integer.
Fig. 2 is represented of a twin-orthogonal graph with the unit weights of the vertices.
We shall be to say that the twin-orthogonal graph $`\stackrel{~}{G}=(\stackrel{~}{X},\stackrel{~}{\mathrm{\Gamma }},\stackrel{~}{})`$ corresponds a graph $`G=(X,\mathrm{\Gamma },)`$, if:
* $`X\stackrel{~}{X}`$;
* $`\mu (x_i)=0`$ for any vertex $`x_i\stackrel{~}{X}X`$;
* any MIS $`UX`$ of $`G`$ can be obtained from some MIS $`\stackrel{~}{U}\stackrel{~}{X}`$ of $`\stackrel{~}{G}`$ by removal of all vertices $`x_i\stackrel{~}{X}`$ such that $`x_iX`$.
It is easy to see that one of twin-orthogonal graphs $`\stackrel{~}{G}=(\stackrel{~}{X},\stackrel{~}{\mathrm{\Gamma }},\stackrel{~}{})`$, corresponding a graph $`G=(X,\mathrm{\Gamma },M)`$, can be constructed as follows.
Let $`X_1`$ be the set of all vertices of the graph $`G`$, not being orthogonal for one vertex of this graph. We join a set of vertices $`{}_{2}{}^{}=\{x_{n+1},\mathrm{},x_{n+p}\}`$ ($`p=rd(X_1)`$) to the graph $`G`$, and each of vertices $`{}_{k}{}^{}_{2}^{}`$ we connect by an edge with one and only one of vertices $`{}_{j}{}^{}X_1`$. We assume that $`\mu (x_k)=0`$ for all $`{}_{k}{}^{}X_2`$.
It is clear that, as a result, a twin-orthogonal graph $`\stackrel{~}{G}=(\stackrel{~}{X},\stackrel{~}{\mathrm{\Gamma }},\stackrel{~}{})`$ will be obtained that is induced on a vertex set $`\stackrel{~}{}=_2`$, where $`\stackrel{~}{\mathrm{\Gamma }}x_i=\mathrm{\Gamma }_i`$ for any vertex $`{}_{i}{}^{}_1`$, $`\stackrel{~}{\mathrm{\Gamma }}x_j=\mathrm{\Gamma }x_j\{_k\}`$ for all $`x_j_1`$ and $`\stackrel{~}{\mathrm{\Gamma }}x_h=\{x_j\}`$ for all $`x_kX_2`$.
It is easy to be convinced that the constructed twin-orthogonal graph $`\stackrel{~}{G}`$ corresponds the initial graph $`G`$.
More simple way for a construction of the twin-orthogonal graph $`\stackrel{~}{G}`$ = $`(\stackrel{~}{X}`$, $`\stackrel{~}{\mathrm{\Gamma }}`$, $`\stackrel{~}{})`$, corresponding a graph $`G`$ = $`(X,\mathrm{\Gamma },M)`$, is based on Theorem 3.1.
We shall join a vertex set <sub>1</sub> ($`rd(X_1)=rd(X)`$) to a graph $`G`$ such that each vertex $`{}_{j}{}^{}_{1}^{}`$ we shall connect by an edge with one and only one of vertices $`{}_{i}{}^{}X`$. We assume $`\mu (_j)=0`$ for all $`{}_{j}{}^{}X_1`$. As a result, obviously, it be also obtained a twin-orthogonal graph $`\stackrel{~}{G}=(\stackrel{~}{X},\stackrel{~}{\mathrm{\Gamma }},\stackrel{~}{})`$ induced on a set of vertices $`\stackrel{~}{}=X_1`$, where $`\stackrel{~}{\mathrm{\Gamma }}_i=\mathrm{\Gamma }_i\{_j\}`$ for any vertex $`{}_{i}{}^{}`$ and $`\stackrel{~}{\mathrm{\Gamma }}x_j=\{_i\}`$ for all $`{}_{j}{}^{}X_1`$.
###### Theorem 3.3
If $`\stackrel{~}{U}\stackrel{~}{}`$ be an optimum solution of the problem $`Z`$ on the twin-orthogonal graph $`\stackrel{~}{\mathrm{\Gamma }}_i=\mathrm{\Gamma }_i\{_j\}`$ corresponding a graph $`G=(X,\mathrm{\Gamma },)`$ then an optimum solution $`\widehat{U}`$ of the problem $`Z`$, given on the graph $`G`$, can be obtained by removal from $`\stackrel{~}{U}`$ of all vertices $`{}_{i}{}^{}\stackrel{~}{U}`$ such that $`{}_{i}{}^{}\stackrel{~}{}`$, and, besides, $`\mu (\widehat{U})=\mu (\stackrel{~}{U})`$.
It follows from the definition of a twin-orthogonal graph $`\stackrel{~}{G}`$, corresponding a graph $`G`$. Q.E.D.
## 4 Some properties of a twin-orthogonal graph
Let $`L_0`$ be the set of the normal twin-orthogonal graphs.
###### Theorem 4.1
If $`U_1`$, $`U_2`$ be the different MISs of a twin-orthogonal graph $`G=(X,\mathrm{\Gamma },)L_0`$ then $`rd(U_1)=rd(U_2)=k`$, where $`rd()=n=2k`$, $`k`$ be a non-negative integer.
It follows from Theorem 3.2. Q.E.D.
A twin-orthogonal graph $`G=(X,\mathrm{\Gamma },)`$ is called trivial if for any orthogonal vertices <sub>i</sub>, $`{}_{j}{}^{}X`$ the relation is fulfilled: $`\mu (_i)=\mu (x_j)`$.
###### Theorem 4.2
If $`G=(X,\mathrm{\Gamma },)`$ be a trivial twin-orthogonal graph then any MIS is also MMIS.
It follows from Theorems 3.2 and 4.1. Q.E.D.
###### Theorem 4.3
If $`x_1`$, $`x_2`$ be the orthogonal vertices of a graph $`G=(,\mathrm{\Gamma },)`$ then any pair of vertices from a set $`\mathrm{\Gamma }x_1\{x_2\}`$ ($`\mathrm{\Gamma }x_2\{x_1\}`$) is not orthogonal.
Assume that the conditions of Theorem 4.3 are satisfied, and we suppose that the vertices $`_{i_1}`$, $`{}_{i_2}{}^{}\mathrm{\Gamma }x_1\{_2\}`$ are orthogonal. Then we have $`{}_{1}{}^{}\mathrm{\Gamma }x_{i_1}`$, and $`{}_{1}{}^{}\mathrm{\Gamma }_{i_2}`$, that is, the relation (4) are not fulfilled for vertices $`x_{i_1}`$, $`x_{i_2}`$. We have obtained the contradiction. Q.E.D.
###### Corollary 4.1
If the vertices <sub>1</sub>, $`x_2`$ of a graph $`G=(X,\mathrm{\Gamma },)`$ are orthogonal then they do not form a three-vertex clique with any vertex $`{}_{i}{}^{}`$ ($`i1`$, $`i2`$).
###### Corollary 4.2
If the vertices $`{}_{1}{}^{},_{2}^{}`$ of a graph $`G=(X,\mathrm{\Gamma },)`$ are orthogonal then $`(\mathrm{\Gamma }x_1\{x_2\})(\mathrm{\Gamma }x_2\{x_1\})=`$.
## 5 A dual problem
Further, for convenience, any two orthogonal vertices of a graph $`G`$ = $`(X`$, $`\mathrm{\Gamma },)`$ $`L_0`$ we shall designate by <sub>i</sub> and $`{}_{}{}^{}{}_{i}{}^{}`$.
A graph $`G^{}=(X,\mathrm{\Gamma }^{},M)L_0`$, obtained from a graph $`G=(X,\mathrm{\Gamma },)L_0`$ by renaming of pairs of orthogonal vertices, is called conjugate for the graph $`G`$.
Thus, any orthogonal vertices <sub>i</sub>, $`{}_{i}{}^{}`$ are adjacent in graphs $`G`$ and $`G^{}`$. The vertices <sub>i</sub>, $`{}_{j}{}^{}X`$, if they are not orthogonal in the graph $`G`$, are adjacent in the graph $`G^{}`$ if and only if corresponding vertices $`x_i^{}`$, $`x_j^{}X`$ are adjacent in the graph $`G`$.
Obviously, that $`(G^{})^{}=G`$.
A problem of finding of a vertex set $`U`$ of a graph $`G^{}=(X,\mathrm{\Gamma }^{},)`$, satisfying conditions (1), (2) and supplying the minimum of the function (3), we shall call dual to the problem $`Z`$.
A MIS $`\stackrel{ˇ}{U}X`$, supplying the minimum of the function (3), is called the minimum independent set of vertices (MNMIS) of a graph $`G^{}`$.
The following statements are proved easily.
###### Theorem 5.1
If $`U`$ be a MIS of a graph $`G=(X,\mathrm{\Gamma },M)`$ then $`\mathrm{\Gamma }U=U`$ be a MIS of the conjugate graph $`G=(X,\mathrm{\Gamma }^{},)`$.
###### Theorem 5.2
If $`U_1`$, $`U_2`$ be MISs of a graph $`G=(X,\mathrm{\Gamma },)`$ then $`\mu (U_1)\mu (U_2)`$ if and only if $`\mu (\mathrm{\Gamma }U_1)\mu (\mathrm{\Gamma }U_2)`$.
The following theorem is a corollary of Theorems 5.1 and 5.2.
###### Theorem 5.3
MIS $`\widehat{U}`$ is MNMIS of a graph $`G=(X,\mathrm{\Gamma },)`$ if and only if $`\mathrm{\Gamma }\widehat{U}=\widehat{U}`$ is MNMIS of a conjugate graph $`G^{}=(X,\mathrm{\Gamma }^{},)`$.
###### Theorem 5.4
Let $`\widehat{U}`$, $`\stackrel{ˇ}{U}`$ be MMIS and MNMIS of a graph $`G`$ = $`(X`$, $`\mathrm{\Gamma }`$, $`M)`$ respectively. Then a relation takes place
$$0\mu (\widehat{U})\mu (\stackrel{ˇ}{U})\underset{x_i,x_i^{}X}{}|\mu (x_i)\mu (x_i^{})|.$$
It is also easy to be convinced in a validity of this statement. |
warning/0003/math-ph0003041.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Clifford algebras (CA) are very important in theoretical and mathematical physics. It is almost impossible to list all its applications but the interested reader can found some of them in references . Now, in order to define the CA of a given vector space $`V`$ one need to endow $`V`$ with a (in general symmetric) bilinear form $`g_{p,q}`$. One interesting fact is that the structure of the CA depends not only on the dimension of $`V`$ but also on the signature of $`g`$. In another words, real Clifford algebras $`Cl_{p,q}`$ and $`Cl_{p^{},q^{}}`$ ($`p+q=p^{}+q^{}=n`$) are in general not isomorphic, that is, when we change the signature we get in general different Clifford algebras. For example, the Clifford algebras $`Cl_{1,3}`$ and $`Cl_{3,1}`$ are not isomorphic – indeed, in terms of matrix algebras, the former is isomorphic to the algebra of $`2\times 2`$ quaternionic matrices while the latter is isomorphic to the algebra of $`4\times 4`$ real matrices.
Changing the signature of a given space may sound in principle an artificial process but undoubtly it is a very important thing in modern physics. For example, the euclidean formulation of field theories is a fundamental tool in modern physics . Indeed, sometimes it seems to be even crucial, as in the theory of instantons, in finite temperature field theory and in lattice gauge theory. Going from an euclidean to a minkowskian theory or vice-versa involves changing the signature of the metric over the spacetime and in general a minkowskian theory is transformed in an euclidean theory by analytical continuation, that is, by making $`tit`$. The interpretation of making $`tit`$ is not trivial – in relation to this, see (and references therein) where it was interpreted as a rotation in a five dimensional spacetime.
Despite the ingenuity of an approach like in interpreting $`tit`$ as a rotation in a five dimensional space, we believe it is unsatisfactory since the use of an additional time coordinate in spacetime appears to us to be meaningless. Indeed it would be much more satisfactory if one could find an approach to describe the signature change where there is no need to introduce extra dimensions. In the case where the problem involves CA the situation is even more problematic since the signature changed CA and the original one may not be isomorphic. A possible way to overcome this situation is to complexify the original real CA since the structure of complex CA depends only on the dimension of $`V`$, but this approach can also be seen as a result of introducing an extra dimension (see below).
The objective of this paper is to introduce an algebraic approach to the change of signature in CA where there is no need to introduce any extra dimension to describe it. The signature change appears as a transformation on the algebraic structure underlying the theory. The idea is to propose an operation that “simulates” the product properties of the signature changed space in terms of the original space or vice-versa. The particular case we have in mind is the one involving the four dimensional spacetime, where we introduce an operation “simulating” the properties of minkowski spacetime in terms of an euclidean spacetime and vice-versa. This operation will be called “vee product” (since it will be denoted by a $``$ in order to distinguish it from the usual Clifford product that will be denoted by juxtaposition). The advantage of this approach is obvious since we can retain the “physics” (the minkowskian properties) in a suitable mathematical world (the euclidean spacetime). The fact that we can define the vee product in terms of the Clifford product means that we can describe the minkowskian properties in terms of euclidean spacetime and vice-versa.
Our approach has been inspired by the work of Lounesto . But there are some differences in the method and interpretation which the reader can compare. Moreover, Lounesto only discussed some problems involving the case of signature change corresponding to opposite signatures, that is, $`(p,q)`$ and $`(q,p)`$. Cases like $`(p,q)`$ and $`(p+q,0)`$ have not been considered by Lounesto and this is the case we have when considering minkowskian and euclidean spacetimes. Some applications of this are discussed here – and some others can be found in .
As an application, we are first interested in this paper in studying the Dirac equation. First of all, the version of Dirac equation we obtain by making $`tit`$ – we shall call it the euclidean Dirac equation – has physical properties that are obviously different from the original Dirac equation – which we shall call the minkowskian Dirac equation. The question we want to address in this context is if it is possible to obtain a new equation that exhibit the same physical properties of the original equation. More specifically, our idea is to write the minkowskian Dirac equation in the euclidean spacetime, which should obviously be different from an euclidean Dirac equation in an euclidean spacetime. Minkowskian and euclidean spacetimes are different worlds, both mathematically and physically speaking, and we want to simulate the minkowskian scenario in an euclidean world. The idea is to write an equation in euclidean spacetime in terms of the vee product such that it is equivalent to the original equation in terms of the Clifford product in Minkowski spacetime. Of course the Dirac equation in terms of the vee product is expected to be different from the euclidean Dirac equation.
In order to introduce our approach we need first of all to take some care. Dirac equation is in general formulated in terms of $`4\times 4`$ complex matrices - the gamma matrices - obeying the relation $`\mathrm{\Gamma }_\mu \mathrm{\Gamma }_\nu +\mathrm{\Gamma }_\nu \mathrm{\Gamma }_\mu =2g_{\mu \nu }`$, where $`g_{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ in the euclidean case and $`g_{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ in the minkowskian case. This algebra is the matrix representation of the complex Clifford algebra $`Cl_{}(4)`$. On the other hand, this algebra is the complexification of the real algebras $`Cl_{1,3}`$ and $`Cl_{4,0}`$ associated with the minkowskian and euclidean spacetimes, respectively, that is, $`Cl_{}(4)Cl_{1,3}Cl_{4,0}`$. Moreover, $`Cl_{}(4)`$ is also isomorphic to the real algebra $`Cl_{4,1}`$, while the real algebras $`Cl_{1,3}`$ and $`Cl_{4,0}`$ \- which are also isomorphic<sup>1</sup><sup>1</sup>1However $`Cl_{3,1}`$ and $`Cl_{4,0}`$ are not isomorphic, and we shall see how to consider this case also. Anyway, it is not the isomorphism between $`Cl_{1,3}`$ and $`Cl_{4,0}`$ that matters in this discussion. \- are isomorphic to the even subalgebra of $`Cl_{4,1}`$ . All these facts show that when we complexify the real algebra $`Cl_{1,3}`$ getting $`Cl_{1,3}Cl_{}(4)Cl_{4,1}`$ we are introducing an extra dimension to the minkowskian spacetime such that the extra dimension is of the type of euclidean time. In the same way, when we complexify the real algebra $`Cl_{4,0}`$ getting $`Cl_{4,0}Cl_{}(4)Cl_{4,1}`$ we are introducing an extra dimension to the euclidean spacetime such that the extra dimension is of the type of minkowskian time.
Now, our idea is to not use any extra dimension, and one certain way to do this is to avoid using any of these complexified structures. This can be achieved using the real formulation of Dirac theory due to Hestenes . One can easily formulate the Dirac theory in terms of the Clifford algebra $`Cl_{1,3}`$, which is isomorphic to the algebra of $`2\times 2`$ quaternionic matrices. A Dirac spinor in this way is represented by a pair of quaternions, but rather we prefer to use the form of Dirac equation called Dirac-Hestenes equation in terms of the so called Dirac-Hestenes spinor . The reason for our choice is simple. The Dirac-Hestenes spinor is represented by an element of the even subalgebra $`Cl_{1,3}^+`$ while a Dirac spinor is an element of an ideal of $`Cl_{1,3}`$. Both formulations are equivalent but if we use the former one we avoid the problem of considering the transformation between different ideals (in fact left and right ideals), which will happen in the latter case, and in order to not do unnecessary work – and even hidden some fundamental facts – we prefer to use the Dirac-Hestenes equation.
As another application, we discuss the problem of finding self-dual and anti-self-dual solutions of gauge fields. Since the group being abelian or not is irrelevant for this matter we shall restrict our attention to the abelian case. As is well-known, in the Minkowski spacetime there does not exist (real) solutions to the problem $`F=\pm F`$, where $`F`$ is the 2-form representing the electromagnetic field and $``$ is the Hodge star operator. However, in an euclidean spacetime this problem has solutions and they are given by $`E=\pm B`$, where $`E`$ and $`B`$ are the electric and magnetic components of $`F`$. Now, using the operation we discussed above, we can define a new Hodge-like operator on Minkowski spacetime such that in relation to this new operator we have self-dual and anti-self-dual solutions for that problem on Minkowski spacetime. Moreover, using this Hodge-like operator we can completely simulate an euclidean metric while still working on Minkowski spacetime. We also show that the relation between those Hodge-like operators is given by the parity operation.
We organized this paper as follows. In section 2 we briefly discuss the Clifford algebras. In section 3 we discuss the transformation from euclidean to minkowskian spacetimes and vice-versa from an algebraic point of view. Our idea is to define a new Clifford product in terms of the original Clifford product that defines the algebra we are working. This new product simulates the algebra of Minkowski spacetime inside the algebra of the euclidean spacetime. Then in section 4 we discuss the minkowskian Dirac equation over the euclidean spacetime using our approach. Of course that one could be interested in the other way, that is, to simulate an euclidean world inside a minkowskian spacetime. Indeed we can consider any other possibility, as we discuss in section 5. Finally in section 6 we discuss how to obtain solutions corresponding to self-dual and anti-self-dual gauge fields and the definition of a Hodge-like operation with many interesting properties.
## 2 Mathematical Preliminaries
There are many different ways to define Clifford algebras , each of them emphasizing different aspects. Our approach has been choosen due to the direct introduction to Clifford product which will be fundamental in the next sections.
Let $`\{𝐞_1,\mathrm{},𝐞_n\}`$ be an orthonormal basis for $`^{p,q}`$, where $`^{p,q}`$ is a real vector space of dimension $`n=p+q`$ endowed with an interior product $`g:^{p,q}\times ^{p,q}`$. Writing the quadratic form $`_{i=1}^n_{j=1}^ng_{ij}x_ix_j`$ as the square of the linear expression $`_{i=1}^nx_i𝐞_i`$ and assuming the distributive property we obtain the well-known expression for the Clifford algebra $`Cl(^{p,q},g)Cl_{p,q}`$,
$$𝐞_i𝐞_j+𝐞_j𝐞_i=2g_{ij}$$
(1)
where $`g_{ij}`$ are the metric components. This defines the Clifford product, which has been denoted by juxtaposition.
There is a product $``$, called the exterior product, underlying the Clifford algebra. It is an associative, bilinear and skew-symmetric product of vectors. Furthermore, by applying it to our orthogonal basis we can construct a new vector space $`\mathrm{\Lambda }^2(^{p,q})`$ whose elements are called bivectors, i.e., $`:^{p,q}\times ^{p,q}\mathrm{\Lambda }^2(^{p,q})`$. The skew-symmetric property allows us to extend the definition to $`\mathrm{\Lambda }^n(^{p,q})`$. In general $`:\mathrm{\Lambda }^k(^{p,q})\times \mathrm{\Lambda }^l(^{p,q})\mathrm{\Lambda }^{k+l}(^{p,q})`$.
If $`\{𝐞_1,\mathrm{},𝐞_n\}`$ is a basis of $`^{p,q}`$ then $`1`$ and the Clifford products $`𝐞_{i_1}\mathrm{}𝐞_{i_k}`$, $`(1i_1<i_2<\mathrm{}<i_kn)`$ will establish a basis for $`Cl_{p,q}`$ which has dimension $`2^n`$. If $`\{𝐞_1,\mathrm{},𝐞_n\}`$ is an orthogonal basis then $`𝐞_1\mathrm{}𝐞_n=𝐞_1\mathrm{}𝐞_n`$, which is usually called volume element. It follows that $`Cl_{p,q}`$ and $`\mathrm{\Lambda }(^{p,q})=_{k=0}^n\mathrm{\Lambda }^k(^{p,q})`$ are isomorphic as vector spaces. Therefore, a general element $`ACl_{p,q}`$ takes the form
$$A=A_0+A_1+\mathrm{}+A_n$$
(2)
where $`A_r`$, called an r-vector, belongs to $`\mathrm{\Lambda }^r(^{p,q})Cl_{p,q}`$, $`(r=0,1,\mathrm{},n)`$. More explicitly,
$$A=a_0+a^i𝐞_i+a^{ij}𝐞_{ij}+\mathrm{}+a^{1\mathrm{}n}𝐞_{1\mathrm{}n}$$
(3)
It is convenient to define a projector $`_r`$ as $`_r:\mathrm{\Lambda }(V)\mathrm{\Lambda }^r(V)`$, i.e., $`A_r=A_r`$.
An important property is that Clifford algebra is a $`Z_2`$-graded algebra, i.e., we can divide it into even ($`Cl_{p,q}^+`$) and odd ($`Cl_{p,q}^{}`$) grades. $`Cl_{p,q}^+`$ is a sub-algebra of $`Cl_{p,q}`$ called even sub-algebra. In our example, $`Cl_{3,0}^+=\{1,𝐞_{12},𝐞_{13},𝐞_{23}\}`$ and $`Cl_{3,0}^{}=\{𝐞_1,𝐞_2,𝐞_3,𝐞_{123}\}`$. Some important identities that we will use later are ($`a^{p,q}`$, $`B\mathrm{\Lambda }^r(^{p,q})Cl_{p,q}`$):
$$aB=aB+aB$$
(4)
$$aBaB_{r+1}=\frac{1}{2}(aB+(1)^rBa)$$
(5)
$$aBaB_{r1}=\frac{1}{2}(aB(1)^rBa)$$
(6)
## 3 The Vee Product
Let $`V`$ be a vector space of dimension $`n=4`$. We have five different Clifford algebras depending on the signature: $`Cl_{4,0}`$, $`Cl_{3,1}`$, $`Cl_{2,2}`$, $`Cl_{1,3}`$ and $`Cl_{0,4}`$. With the (possible) exception of $`Cl_{2,2}`$ the importance of the others in modern physics is more than obvious.
First we shall consider the case involving the algebras $`Cl_{1,3}`$ and $`Cl_{4,0}`$. Let $`A,BCl_{4,0}`$ and $`AB`$ be its Clifford product. Now we define a new product, which we call a vee product, $`AB`$ simulating the $`Cl_{1,3}`$ Clifford product in $`Cl_{4,0}`$. After the selection in $`^{4,0}`$ of an arbitrary unit vector $`𝐞_0`$ to represent the fourth dimension and the completion of the basis with three other orthonormal vectors $`𝐞_i`$, we define for $`𝐮,𝐯^{4,0}`$
$$𝐮𝐯:=(1)[\mathrm{𝐯𝐮}2(𝐯𝐞_0)(𝐞_0𝐮)]$$
(7)
Using this product and the $`Cl_{4,0}`$ standard basis it is easy to prove that $`𝐞_0𝐞_0=1`$ and $`𝐞_i𝐞_i=1`$ ($`i=1,2,3`$), while $`𝐞_\mu ^2=𝐞_\mu 𝐞_\mu =1`$ ($`\mu =0,1,2,3`$). Moreover, for $`𝐮,𝐰^{4,0}`$,
$$\mathrm{𝐮𝐰}+\mathrm{𝐰𝐮}=2𝐮𝐰=2(u_0w_0+u_1w_1+u_2w_2+u_3w_3)$$
(8)
but if we use the vee product then
$`𝐮𝐰+𝐰𝐮`$ $`=`$ $`\mathrm{𝐰𝐮}+2(𝐰𝐞_0)(𝐞_0𝐮)\mathrm{𝐮𝐰}+2(𝐮𝐞_0)(𝐞_0𝐰)`$
$`=`$ $`2(u_0w_0u_1w_1u_2w_2u_3w_3)`$
A little bit more general case is when one has a vector and a $`k`$-graded element, i.e., $`𝐯`$ and $`B_k`$
$$B_k𝐯=(1)^k[𝐯B_k2(𝐯𝐞_0)(𝐞_0B_k)]$$
(9)
$$𝐯B_k=(1)^k[B_k𝐯2(B_k𝐞_0)(𝐞_0𝐯)]$$
(10)
Now with the help of (6) one can see that $``$ is associative, that is, $`𝐯(𝐮𝐰)=(𝐯𝐮)𝐰`$. Moreover, the vee product preserves the multivectorial structure since
$$\frac{1}{2}[𝐮,𝐯]_{}=\frac{1}{2}(𝐮𝐯𝐯𝐮)=\frac{1}{2}(\mathrm{𝐯𝐮}+2u_0v_0+\mathrm{𝐮𝐯}2u_0v_0)=\frac{1}{2}(\mathrm{𝐮𝐯}\mathrm{𝐯𝐮})=\frac{1}{2}[𝐮,𝐯]=𝐮𝐯$$
Finally, we can generalize those expression as
$$A_lB_k=(1)^{kl}[B_kA_l2(B_k𝐞_0)(𝐞_0A_l)]$$
(11)
## 4 Dirac equation and vee product
In this section we come to consider the results recently introduced and their applications to Dirac equation. We shall use Dirac-Hestenes equation as discussed in the introduction. As starting point for our analysis we take the Dirac-Hestenes equation in $`Cl_{1,3}`$
$$\psi \gamma _{21}m\psi \gamma _0=0$$
(12)
where $``$ denotes the Dirac operator, that is,
$$=\gamma _0_0\gamma _1_1\gamma _2_2\gamma _3_3$$
and $`\gamma _\mu `$ ($`\mu =0,1,2,3`$) are interpreted as vectors in $`Cl_{1,3}`$ and $`\psi =\psi (x)Cl_{1,3}^+,xM`$, where $`M`$ is the minkowskian manifold. Now let us to ask a question: How can one simulate the minkowskian Dirac equation in an euclidean formulation? The answer is given for the vee product.
Firstly, multiplying on the right for $`\gamma _{12}`$ we can write (12) as
$$\psi m\psi \gamma _{012}=0$$
(13)
Considering $`\psi Cl_{4,0}^+,xM`$ and using $`𝐞`$-notation for the $`Cl_{4,0}`$ elements, the Dirac equation (13) can be written in the euclidean spacetime using the $``$ product as
$$\psi m\psi 𝐞_{012}=0$$
(14)
where $`=𝐞^\mu _\mu `$ with $`𝐞^\mu =𝐞_\mu `$. Note that $`𝐞_{012}=𝐞_0𝐞_1𝐞_2=𝐞_0𝐞_1𝐞_2`$.
Let us see how this equation appears in terms of the original Clifford product in euclidean spacetime. First we split the Dirac operator in temporal and space parts,
$$\psi =𝐞_0_0\psi +𝐞_i_i\psi $$
Working on the temporal part we have
$`𝐞_0_0\psi `$ $`=`$ $`_0\psi 𝐞_02(_0\psi 𝐞_0)(𝐞_0𝐞_0)`$
$`=`$ $`_0\psi 𝐞_02[{\displaystyle \frac{1}{2}}(_0\psi 𝐞_0𝐞_0_0\psi )]=𝐞_0_0\psi `$
where we have used $`_0\psi 𝐞_0=\frac{1}{2}(_0\psi 𝐞_0𝐞_0_0\psi )`$ For the space part we have
$$𝐞_i_i\psi =_i\psi 𝐞_i2[(_i\psi )(𝐞_0𝐞_i)]=_i\psi 𝐞_i$$
Therefore
$$\psi =𝐞_0_0\psi +_i\psi 𝐞_i$$
(15)
It is easy to see that
$$\psi =\mathrm{}_M\psi $$
(16)
where $`\mathrm{}_M=_0^2_{i=1}^3_i^2`$, while $`^2\psi =\mathrm{}_E\psi `$ with $`\mathrm{}_E=_0^2+_{i=1}^3_i^2`$.
The massive term in (13) will be
$$m\psi 𝐞_{012}=m𝐞_{12}\psi 𝐞_0$$
(17)
Now we can write an equation simulating the minkowskian Dirac equation in $`Cl_{4,0}`$:
$$𝐞_0_0\psi +_i\psi 𝐞_im𝐞_{12}\psi 𝐞_0=0$$
(18)
If one considers a charged fermion field $`\psi `$ in interaction with the electromagnetic field $`A`$ we will add to the Dirac equation (13) the term $`eA\psi \gamma _{12}`$ where $`A^{1,3}`$, applying the vee product we get
$$A\psi 𝐞_{12}=𝐞_{12}(\psi A2(\psi 𝐞_0)A_0)$$
(19)
It is important to remark that we have never changed the algebra – we have working with $`Cl_{4,0}`$. It can be proved that solutions of (13) are solutions of (18) too. For a general $`\psi Cl_{1,3}^+,xM`$, we get in (13) eight coupled differential equations; transforming this original $`\psi `$ solution to its $`Cl_{4,0}^+`$ version and checking it with (18) we recover exactly the same coupled system.
## 5 The General Case
Lounesto has studied the case involving the opposite signatures $`(+)`$ and $`(+++)`$. For $`a,bCl_{1,3}`$ the tilt transformation, based in the even-odd decomposition $`Cl_{1,3}=Cl_{1,3}^+Cl_{1,3}^{}`$ first emphasized by Clifford , was defined as
$$\underset{Cl_{1,3}product}{\underset{}{ab}}\underset{Cl_{3,1}product}{\underset{}{b_+a_++b_+a_{}+b_{}a_+b_{}a_{}}}$$
where the subscripts minus and plus are the odd and even parts. We reinterpret this transformation as a new product $`_t`$ given by
$$A_l_{t(3,1)}B_k=(1)^{kl}B_kA_l$$
(20)
The meaning of this expression is that we are defining a mapping from an algebra $`Cl`$ to its opposite algebra $`Cl^{\mathrm{opp}}`$, and it is not difficult to show that $`Cl_{p,q}^{\mathrm{opp}}=Cl_{q,p}`$.
In order to consider the general case, we consider the vee and tilt products. There are indeed some similarities
$$\mathrm{Vee}\mathrm{product}:A_lB_k=(1)^{kl}[B_kA_l2(B_k𝐞_0)(𝐞_0A_l)]$$
(21)
$$\mathrm{Tilt}\mathrm{product}:A_lB_k=(1)^{kl}B_kA_l$$
(22)
It is interesting to see which is the purpose of the term with the temporal component $`𝐞_0`$. The vee product simulates the change of signature as
$$\underset{Cl_{4,0}}{\underset{}{(++++)}}\underset{Cl_{1,3}}{\underset{}{(+)}}$$
where these signs correspond to the square of the canonical basis elements. All the squares change except for the temporal component. If we set up the inverse problem $`(+)(++++)`$, the general expression for the corresponding vee product will be the same. This is very welcome since, for example, we can get back the euclidean Dirac equation from the minkowskian one as in (18) just by using again the $``$ product. The same holds in the opposite direction, of course.
For the tilt product,
$$\underset{Cl_{1,3}}{\underset{}{(+)}}\underset{Cl_{3,1}}{\underset{}{(+++)}}$$
we change all the squares and now we do not need to subtract any temporal term. Again the opposite problem $`(+++)(+)`$ keeps up the same form. For the most interesting cases in physics we have
$$\begin{array}{cc}(++++)& (+)\\ & \\ ()& (+++)\end{array}$$
As we know the change of signature in the same row is done using the vee product and between rows of the same column using the tilt product. Now if we want to do the change in a diagonal way, i.e. $`(++++)(+++)`$ then we will need to compose vee and tilt product. All this products can be extended to anothers signatures. For example
$$\begin{array}{ccccc}(++++)& (+)& (++)& (+)& ()\\ & & & & \\ ()& (+++)& (++)& (+++)& (++++)\end{array}$$
If we want to change the signature along the same row we only need to know the canonical basis element $`𝐞_\mu `$ which doesn’t change its square, then
$$A_lB_k=(1)^{kl}[B_kA_l2(B_k𝐞_\mu )(𝐞^\mu A_l)]$$
(23)
where $`\mu `$’s are not summed. And for the change between rows in the same column
$$A_l_tB_k=(1)^{kl}B_kA_l$$
(24)
As one can easily see there is no condition limiting this work to four dimensions, so our scheme holds in any dimension.
## 6 Other Applications
While in our discussion of the Dirac equation we were interested in simulating “minkowskian properties” in an euclidean spacetime, in this section we are interested in simulating “euclidean properties” in Minkowski spacetime. Therefore in relation to the notation we are interested in simulating the product of $`Cl_{4,0}`$ in terms of the generators $`\{\gamma _\mu \}`$ ($`\mu =0,1,2,3`$) of $`Cl_{1,3}`$.
### 6.1 The Hodge Star Operator
The Hodge star operator $``$ in Minkowski spacetime can be written using $`Cl_{1,3}`$ as
$$\mathrm{\Phi }=\stackrel{~}{\mathrm{\Phi }}\gamma _5,$$
(25)
where $`\mathrm{\Phi }Cl_{1,3}`$ is a multivector field and $`\gamma _5=\gamma _0\gamma _1\gamma _2\gamma _3`$ and by the tilde we denoted the reversion operation such that
$$\stackrel{~}{A}_k=(1)^{k(k1)/2}A_k$$
(26)
for $`A_k\mathrm{\Lambda }^k`$.
Using the vee product we can write the Hodge star operator $``$ corresponding to euclidean spacetime in terms of the algebra $`Cl_{1,3}`$ of Minkowski spacetime as<sup>2</sup><sup>2</sup>2Note that the Hodge operator in euclidean spacetime is denoted by an asterisk while the one in Minkowski spacetime is denoted by a star.
$$\mathrm{\Phi }=\stackrel{~}{\mathrm{\Phi }}\gamma _5.$$
(27)
It is easy to see that $`\gamma _0\gamma _1\gamma _2\gamma _3=\gamma _0\gamma _1\gamma _2\gamma _3`$ and that $`\gamma _5\gamma _5=1`$ while $`\gamma _5\gamma _5=1`$. In order to rewrite the above expression using the definition of the vee product it is convenient to split $`\mathrm{\Phi }`$ into even and odd parts, that is, $`\mathrm{\Phi }=\mathrm{\Phi }_++\mathrm{\Phi }_{}`$ where $`\mathrm{\Phi }_\pm =\pm \widehat{\mathrm{\Phi }}_\pm `$, where by the hat we denoted the graded involution, that is,
$$\widehat{A}_k=(1)^kA_k,$$
(28)
where $`A_k\mathrm{\Lambda }^k`$. Now we can write using the definition of the vee product that
$$\begin{array}{cc}\hfill \mathrm{\Phi }\gamma _5& =\mathrm{\Phi }_+\gamma _5+\mathrm{\Phi }_{}\gamma _5\hfill \\ & =\gamma _5\mathrm{\Phi }_++2\gamma _{123}(\gamma _0\mathrm{\Phi }_+)+\gamma _5\mathrm{\Phi }_{}+2\gamma _{123}(\gamma _0\mathrm{\Phi }_{}),\hfill \end{array}$$
(29)
and
$$\mathrm{\Phi }_+\gamma _5=\gamma _5\gamma _0\mathrm{\Phi }_+\gamma _0,\mathrm{\Phi }_{}\gamma _5=\gamma _5\gamma _0\mathrm{\Phi }_{}\gamma _0.$$
(30)
We have therefore
$$\mathrm{\Phi }\gamma _5=\gamma _5\gamma _0\widehat{\mathrm{\Phi }}\gamma _0,$$
(31)
and the euclidean Hodge star operator $``$ can be written as
$$\mathrm{\Phi }=\gamma _5\gamma _0\overline{\mathrm{\Phi }}\gamma _0,$$
(32)
where $`\overline{\mathrm{\Phi }}=\widehat{\stackrel{~}{\mathrm{\Phi }}}=\stackrel{~}{\widehat{\mathrm{\Phi }}}`$. Moreover, since the parity operation $`𝒫`$ in $`Cl_{1,3}`$ is given by
$$𝒫(\mathrm{\Phi })=\gamma _0\mathrm{\Phi }\gamma _0,$$
(33)
and if we rewrite (32) as
$$\mathrm{\Phi }=\gamma _0\stackrel{~}{\mathrm{\Phi }}\gamma _5\gamma _0,$$
(34)
we have that
$$\mathrm{\Phi }=𝒫(\mathrm{\Phi }).$$
(35)
This expression clearly shows that $``$ is written only in terms of operations on Minkowski spacetime.
### 6.2 Differential and Codifferential Operators
In Minkowski spacetime the differential $`\mathrm{d}`$ and codifferential $`\delta `$ operators can be written in terms of Dirac operator $``$ as
$$\mathrm{d}\mathrm{\Phi }=\frac{1}{2}\left(\mathrm{\Phi }+\widehat{\mathrm{\Phi }}\stackrel{}{}\right),\delta \mathrm{\Phi }=\frac{1}{2}\left(\mathrm{\Phi }\widehat{\mathrm{\Phi }}\stackrel{}{}\right),$$
(36)
and such that $`=\mathrm{d}+\delta `$, where $`=\gamma ^\mu _\mu `$ and the right action of the Dirac operator denoted by $`\stackrel{}{}`$ is defined as $`\mathrm{\Phi }\stackrel{}{}=(_\mu \mathrm{\Phi })\gamma ^\mu `$. With the above definition we have that $`\delta =\mathrm{d}`$, where $``$ is the Hodge star operator in Minkowski spacetime<sup>3</sup><sup>3</sup>3One can find different definitions for the codifferential operator $`\delta `$ but this is completely irrelevant for our purpose that is to give examples of applications of our method..
Now we want to write the differential and codifferential operators corresponding to the case of euclidean spacetime using the algebra of Minkowski spacetime. Let us denote these operators by $`\stackrel{ˇ}{\mathrm{d}}`$ and $`\stackrel{ˇ}{\delta }`$, respectively, the check being used to distinguish them from their counterparts in Minkowski spacetime. Since the differential operator is defined independently of the existence of a metric structure on a manifold, we have that $`\stackrel{ˇ}{\mathrm{d}}=\mathrm{d}`$; on the other hand, the codifferential operator requires a metric for its definition and therefore we have $`\stackrel{ˇ}{\delta }\delta `$. In order to write the euclidean version of these operators in Minkowski spacetime the recipe is to replace the usual Clifford product by the vee product in formulas (36), that is,
$$\stackrel{ˇ}{\mathrm{d}}\mathrm{\Phi }=\frac{1}{2}\left(\mathrm{\Phi }+\widehat{\mathrm{\Phi }}\stackrel{}{}\right),\stackrel{ˇ}{\delta }\mathrm{\Phi }=\frac{1}{2}\left(\mathrm{\Phi }\widehat{\mathrm{\Phi }}\stackrel{}{}\right),$$
(37)
The expressions $`\mathrm{\Phi }=\gamma ^\mu _\mu \mathrm{\Phi }`$ and $`\mathrm{\Phi }\stackrel{}{}=_\mu \mathrm{\Phi }\gamma ^\mu `$ can be calculated as before. We have for $`\mathrm{\Phi }_k\mathrm{\Lambda }^k`$ that ($`i=1,2,3`$)
$$\mathrm{\Phi }_k=\gamma ^0_0\mathrm{\Phi }+(1)^k(_i\mathrm{\Phi }_k)\gamma ^i,\mathrm{\Phi }_k\stackrel{}{}=(1)^k\gamma ^i_i\mathrm{\Phi }_k+_0\mathrm{\Phi }_k\gamma ^0,$$
(38)
and therefore
$$\mathrm{\Phi }=\gamma ^0_0\mathrm{\Phi }+_i\widehat{\mathrm{\Phi }}\gamma ^i,\mathrm{\Phi }\stackrel{}{\mathrm{\Phi }}=\gamma ^i_i\widehat{\mathrm{\Phi }}+_0\mathrm{\Phi }\gamma ^0.$$
(39)
Now, using the expression for $`\stackrel{ˇ}{\mathrm{d}}`$ given in (37) we see that
$$\begin{array}{cc}\hfill \stackrel{ˇ}{\mathrm{d}}\mathrm{\Phi }& =\frac{1}{2}\left(\gamma ^0_0\mathrm{\Phi }+_i\widehat{\mathrm{\Phi }}\gamma ^i+\gamma ^i_i\mathrm{\Phi }+_0\widehat{\mathrm{\Phi }}\gamma ^0\right)\hfill \\ & =\frac{1}{2}\left(\gamma ^\mu _\mu \mathrm{\Phi }+_\mu \widehat{\mathrm{\Phi }}\gamma ^\mu \right)=\mathrm{d}\mathrm{\Phi },\hfill \end{array}$$
(40)
that is, $`\stackrel{ˇ}{\mathrm{d}}=\mathrm{d}`$ as expected.
On the other hand, using the expression for $`\stackrel{ˇ}{\delta }`$ given in (38) we have that
$$\begin{array}{cc}\hfill \stackrel{ˇ}{\delta }\mathrm{\Phi }& =\frac{1}{2}\left(\gamma ^0_0\mathrm{\Phi }+_i\widehat{\mathrm{\Phi }}\gamma ^i\gamma ^i_i\mathrm{\Phi }_0\widehat{\mathrm{\Phi }}\gamma ^0\right)\hfill \\ & =\frac{1}{2}\left[\left(\gamma ^0_0\mathrm{\Phi }\gamma ^i_i\mathrm{\Phi }\right)\left(_0\widehat{\mathrm{\Phi }}\gamma ^0_i\widehat{\mathrm{\Phi }}\gamma ^i\right)\right],\hfill \end{array}$$
(41)
and clearly $`\stackrel{ˇ}{\delta }\delta `$, as expected. Using (32) for the euclidean Hodge star operator $``$ written in Minkowski spacetime we have that
$$\mathrm{d}\mathrm{\Phi }=\frac{1}{2}\left(\gamma _5\gamma ^\mu \gamma _0_\mu \overline{\mathrm{\Phi }}\gamma _0+\gamma _5\gamma _0_\mu \stackrel{~}{\mathrm{\Phi }}\gamma _0\gamma ^\mu \right)$$
(42)
and therefore
$$\begin{array}{cc}\hfill \mathrm{d}\mathrm{\Phi }& =\frac{1}{2}\left(\gamma _5_\mu \mathrm{\Phi }\gamma _0\gamma ^\mu \gamma _5\gamma _0+\gamma _5\gamma _0\gamma ^\mu \gamma _0_\mu \widehat{\mathrm{\Phi }}\gamma _5\right)\hfill \\ & =\frac{1}{2}\left(\gamma _0\gamma ^\mu \gamma _0_\mu \mathrm{\Phi }_\mu \widehat{\mathrm{\Phi }}\gamma _0\gamma ^\mu \gamma _0\right)\hfill \\ & =\frac{1}{2}\left[\left(\gamma ^0_0\mathrm{\Phi }\gamma ^i_i\mathrm{\Phi }\right)\left(_0\widehat{\mathrm{\Phi }}\gamma ^0_i\widehat{\mathrm{\Phi }}\gamma ^i\right)\right],\hfill \end{array}$$
(43)
and comparing this with (37) we have
$$\stackrel{ˇ}{\delta }=\stackrel{ˇ}{\mathrm{d}}=\mathrm{d},$$
(44)
as expected.
### 6.3 Self-Dual and Anti-Self-Dual Solutions of Gauge Field Equations
Since for what follows it is irrelevant whether or not we are considering abelian or non-abelian gauge fields, we shall consider the electromagnetic field – U(1) gauge fields – as an example for our discussion.
Let us consider the free Maxwell equations in Minkowski spacetime,
$$\mathrm{d}F=0,\delta F=0.$$
(45)
These equations don’t have self-dual and/or anti-self-dual solutions, that is, solutions satisfying the conditions $`F=\pm F`$ in the real case. Such solutions are possible in Minkowski spacetime only if we complexify the underlying algebra considering a complex $`F`$. On the other hand, self-dual and/or anti-self-dual solutions of Maxwell equations exist in euclidean spacetime without the need of any complexification. This can be described in our scheme by writing the equivalent Maxwell equations in Minkowski spacetime using the vee product.
The equations
$$\stackrel{ˇ}{\mathrm{d}}F=0,\stackrel{ˇ}{\delta }F=0,$$
(46)
are written in terms of the real Clifford algebra of Minkowski spacetime and admit self-dual and anti-self-dual solutions satisfying $`F=\pm F`$. This condition reads
$$F=\pm \gamma _5\gamma _0\overline{F}\gamma _0=\pm \gamma _5\gamma _0\stackrel{~}{F}\gamma _0.$$
(47)
Now writting
$$F=E+\gamma _5B,$$
(48)
where
$$E=\frac{1}{2}(F\gamma _0F\gamma _0),\gamma _5B=\frac{1}{2}(F+\gamma _0F\gamma _0),$$
(49)
and such that $`E\gamma _0=\gamma _0E`$ and $`\gamma _5B\gamma _0=\gamma _0\gamma _5B`$, we can see that (47) is satisfied for
$$E=\pm B,$$
(50)
which are the usual self-dual and anti-self-dual solutions of the euclidean case but now written in terms of Minkowski spacetime.
## 7 Conclusions
Given the Clifford algebra of a quadratic space with a given signature, we have defined a new product in this structure such that it simulates the Clifford product of a quadratic space with another arbitrary signature different from the original one. We have used this in order to give an algebraic approach to the so called Wick rotation. We have used this new product in order to simulate the product associated with the Minkowski spacetime in terms of the Clifford algebra of the euclidean spacetime. We have also shown how to write the minkowskian Dirac equation in euclidean spacetime and in the other way how to write the Hodge star operator and the differential and codifferential operators corresponding to the euclidean case in terms of Minkowski spacetime, discussing self-dual and anti-self-dual solutions for the gauge field equations in this case.
Acknowledgments: D.M. is grateful to Departamento de Matemática Aplicada, Universidade Estadual de Campinas (UNICAMP), for hospitality and support during the preparation of this work. J.M.P. acknowledges support from the Spanish Ministry of Education contract No. PB96-0384 and the Institut d’Estudis Catalans. |
warning/0003/astro-ph0003344.html | ar5iv | text | # Two-dimensional galaxy-galaxy lensing: a direct measure of the flattening and alignment of light and mass in galaxies
## 1. Introduction
Current observations require the existence of dark matter halos for galaxies. However, fundamental parameters such as their total mass and spatial extent are not well constrained. The mass distribution in galaxies is primarily probed via dynamical tracers of the galactic potential on various scales: stars in the inner regions, HI gas in regions outside the optical radius, and orbital motions of bound satellites in the outer-most regions (e.g. Zaritsky & White 1994). Probes of halo structure at radii devoid of any luminous tracers are therefore needed, weak gravitational lensing offers precisely that.
Galaxy-galaxy lensing, the preferential tangential alignment of the images of background galaxies around bright foreground ones, is detected statistically by stacking galaxies. The first observational attempt to look for galaxy-galaxy lensing was made by Tyson et al. (1984). Recent studies have been very successful, and a signal at the 99.5% confidence level was first reported by Brainerd, Blandford & Smail (1996) using deep ground-based CCD data. Several subsequent studies using Hubble Space Telescope images and ground-based data (Griffiths et al. 1996; Dell’ Antonio & Tyson 1996; Hudson et al. 1998; Casertano 1999; Ebbels et al. 2000; Fischer et al. 1999 \[Sloan Digital Sky Survey (SDSS hereafter) commissioning data\]) report unambiguous detection of a galaxy–galaxy lensing signal.
In the method proposed in this letter, additional information that is available but not exploited in current galaxy-galaxy lensing studies is utilized, namely, the light distribution of galaxies selected to be foreground lenses. The general results derived by Schneider & Bartelmann (1997) are used to relate the shear field to the mass multipole moments. We show that the shapes and orientations of the foreground galaxies (probes of the light) can be compared statistically to that of the shear field (probe of the mass), thus providing a direct method to compare the ellipticity of the light to that of the mass as well any potential misalignments between them reliably. These parameters offer an important clue to the galaxy formation process, since they provide a quantitative measure of the importance of dissipation in the assembly of galaxies. Any variation of the flattening of the total mass (predominantly the dark matter component) with radius is a probe of the efficiency of angular momentum transfer to the dark halo and might provide insights into the structure and composition of galaxy halos. With regard to the relative orientation of the light and mass in galaxies, the two components are likely to be aligned on average and any misalignments might occur transiently, for instance, after a violent merger.
The outline of the paper is as follows: in §2, the current status of modeling in galaxy-galaxy lensing studies is reviewed. In §3 the formalism used to extract the shape of the mass distribution is presented, with the application to an elliptical isothermal mass model in §4. The feasibility of detection of the flattening of the mass is examined in §5, and a measure of the alignment of mass and light is discussed in §6. We conclude in §7 with a discussion of the importance of studying shapes of dark halos for understanding key issues in the formation and structure of galaxies and the relation between the luminous and dark component in galaxies.
## 2. Current Status of galaxy-galaxy lensing
The galaxy-galaxy lensing signal is measured ideally from a deep image by selecting a population of brighter (assumed to be foreground) galaxies as lenses and measuring the induced shape distortion in the fainter (background) galaxies. The shear signal obtained from direct averaging in radial bins around the bright foreground galaxies is then stacked to obtain the radial profile of the shear $`\gamma `$ as a function of distance from the lens. This provides reasonable constraints on the circular velocity of a fiducial halo, found to be in the range 210 km s<sup>-1</sup> – 250 km s<sup>-1</sup> for a typical $`L^{}`$ galaxy, but is found to be fairly insensitive to the radial extent of the halos (consistent with halo sizes in excess of 100h<sup>-1</sup> kpc). In analyses to date, the errors are dominated by shot noise and error arising from the unknown redshift distribution of galaxies. However, with forthcoming surveys like the SDSS, which plan to image over a hundred million galaxies in many bands and provide reliable photometric redshifts, the prospects for galaxy-galaxy lensing studies are extremely good (Fischer et al. 1999).
## 3. Extracting shape parameters for the mass distribution
The mass distribution of a lensing galaxy is described by the convergence field $`\kappa (𝐱)`$, defined to be the surface mass-density $`\mathrm{\Sigma }(𝐱)`$ expressed in units of the critical surface mass density $`\mathrm{\Sigma }_{\mathrm{crit}}`$. The critical surface mass density depends on the precise configuration, i.e. on the angular diameter distances from the lens to the source $`D_{\mathrm{ls}}`$, observer to source $`D_{\mathrm{os}}`$ and observer to lens $`D_{\mathrm{ol}}`$. It is given by, $`\mathrm{\Sigma }_{\mathrm{crit}}=\frac{c^2}{4\pi G}\frac{D_{\mathrm{os}}}{D_{\mathrm{ls}}D_{\mathrm{ol}}}`$. Standard galaxy-galaxy lensing provides a measure of the mass $`M`$ within an aperture, which is given by
$$M=d^2xw(x)\kappa (𝐱),$$
(1)
where $`w(x)`$ is a weight function normalized so that $`d^2xw(x)=1`$. It is chosen to be continuous, differentiable and is required to fall-off rapidly to zero outside the aperture scale $`\beta `$. The shape parameters of the mass distribution are characterized by the quadrupole moments of the convergence $`\kappa (𝐱)`$ within the aperture (Schneider & Bartelmann 1997), which are defined as
$`Q_{ij}{\displaystyle d^2xx_ix_jw(x)\kappa (𝐱)},`$ (2)
This tensor can be decomposed into a trace-free piece $`Q`$ and a trace $`T`$ defined as:
$`Q=Q_{11}Q_{22}+\mathrm{\hspace{0.17em}2}iQ_{12},T=Q_{11}+Q_{22}.`$ (3)
The ellipticity of the mass $`ϵ_\kappa `$ is then simply,
$`ϵ_\kappa ={\displaystyle \frac{Q}{T}}={\displaystyle \frac{(a_\kappa ^2b_\kappa ^2)}{(a_\kappa ^2+b_\kappa ^2)}}e^{i\phi _\kappa },`$ (4)
where $`a_\kappa `$ and $`b_\kappa `$ are, respectively, the major and minor axes of the mass distribution, and $`\phi _\kappa `$ is its position angle relative to the positive $`x`$-axis.
Schneider & Bartelmann (1997) have shown that the multipole moments of $`\kappa (𝐱)`$ computed from the observed shear $`\gamma (𝐱)`$ field. In particular, the quadrupole moments (Eq. ) can be expressed as
$`Q={\displaystyle d^2xe^{2i\phi }\left[g_t(x)\gamma _t(𝐱)+ig_\times (x)\gamma _\times (𝐱)\right]},`$ (5)
where the rotated shear components $`\gamma _t`$ and $`\gamma _\times `$ correspond, respectively, to a tangential and curl-type shear pattern about the center of mass of the lens (see Rhodes, Refregier & Groth 2000 for an illustration). They are related to the unrotated components by,
$`\gamma _t`$ $`=`$ $`\left[\mathrm{cos}(2\phi )\gamma _1+\mathrm{sin}(2\phi )\gamma _2\right]`$
$`\gamma _\times `$ $`=`$ $`\left[\mathrm{sin}(2\phi )\gamma _1+\mathrm{cos}(2\phi )\gamma _2\right]`$ (6)
where $`\phi `$ is the polar angle from the $`x`$-axis. The associated aperture functions $`g_t(x)`$ and $`g_\times (x)`$ are given by
$`g_t(x)=2V_2(x)x^2w(x),g_\times (x)=2V_2(x),`$ (7)
where $`V_n(x)=x^2_0^x𝑑x^{}x^{n+1}w(x^{})`$. Similarly, the trace part $`T`$ and the mass $`M`$ can also be written as
$`T={\displaystyle d^2xg_t(x)\gamma _t(𝐱)},M={\displaystyle d^2xh_t(x)\gamma _t(𝐱)},`$ (8)
where $`h_t(x)=2V_0(x)w(x)`$.
## 4. Application to the elliptical isothermal model
We consider an isothermal model with concentric elliptical equipotentials (Natarajan & Kneib 1996). The projected potential for this model is $`\psi =\alpha r`$, where $`\alpha `$ is the Einstein radius and $`r`$ is a generalized elliptical radius. If the $`x`$-axis is aligned with the major axis of the potential, the generalized radius is given by $`r^2=\frac{x_1^2}{1+ϵ}+\frac{x_2^2}{1ϵ}`$, where $`ϵ`$ is the ellipticity of the equipotentials. The Einstein radius is related to the velocity dispersion of the galaxy $`\sigma _v`$ by $`\alpha =4\pi (\frac{\sigma _v}{c})^2(\frac{D_{\mathrm{ls}}}{D_{\mathrm{os}}})`$, and is of the order of $`1^{\prime \prime }`$ for galaxies.
For a weakly elliptical model ($`ϵ1`$), the potential has the form,
$`\psi (𝐱)\alpha x[1{\displaystyle \frac{ϵ}{2}}\mathrm{cos}\mathrm{\hspace{0.17em}2}(\phi \phi _0)];`$ (9)
where $`\phi _0`$ is the position angle of the potential, and reduces to that of a singular isothermal sphere in the circular limit ($`ϵ=0`$). Current observational limits on ellipticities of halos have been compiled in a comprehensive recent review by Sackett (1999), however, there are no constraints on the dispersion in the shape parameters, therefore, for the purposes of this calculation we have used the central value of $`ϵ=0.3`$. Besides, higher order terms will be approximately $`(0.3)^2=0.09`$ which are still only corrections at the 10% level, consistent with the assumption of small $`ϵ`$.
Restricting our analysis to the weak regime, to first order in $`ϵ`$, the associated convergence $`\kappa =^2\psi /2`$, where $`\psi =\alpha r`$, is given by,
$`\kappa (𝐱){\displaystyle \frac{\alpha }{2x}}[1+{\displaystyle \frac{3ϵ}{2}}\mathrm{cos}\mathrm{\hspace{0.17em}2}(\phi \phi _0)],`$ (10)
and the complex shear $`\gamma =\gamma _1+i\gamma _2=[_1^2_2^2+2i_1_2]\psi /2`$ is
$`\gamma {\displaystyle \frac{\alpha }{2x}}[1+{\displaystyle \frac{3ϵ}{2}}\mathrm{cos}2(\phi \phi _0)]e^{2i\phi },`$ (11)
The rotated shear components are thus
$`\gamma _t{\displaystyle \frac{\alpha }{2x}}[1+{\displaystyle \frac{3ϵ}{2}}\mathrm{cos}2(\phi \phi _0)],\gamma _\times =0,`$ (12)
yielding a tangential shear modulated by an elliptical pattern.
The ellipticity of the underlying mass distribution $`\kappa (𝐱)`$ needs to be related to that of the projected (2-d) potential $`\varphi (𝐱)`$. Using a normalized gaussian as the weight function $`w(x)=e^{x^2/2\beta ^2}/(2\pi \beta ^2)`$, we evaluate the integral for the quadrupole (Eq. ) and monopole (Eq. ) moments of $`\kappa (𝐱)`$,
$`M=\sqrt{{\displaystyle \frac{\pi }{8}}}{\displaystyle \frac{\alpha }{\beta }},|Q|={\displaystyle \frac{3}{8}}\sqrt{{\displaystyle \frac{\pi }{2}}}\alpha \beta ϵ,T=\sqrt{{\displaystyle \frac{\pi }{8}}}\alpha \beta `$ (13)
The ellipticity of the mass (Eq. ) is thus $`ϵ_\kappa =\frac{3ϵ}{4}`$. The ellipticity of the potential $`ϵ_\psi `$, obtained similarly by taking moments of $`\psi `$, is $`ϵ_\psi =ϵ/4`$. Note that the ellipticity of the potential $`ϵ_\varphi `$ computed above is smaller than $`ϵ`$ (by a factor of 4), as it is weighted by the circular gaussian window function. Comparing, the weighted ellipticities, we find that $`ϵ_\kappa >ϵ_\psi `$, as expected since equi-potentials are always rounder than the mass contours.
## 5. Measuring the flattening of the mass distribution
We now show how these results can be used to measure the flattening of the mass distribution. As in ordinary galaxy-galaxy lensing, the galaxy catalog is separated into a foreground and a background subsample, using magnitude, colors or photometric redshifts. The ellipticity of the galaxies in both subsamples is then measured in the by taking second moments of the light distribution. The ellipticities of the foreground sample yields the ellipticity of the light $`ϵ_l`$ associated with each lens. While $`ϵ_l`$ is ignored in ordinary galaxy-galaxy lensing, we instead align the foreground galaxies along their major axes before stacking. We then measure the average ellipticity of the mass $`ϵ_\kappa `$ as described above, by replacing the integrals in equations (5) and (8) by sums over the sheared background galaxies. This yields a measurement of the component of the average ellipticity of the mass, $`ϵ_\kappa `$ parallel to the that of the light, i.e.
$$ϵ_\kappa =\mathrm{Re}ϵ_\kappa ^{}\widehat{ϵ}_l,$$
(14)
where the ellipticities are taken to be complex numbers with $`ϵ=ϵ_1+iϵ_2`$, denotes complex conjugation, and $`\widehat{ϵ}_l=ϵ_l/|ϵ_l|`$ is the unit ellipticity of the light.
We now compute the uncertainty in measuring $`ϵ_\kappa `$, by taking the square of the mean of the discrete estimators for $`M`$, $`T`$ and $`Q_{}=\mathrm{Re}(Q)`$, and converting back into integrals (Schneider & Bartelmann 1997). In the absence of lensing, we find
$`\sigma ^2[M]`$ $`=`$ $`{\displaystyle \frac{\sigma _ϵ^2}{n_bn_fA}}{\displaystyle d^2xh_r^2(x)},`$
$`\sigma ^2[T]`$ $`=`$ $`{\displaystyle \frac{\sigma _ϵ^2}{n_bn_fA}}{\displaystyle d^2xg_r^2(x)},`$
$`\sigma ^2[Q_{}]`$ $`=`$ $`{\displaystyle \frac{\sigma _ϵ^2}{2n_bn_fA}}{\displaystyle d^2x\left[g_r^2(x)+g_\times ^2(x)\right]},`$ (15)
where $`\sigma _ϵ^2=ϵ_r^2=ϵ_\times ^2`$ is the variance of the intrinsic ellipticity distribution of galaxies ($`0.3^2`$), $`n_b`$ and $`n_f`$ are respectively the number density of background and foreground galaxies, and $`A`$ is the area covered by the survey.
For the elliptical isothermal model with the gaussian weight function, we can evaluate these integrals and find,
$`\sigma ^2[M]`$ $`=`$ $`{\displaystyle \frac{\sigma _ϵ^2}{4\pi n_bn_fA\beta ^2}},\sigma ^2[T]={\displaystyle \frac{\sigma _ϵ^2\beta ^2}{2\pi n_bn_fA}},`$
$`\sigma ^2[Q_{}]`$ $`=`$ $`{\displaystyle \frac{3\sigma _ϵ^2\beta ^2}{4\pi n_bn_fA}}`$ (16)
By propagating these errors in the definition of the ellipticity of the mass (Eq. ), we can compute the signal to noise ratio $`\left(S/N\right)_{ϵ_\kappa }=ϵ_\kappa /\sigma [ϵ_\kappa ]`$ for measuring $`ϵ_\kappa `$, and find, to first order in $`ϵ`$,
$`\left({\displaystyle \frac{S}{N}}\right)_{ϵ_\kappa }`$ $``$ $`4.6\left({\displaystyle \frac{ϵ_\kappa }{0.3}}\right)\left({\displaystyle \frac{\alpha }{0.5^{\prime \prime }}}\right)\left({\displaystyle \frac{n_b}{1.5\mathrm{arcmin}^2}}\right)^{\frac{1}{2}}`$
$`\left({\displaystyle \frac{n_f}{0.035\mathrm{arcmin}^2}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{0.3}{\sigma _ϵ}}\right)\left({\displaystyle \frac{\mathrm{A}}{1000\mathrm{d}\mathrm{e}\mathrm{g}^2}}\right)^{\frac{1}{2}}.`$
Here, we have chosen to scale $`ϵ_\kappa `$ in units of 0.3, which is reasonable given current observational limits on the flattening of dark matter halos (see Table 3 in the Sackett (1999) review and references therein; Buote & Canizares 1998). Additionally, in these scalings, we have used the survey specifications (ellipticity dispersion, number density of foreground lenses $`n_f`$ and the number density of background galaxies $`n_b`$, and approximate observed Einstein radius) quoted for the SDSS commissioning run provided by Fischer et al. (1999), with a modestly expanded area (1000 deg<sup>2</sup>) from the current area of 225 deg<sup>2</sup>. In the context of estimating the scatter arising in the mass estimates from galaxy-galaxy lensing due to halo shapes, Fischer et al. (2000) mention in passing that with 10 times more data than the commissioning run, halo shapes can be measured; this figure is comparable to our estimate of the $`S/N`$.
Note that since these numbers have been taken from the SDSS commissioning run (which suffered from poor seeing), are conservative and do take into account several observational errors. The dispersion of 0.3 in the ellipticity distribution includes a correction for the seeing (see for instance weak lensing observations of Rhodes et al. (2000) and Bacon et al. (2000) and references therein). The effect of seeing and noise are also reflected in the modest number density assumed for background sources. Note, however, that in contrast to the case of ordinary galaxy-galaxy lensing, foreground galaxies will need to be sufficiently resolved to measure their shapes prior to alignment and stacking. This could induce some degradation in the signal, but since the foreground galaxies are typically brighter and larger, this effect is expected to be small.
The shape parameters of the mass will therefore be easily detectable with SDSS in the near future (McKay et al. 2000). For the total SDSS area of $`10^4`$ deg<sup>2</sup>, the significance rises to 15$`\sigma `$. This in fact implies that potentially, even the radial dependence of the flattening can be studied by considering annuli-shaped weight functions (for instance, the difference of two gaussians). Moreover, the degree of flattening as a function of the morphological galaxy type can also be studied.
It is interesting to compare the $`(S/N)_{ϵ_\kappa }`$ expected for measuring $`ϵ_\kappa `$ estimated above with that of the usual galaxy-galaxy lensing $`(S/N)_M=M/\sigma [M]`$ which measures the mass enclosed within an aperture. For the model considered here, we find the following relation,
$`\left({\displaystyle \frac{S}{N}}\right)_{ϵ_\kappa }=\mathrm{\hspace{0.17em}0.17}\left({\displaystyle \frac{ϵ_\kappa }{0.3}}\right)\left({\displaystyle \frac{S}{N}}\right)_M.`$ (18)
Therefore, shape parameters can be measured with a significance which is smaller but comparable to that of the enclosed mass. Note that for the current SDSS survey area of $`A=225`$ deg<sup>2</sup>, we find $`(S/N)_M13`$ which agrees roughly with the significance of the reported Fischer et al. (1999) results, when averaged over all radial bins.
## 6. Measuring the alignment of light and mass
Given the good prospects expected from the above results, one can be more ambitious and try to characterize the alignment of mass and light in more detail. We can make use of the amplitude of the light ellipticity $`ϵ_l`$, which we have not used until now. This can be done by grouping the lens galaxies into several $`ϵ_l`$-bins, and computing $`ϵ_\kappa `$ separately for each bin. A more direct approach would be to consider the correlation function of the ellipticities of the mass and light, defined as
$$C_{\kappa l}=\mathrm{Re}ϵ_\kappa ^{}ϵ_l,$$
(19)
where the average is over all lens galaxies, and $`ϵ_\kappa `$ is an estimate of the mass ellipticity of each lens derived from its associated background galaxies. While $`ϵ_\kappa `$ for an individual lens is rather noisy, a significant measurement of $`C_{\kappa l}`$ can be obtained by averaging over a large number of lens galaxies. This correlation function could also be computed for several annuli, and, would therefore, yield a direct measure of the radial dependence of the alignment of mass and light.
## 7. Discussion
The shapes of dark matter halos (see Sackett (1999) for a more comprehensive review) have been probed via many techniques and the consensus from these studies is that the precise shapes offer important clues to both the galaxy formation process and perhaps, even to the nature of dark matter. Cosmological N-body simulations suggest that dark matter halos are triaxial and that dissipation determines their shape. High resolution simulations find that the effect of dissipation (Katz & Gunn 1991; Dubinski 1994) is to transform an initially triaxial halo from prolate-triaxial to oblate-triaxial, while preserving the degree of flattening, yielding on average rounder and more oblate dark halos than those in dissipationless simulations.
Comparing the shape of the mass profile inferred from X-ray data for a sample of ellipticals with that of the light, Buote & Canizares (1994; 1998) find that the dark matter is at least as flattened as the light and is definitely more extended. The origin of the X-ray isophotal twist in the case of NGC720, they argue reflects an intrinsic misalignment of the stars with the dark matter. Keeton, Kochnek & Seljak (1997) incorporate the shape of the light distribution as a constraint in modeling individual lenses (that produce multiple images of background quasars) and find that an additional component of shear is required to match the observations. They speculate that this component could arise from a misalignment between the luminous galaxy and its dark matter halo. Our proposed technique will provide reliable measurements of the shape and orientation of light and mass in galaxies, thereby aiding in the understanding of the coupling of baryons with the dark matter.
We thank Martin Rees and Donald Lynden-Bell for useful discussions. PN acknowledges support from a Trinity College Research Fellowship and AR from the EEC Lensing Network for a TMR post-doctoral fellowship and a Wolfson College Research Fellowship. |
warning/0003/astro-ph0003189.html | ar5iv | text | # Angular size in “quintessence” cosmology
## 1 Introduction
Recent data from SNe Ia have provided strong evidence for an expanding Universe speeding up, rather than slowing down (Riess et al. 1998; Perlmutter et al. 1998). These observational evidences have stimulated great interest in a more general class of cosmological models driven by nonrelativistic matter and a “quintessence” component, i.e., an exotic fluid with an arbitrary equation of state $`p_x=\omega _x\rho _x`$ ($`\omega _x1`$), which probably dominates the bulk of matter in the observed Universe. Examples of these models include the evolving scalar field (Ratra & Peebles 1988; Frieman et al. 1998; Caldwell et al. 1998), the smooth noninteracting component (xCDM) (Turner & White 1997; Chiba et al. 1997), and still the frustated network of topological defects in which $`\omega _x=\frac{n}{3}`$, being $`n`$ the dimension of the defect (Spergel & Pen 1997). Some observational aspects of these models have extensively been analyzed in the literature. For example, Waga & Miceli (1999), combining statistics of gravitational lenses and SNe Ia data have found $`\omega _x<0.7`$ ($`68\%`$ cl) for a spatially flat Universe. Efstathiou (1999), by using high-z Type Ia supernovae and cosmic microwave background anisotropies, has found $`\omega _x<0.6`$ (2$`\sigma `$) if the Universe is assumed to be spatially flat, or $`\omega _x<0.4`$ (2$`\sigma `$) for universes of arbitrary spatial curvature. Perlmutter et al. (1999) constrained $`\omega _x<0.6`$ ($`95\%`$ cl) using large-scale structure and SNe Ia in a spatially flat geometry. However, although carefully investigated in many of their theoretical and observational aspects, the influence of a “quintessence” component in some kinematic tests like the angular size-redshft relation still remains to be analyzed. In principle, the lensing effect of the expanding Universe may provide strong limits on the free parameter describing this exotic component. Therefore, it is interesting to explore how uncertaints in distance measures of extragalactic objects and their underlying evolutionary effects may alter the standard cold dark matter results.
On the other hand, the existing angular size data for distant objects are until nowadays somewhat controversial, specially because they envolve at least two kinds of observational dificulties. First, any high redshift object may have a wide range of proper sizes, and, second, evolutionary and selection effects probably are not negligible. Indeed, the $`\mathrm{\Theta }(z)`$ relation for some extended sources samples seems to be quite imcompatible with the predictions of the standard FRW model when the latter effects are not taken into account (Kapahi 1987;1989). There have also been some claims that the best fit model for the observed distribution of high redshifts extended objects is provided by the standard Einstein-de Sitter universe ($`q_o=\frac{1}{2}`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$) with no significant evolution (Buchalter et al. 1998). However, all these results are in contradiction with the recent observations from type Ia supernovae. Indeed, such data seem to ruled out world models filled only by baryonic matter, and more generally, any model with positive deceleration parameter. The same happens with the corresponding bounds using the ages of old high redshift galaxies (Dunlop et al. 1996; Krauss 1997; Alcaniz & Lima 1999).
The case for compact radio sources is also of great interest. These objects seem to be less sensitive to evolutionary effects since they are short-lived ($`10^3yr`$) and much smaller than their host galaxy. Initially, the data from a sample of 82 objects gave remarkable suport for the Einstein-de Sitter Universe (Kellerman 1993). However, some analysis suggest that, although compatible with an Einstein-de Sitter Universe, the Kellerman data cannot rule out a significant part of the $`\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$ plane (Kayser 1995). Some authors have also argued that models where $`\mathrm{\Theta }(z)`$ diminishes and after a given $`z`$ remains constant may also provide a good fit to Kellerman’s data. In particular, by analysing a subset of 59 compact sources within the same sample, Dabrowski et al. (1995) found that no useful bounds on the value of the deceleration parameter $`q_o`$ can be derived. Indeed, even considering that Euclidean angular sizes ($`\mathrm{\Theta }z^1`$) are excluded at 99$`\%`$ confidence level, and that the data are consistent with $`q_o=1/2`$, they apparently do not rule out higher values of the deceleration parameter (Stephanas & Saha 1995). More recently, based in a more complete sample of data, which include the ones originally obtained by Kellermann, it was argued that the $`\mathrm{\Theta }(z)`$ relation may be consistent with any model of the FRW class with deceleration parameter $`0.5`$ (Gurvits et al. 1999).
In this context, we discuss the influence of a “quintessence” component (Q-model) on the angular size-redshift relation. Particular emphasis is given for the critical redshift at which the angular size of an extragalactic source takes its minimal value. In the limiting case ($`\omega _x=1`$), the results previously derived by Krauss & Schramm (1993) for a flat universe with cosmological constant ($`\mathrm{\Lambda }`$CDM) are recovered. For comparison, we also consider the case of an open model dominated by nonrelativistic matter (OM).
## 2 Angular size and “quintessence”
Let us now consider the FRW line element $`(c=1)`$
$$ds^2=dt^2R^2(t)[d\chi ^2+S_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)],$$
(1)
where $`\chi `$, $`\theta `$, and $`\varphi `$ are dimensionless comoving coordinates, $`R(t)`$ is the scale factor, and $`S_k(\chi )`$ depends on the curvature parameter ($`k=0`$, $`\pm 1`$). The later function is defined by one of the following forms: $`S_k(\chi )=\mathrm{sinh}(\chi )`$, $`\chi `$, $`\mathrm{sin}\chi `$, respectively, for open, flat and closed Universes.
In this background, the angular size-redshift relation for a rod of intrinsic length $`D`$ is easily obtained by integrating the spatial part of the above expression for $`\chi `$ and $`\varphi `$ fixed. One finds
$$\theta (z)=\frac{D(1+z)}{R_oS_k(\chi )}.$$
(2)
The dimensionless coordinate $`\chi `$ is given by
$$\chi (z)=\frac{1}{H_oR_o}_{(1+z)^1}^1\frac{dx}{xE(x)},$$
(3)
where $`x=\frac{R(t)}{R_o}=(1+z)^1`$ is a convenient integration variable. For flat Q-models, the dimensionless function $`E(x)`$ takes the following form
$$E_Q(x)=\left[(1\mathrm{\Omega }_x)x^1+\mathrm{\Omega }_xx^{(1+3\omega _x)}\right]^{\frac{1}{2}},$$
(4)
where $`\mathrm{\Omega }_x=\frac{8\pi G\rho _x}{3H_o^2}`$ is the present day density parameter associated with the “quintessence” component. Observe that the flat constraint condition, $`\mathrm{\Omega }_M+\mathrm{\Omega }_x=1`$, where $`\mathrm{\Omega }_M=\frac{8\pi G\rho _M}{3H_o^2}`$ has been explicitly used in the derivation of (4).
Before proceed further, it is interesting to make explicit the connection with some special cases already established in the literature. If $`\mathrm{\Omega }_x=0`$, or still if $`\mathrm{\Omega }_x=1`$ and $`\omega _x=0`$, one obtains from (2)-(4) the angular diameter expression of the Einstein-de Sitter universe (Sandage 1988)
$$\mathrm{\Theta }(z)=\frac{DH_o(1+z)^{\frac{3}{2}}}{2\left[(1+z)^{\frac{1}{2}}1\right]}.$$
(5)
If $`\omega _x=1`$, the Q-model reduces to a $`\mathrm{\Lambda }`$CDM universe, the details of which has been analysed by Krauss & Schramm (1993). In particular, if the pair $`(\mathrm{\Omega }_x,\omega _x)=(1,1)`$, the angular diameter of this Q-model is the same of a flat universe with a pure cosmological constant, namely
$$\mathrm{\Theta }(z)=\frac{DH_o(1+z)}{z}.$$
(6)
We recall that expression (5) yields a well-known result that the angular diameter in Einstein-de Sitter model has a minimum at $`z_m=5/4`$ (Hoyle 1959), whereas (6) shows us that the extreme Q-model, ($`\mathrm{\Omega }_x,\omega _x`$)=($`1,1`$), has no minimum at all ($`z_m=\mathrm{}`$). Indeed, for any expanding FRW type cosmology, the typical behavior of the angular size relation is the existence of a critical redshift greater than the above Einstein-de Sitter value. We also observe that a new analytical result is obtained by taking $`\omega _x=1/3`$ for arbitraries values of $`\mathrm{\Omega }_x`$. From (2)-(4) one finds
$`\mathrm{\Theta }(z)`$ $`=`$ $`{\displaystyle \frac{DH_o(1+z)}{2\sqrt{\mathrm{\Omega }_x}}}\{\mathrm{ln}[\sqrt{\alpha }+\sqrt{\sqrt{\alpha }+1}`$ (7)
$`{\displaystyle \frac{\sqrt{\alpha }}{\sqrt{(1+z)}}}\sqrt{{\displaystyle \frac{\sqrt{\alpha }}{\sqrt{(1+z)}}}+1}]\}^1.`$
where $`\alpha =\frac{\mathrm{\Omega }_x}{1\mathrm{\Omega }_x}`$.
The expression (2) for $`\mathrm{\Theta }(z)`$ cannot be written in simple analytical form, unless the pair of parameters ($`\mathrm{\Omega }_x`$, $`\omega _x`$) take the above mentioned values. For generic cases, the results can be obtained only by numerical treatment.
In Fig.1 we show a log-log plot of angular size versus redshift for flat Q-models with $`\mathrm{\Omega }_x=0.7`$ and some selected values of $`\omega _x`$. For comparison we have also considered the standard OM cosmology ($`\mathrm{\Omega }_M=0.3`$). As can be seen there, for all values of $`\omega _x`$, the angular size initially decreases with increasing $`z`$, reaches its minimum value at a given $`z_m`$, and eventually begins to increase for fainter magnituds. Note also that the standard OM behavior may be interpreted as an intermediary case between $`\mathrm{\Lambda }`$CDM ($`\omega _x=1`$) and a Q-model with $`\omega _x0.5`$, though its critical redshift is displaced to higher values.
## 3 The critical redshift
As widely known, the existence of a critical redshift $`z_m`$ on the angular size-redshift relation may qualitatively be understood in terms of an expanding space. The light observed today from a source at high $`z`$ was emitted when the object was closer. The relevant aspect here is how this effect may be quantified in terms of the $`\omega _x`$ parameter. To analyze the sensivity of the critical redshift to “quintessence”, we addopt here an approach different of the one applied by Krauss & Schramm (1993) to the case of a flat $`\mathrm{\Lambda }`$CDM universe.
The redshift $`z_m`$ at which the angular size takes its minimal value is the one cancelling out the derivative of $`\mathrm{\Theta }`$ with respect to $`z`$. Hence, from (2) we have the condition
$$S_k(\chi _m)=(1+z_m)S_k^{}(\chi _m),$$
(8)
where $`S_k^{}(\chi )=\frac{S_k}{\chi }\frac{\chi }{z}`$, a prime denotes differentiation with respect to $`z`$ and by definition $`\chi _m=\chi (z_m)`$. Observe also that (3) can readily be differentiated yielding
$$(1+z_m)\chi _m^{}=(R_oH_o)^1S_Q(\mathrm{\Omega }_x,\omega _x,z_m),$$
(9)
where
$`S_Q(\mathrm{\Omega }_x,\omega _x,z_m)`$ $`=`$ $`[(1\mathrm{\Omega }_x)(1+z_m)+`$ (10)
$`\mathrm{\Omega }_x(1+z_m)^{1+3\omega _x}]^{\frac{1}{2}}.`$
Now, combining equations (8)-(10), we find
$$_{(1+z_m)^1}^1\frac{dx}{xE_Q(x)}=S_Q(\mathrm{\Omega }_x,\omega _x,z_m).$$
(11)
The meaning of the above equation is self evident. It represents an implicit integro-algebraic equation for the critical redshift $`z_m`$ as a function of the parameters defining the flat Q-models. In general, this expression cannot be solved in closed analytical form for $`z_m`$. However, by taking the limit $`\mathrm{\Omega }_x=0`$ in (11), the value $`z_m=1.25`$ is readily obtained as should be expected. The interesting point here is that (11) is quite convenient for a numerical treatment. A similar equation can also be derived for an open cold dark matter universe (OM). We find
$$\mathrm{\Delta }^1\mathrm{tanh}\left[\mathrm{\Delta }_{(1+z_m)^1}^1\frac{dx}{xE_{OM}(x)}\right]=F_{OM}(\mathrm{\Omega }_M,z_m),$$
(12)
where $`\mathrm{\Delta }=(1\mathrm{\Omega }_M)^{\frac{1}{2}}`$ and the functions $`E_{OM}`$, $`F_{OM}(\mathrm{\Omega }_M,z_m)`$ are given by
$$E_{OM}(x)=\left[1\mathrm{\Omega }_M+\mathrm{\Omega }_Mx^1\right]^{\frac{1}{2}}$$
(13)
$$F_{OM}(\mathrm{\Omega }_M,z_m)=\left[1\mathrm{\Omega }_M+\mathrm{\Omega }_M(1+z_m)\right]^{\frac{1}{2}}.$$
(14)
In Fig.2 we show the diagrams of $`z_m`$ as a function of the density parameter $`\mathrm{\Omega }_x`$, and some selected values of $`\omega _x`$ (at this point the reader should compare our results with the alternative numerical method developed by Krauss & Schramm (1993) for a $`\mathrm{\Lambda }`$CDM universe). Note that equation (12) has also been used to plot the case for the open universes (solid line). In the former case, the curves show us clearly that all the Q-models belongs to the same class, which contains the case of a pure cosmological constant. The smallest value of the critical redshift is exactly the one given by Einstein-de Sitter universe ($`\mathrm{\Omega }_x=0`$). This value is pushed to the right direction, that is, for any value of $`\omega _x`$ it is displaced to higher redshifts as the $`\mathrm{\Omega }_x`$ parameter increases. For instance, consider that $`\omega _x=0.5`$. By taking $`\mathrm{\Omega }_x=0.5`$ and $`\mathrm{\Omega }_x=0.8`$, we find $`z_m=1.42`$ and $`z_m=1.65`$, respectively. In the opposite extreme ($`\mathrm{\Omega }_x1`$) the critical redshift is finite unless the parameter of the equation of state take the extreme value for a pure $`\mathrm{\Lambda }`$ ($`\omega _x=1`$). Note also that for a given value of $`\mathrm{\Omega }_x`$, the minimum is also displaced for higher redshifts when the $`\omega _x`$ parameter diminishes. However, we see that the effect is small if the density parameter $`\mathrm{\Omega }_x`$ is low, say, smaller than 0.2, since the curves are nearly flat and practically coincid below this limit. This coincidence is even more surprizing for $`\omega _x0.7`$. Note that open models driven by cold dark matter (OM) affects strongly the angular size, however, in a somewhat different manner as compared to what happens in a generic Q-model (see Fig.2). In particular, if the density parameters are small, say, $`\mathrm{\Omega }_M`$, $`\mathrm{\Omega }_x`$, smaller than 1/3, the critical redshift is much bigger in the former than in the later. This is easy to understand physically, because at this limit the contibution of the quintessence is small, leading to results close to the flat Einstein-de Sitter universe.
For a large class of Q-models considered in this letter, the critical redshifts $`z_m`$ are displayed in Table 1. The OM results have also been quoted for comparison. The third column with $`\omega _x=1`$ corresponds to a flat $`\mathrm{\Lambda }`$CDM model (see Krauss & Schramm 1993). From these results one arrives to a inevitable conclusion: even neglecting evolution, the redshift at which the angular size is minimal cannot alone discriminate between world models since different scenarios may provide the same $`z_m`$ values. In particular, for the observationally favoured open universe ($`\mathrm{\Omega }_M=0.3`$) we find $`z_m=1.89`$, a value that may also be obtained for Q-models having $`0.85\mathrm{\Omega }_x0.93`$ and $`1\omega _x0.5`$. However, if the angular diameter results are combinated with other tests, some interesting cosmological constraints may be obtained. For instance, at galactic scales, the observed COBE normalized pattern of density fluctuations is more difficult to fit within a low-density open universe than in Q-models (Caldwell et al. 1998). Another real possibility is that the universe is actually in an accelerated expansion state ($`q_o<0`$), as indicated recently by measurements using Type Ia supernovae (Riess et al. 1998; Perlmutter et al. 1998). In this case, any model of the standard FRW class is ruled out regardless of its curvature parameter. However, the bidimensional parameter space ($`\mathrm{\Omega }_x,\omega _x`$) is still large enough to accomodate Q-models predicting both an acelerated expansion (if $`q_o<\frac{1}{3\mathrm{\Omega }_x}`$), and high values of the critical redshift, say, close to the values of $`z_m`$ given by the open models.
The same analytical procedure developed here may be applied when evolutionary and/or selection effects due to a linear size-redshift or to a luminosity-redshift dependence are taken into account<sup>1</sup><sup>1</sup>1For a more detailed discussion on these effects see Buchalter et al. 1998 . As widely believed, a plausible way of standing for such effects is to consider that the intrinsic linear size has a similar dependence on the redshift as the coordinate dependence, i.e., $`D=D_o(1+z)^c`$, being $`c<0`$ (Ubachukwu 1995; Buchalter et al. 1998). In this case, Eq.(11) is still valid but the function $`S_Q(\mathrm{\Omega }_x,\omega _x,z_m)`$ must be divided by a factor $`(1+c)`$. The displacement of $`z_m`$ relative to the case with no evolution ($`c=0`$) due to the effects above mentioned may be unexpectedly large. For example, if one takes $`c=0.8`$ as found by Buchalter et al. 1998, the redshift of the minimum angular size for the Einstein-de Sitter case ($`\mathrm{\Omega }_x=0`$) moves from $`z_m=1.25`$ to $`z_m=11.25`$. In this way, the minimal is clearly removed for all practical purposes. This result may be a possible explanation why the data of Gurvits et al. (1999), although apparently in agreement with the Einstein-de Sitter universe, do not show clear evidence for a minimal angular size close to $`z=1.25`$, as should be expected for this model. This sort of effect is even greater when an additional “quintessence” component is also considered. In the same vein, since evolution is not forbidden from any principle, we stress that constraints from angular size redshift relation should be taken with some caution.
Acknowledgments: This work was partially supported by the project Pronex/FINEP (No. 41.96.0908.00) and Conselho Nacional de Desenvolvimento Científico e Tecnológico - CNPq (Brazilian Research Agency). |
warning/0003/hep-th0003142.html | ar5iv | text | # Non-supersymmetric cousins of supersymmetric gauge theories: quantum space of parameters and double scaling limits
## a Introduction and overview of the results
In recent years, notable successes have been achieved in understanding strongly coupled quantum field theories and string theories in various space-time dimensions, by combining old heuristic ideas with the power of supersymmetry. Perhaps most striking amongst these are the results on four dimensional, asymptotically free, gauge theories. In , Seiberg and Witten, and Seiberg, were able to compute the exact quantum corrections to the moduli space of vacua of various $`N=1`$ and $`N=2`$ gauge theories, using a subtle generalization of Montonen-Olive electric/magnetic duality valid at low energy. In a concrete proposal was made, in some particular supersymmetric examples, for the long-suspected string description of gauge theories , particularly when the number of colors is large. Unfortunately, these works rely heavily on very special mathematical properties of supersymmetric theories, and it has been impossible so far to assess their relevance to the non-supersymmetric world.
The purpose of this letter is to provide a framework where the relevance of supersymmetric models for non-supersymmetric gauge theories can be precisely studied. The idea is to consider a class of asymptotically free quantum field theories in two space-time dimensions which are distinguished by the fact that the low energy coupling can be changed by varying mass parameters, which thus play the rôle of Higgs expectation values. Versions of these theories with four supercharges show quantitative similarities with $`N=2`$ super Yang-Mills , and are thus the best possible toy models for the four dimensional supersymmetric gauge theories.
In this work, two main results are obtained. First, we will see that non-supersymmetric versions of the models discussed in are tractable as well, and do continue to display qualitatively the same physics. The quantum corrections to the space of mass parameters $``$ (which is the analogue of the moduli space of four dimensional gauge theories) can be computed. Convincing evidence is found that the physics unraveled in supersymmetric gauge theories can be relevant in genuinely non-supersymetric models as well. Second, we will show that one can take a double scaling limit, in the sense of the “old” matrix models , when approaching the singularities on the quantum space of mass parameters. This possibility was not anticipated on the gauge theory side, and, if correct in this context, could potentially be of great theoretical interest to understand the field theory/string theory duality in four dimensions.
Our claims will be exemplified in this letter by studying a simple purely bosonic theory displaying a rich physics akin to what was found in the four dimensional supersymmetric gauge theories. The model is a natural generalization of both the $`\mathrm{O}(N)`$ non-linear sigma model and the sine-Gordon model. The fields parametrize the $`N1`$ dimensional sphere $`\mathrm{S}^{N1}`$, and the lagrangian is the sum of the standard $`\mathrm{O}(N)`$ symmetric kinetic term and an additional interaction term which, for $`N=2`$, reduces to the sine-Gordon potential. Classically, this latter term gives a mass $`m`$ to the would-be Goldstone bosons. Quantum mechanically, when $`m`$ is much larger that the dynamically generated “hadronic” mass scale $`\mathrm{\Lambda }`$, the theory is weakly coupled and much of the physics can be read off from the lagrangian. On the contrary, when $`m\mathrm{\Lambda }`$, strong quantum corrections are expected. When $`m=0`$ we recover the standard $`\mathrm{O}(N)`$ non-linear sigma model, which has a mass gap and a spectrum made of a single particle in the vector representation of $`\mathrm{O}(N)`$ .
The quantum space of parameters $`_\mathrm{q}`$ can be worked out by using various techniques, including a large $`N`$ approximation. The submanifold of singularities $``$, defined to be the locus in $``$ where some of the degrees of freedom are massless, drastically changes when one goes from the classical to the quantum regime. Typically, either $`_{\mathrm{cl}}`$ locally splits into two, or does not change its shape, but in both cases the low energy physics is different in the classical and the quantum theory. Globally we obtain a hypersurface of singularities $`_\mathrm{q}`$ which delimits two regions on $`_\mathrm{q}`$, one at strong coupling and the other extending to weak coupling. On some regions of $`_\mathrm{q}`$, both a soliton, which is in our model a generalization of the sine-Gordon soliton, and a bound state of the elementary fields, become massless. The physics in the infrared is then governed by a non-trivial conformal field theory, either an Ising model or an $`\mathrm{O}(2)`$ symmetric Ashkin-Teller model. Kramers-Wannier duality exchanges the soliton and the bound state, and inside $`_\mathrm{q}`$ the notion of a topological charge is ambiguous. The dictionary with phenomena in four dimensional supersymmetric gauge theories is the following: our non-trivial CFTs are like Argyres-Douglas CFTs , the sine-Gordon soliton corresponds to a ’t Hooft-Polyakov monopole and Kramers-Wannier duality is mapped onto Montonen-Olive duality.
Another aspect of the model is that its $`1/N`$ expansion can be interpreted as a sum over topologies for randomly branched polymers. In particular, for some values of the parameters, the model has an $`\mathrm{O}(N1)`$ symmetry and can be viewed as a $`(N1)`$-vector model with an infinite number of interactions (bonds involving an arbitrary number of molecules in the branched polymers). One can approach the critical surface $`_\mathrm{q}`$ by taking a suitable double scaling limit, which on the one hand gives a description of the model in terms of extended objects (the polymers), and on the other hand defines non-perturbatively a continuous theory of polymers in a fully consistent context. In particular, the model overcomes notorious difficulties with double scaling limits in two dimensions .
## b The model and its semi-classical properties
We will work with a space-time of euclidean signature and write the lagrangian of the model as
$$L=\frac{1}{2}\underset{i=1}{\overset{N}{}}_\mu \varphi _{i,0}_\mu \varphi _{i,0}+\frac{\alpha }{2}\left(\underset{i=1}{\overset{N}{}}\varphi _{i,0}^2\frac{1}{g_0^2}\right)+V_\mathrm{m}.$$
(1)
$`\alpha `$ is a Lagrange multiplier implementing the constraint that the target space is a sphere of radius $`1/g_0^2`$. Without the mass term $`V_\mathrm{m}`$, the theory is made UV finite by simple multiplicative renormalizations of the fields and coupling . In the leading $`1/N`$ approximation, only the coupling constant renormalization is needed, and we will take the renormalized fields $`\varphi _i=\varphi _{i,0}`$, and coupling constant $`g`$ such that
$$\frac{1}{g^2}=\frac{1}{g_0^2}+\frac{N}{2\pi }\mathrm{log}\frac{\mu }{\mathrm{\Lambda }_0}=\frac{N}{2\pi }\mathrm{log}\frac{\mu }{\mathrm{\Lambda }}\text{,}$$
(2)
where $`\mathrm{\Lambda }_0`$ is the UV cut-off, $`\mu `$ is a sliding scale, and $`\mathrm{\Lambda }`$ is the dynamically generated mass scale of the theory. A mass term $`V_\mathrm{m}`$ is characterized by the canonical dimension of the mass parameters and the way they transform under $`\mathrm{O}(N)`$. For example, a magnetic field has canonical dimension 2 and transforms in the vector representation. Once these data are fixed, the explicit form of $`V_\mathrm{m}`$ is deduced from renormalization theory. In our model, the mass parameters will be taken to have canonical dimension 2 and to transform as a symmetric traceless rank two tensor $`h_{ij}`$ (they are like a tensor magnetic field). In general, $`h_{ij}`$ is multiplicatively renormalized; no renormalization is actually needed in the leading $`1/N`$ expansion. By diagonalizing $`h_{ij}`$ we can write
$$V_\mathrm{m}=\frac{1}{2}\underset{i=1}{\overset{N}{}}h_i\varphi _i^2.$$
(3)
The trace part of $`h_{ij}`$ would correspond to a constant term in the lagrangian, and can thus be taken to be non zero without affecting the physics. We will use the $`N1`$ independent dimensionless physical parameters $`v_i=(h_Nh_i)/\mathrm{\Lambda }^2`$ and $``$ will be the $`v`$-space. Using the permutation symmetry amongst the $`h_i`$s, we will restrict ourselves to the region $`^+`$ of positive $`v_i`$s unless explicitly stated otherwise. The classical masses of the $`N1`$ independent elementary fields $`\varphi _i`$, $`1iN1`$, are then $`m_i=\mathrm{\Lambda }\sqrt{v_i}`$.
The model has always $`N`$ $`_{2(i)}`$ symmetries $`\varphi _i\varphi _i`$, and can also have additional $`\mathrm{O}(p)`$ symmetries when $`p`$ of the $`v_i`$s coincide.
Singularities on $`^+`$ are found classically when some, say $`p`$, of the $`v_i`$s vanish. The low energy physics is then governed by a standard $`\mathrm{O}(p+1)`$ non linear sigma model, and the $`p`$ massless states are the $`p`$ classical Goldstone bosons for the breaking of $`\mathrm{O}(p+1)`$ down to $`\mathrm{O}(p)`$. $`_{\mathrm{cl}}`$ in $`^+`$ thus coincide with the hyperplanes $`v_i=0`$.
Quantum mechanically, in the weakly coupled region $`v_i1`$, we can use semiclassical techniques to investigate the spectrum of particles further. It can be shown that a bound state $`\varphi `$-$`\varphi `$ is associated with the operator $`\varphi _N^2=1/g_0^2_{i=1}^{N1}\varphi _i^2`$. We will compute the mass of this state, for any $`v_i`$s, in the next section. One can also show that the model admits solitons connecting the two degenerate minima of the potential (3) at $`\varphi _N=\pm 1/g^2`$. These two minima are related by the spontaneously broken $`_{2(N)}`$ symmetry. All the solitonic (time independent, finite energy) solutions can be explicitly found . In the simple $`\mathrm{O}(N1)`$ symmetric case where $`v_1=\mathrm{}=v_{N1}=v1`$, they correspond to trajectories joining the two poles at $`\varphi _N=\pm 1/g^2`$ along a meridian of the target space sphere. They are standard sine-Gordon solitons of masses $`M_{\mathrm{cl}}=2\mathrm{\Lambda }\sqrt{v}/g^2`$. The semi-classical quantization shows that the solitons are particles filling multiplets of $`\mathrm{O}(N1)`$ corresponding to the completely symmetric traceless tensor representations. The rank $`J`$ tensor has a mass $`M_J=M_{\mathrm{cl}}+J(J+N3)v\mathrm{\Lambda }^2/(2M_{\mathrm{cl}})`$.
## c The large N approximation
A useful technique to study our model is to use a large $`N`$ approximation (for a recent review, see ), where $`N`$ is sent to infinity while the scale $`\mathrm{\Lambda }`$ is held fixed . The large $`N`$ expansion is nothing but a standard loop expansion for a non-local effective action $`S_{\mathrm{eff}}`$ obtained by integrating exactly a large number of elementary fields $`\varphi _i`$ from (1). For our purposes, it will be useful to keep explicitly the order parameter $`\varphi _N=\sqrt{N}\phi `$ for $`_{2(N)}`$ in the action. The effective action for large $`N`$ is then
$`{\displaystyle \frac{S_{\mathrm{eff}}}{N}}={\displaystyle d^2x\left\{\frac{1}{2}_\mu \phi _\mu \phi +\frac{\alpha h_N}{2}\phi ^2\frac{\alpha }{2Ng^2}\right\}}+{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N1}{}}}s\left[\alpha h_i\right],`$ (4)
where $`s[f]=(1/2)tr\mathrm{ln}(^2+f)\mathrm{ln}(\mathrm{\Lambda }_0/\mu )d^2xf/(4\pi )`$ can be expanded in terms of ordinary Feynman diagrams.
It is useful to compute, in this framework, the masses of the particles we found previously in a semi-classical approximation. At leading order, this is done by looking at poles in the two-point functions derived from (4). In the $`\mathrm{O}(N1)`$ symmetric case $`v_1=\mathrm{}=v_{N1}=v`$, and for $`v>1`$, the $`N1`$ elementary fields $`\varphi _i`$ have a mass $`m_\varphi =\mathrm{\Lambda }\sqrt{v}`$, and the mass $`m_\mathrm{b}`$ of the $`\varphi `$-$`\varphi `$ bound state (the field $`\phi `$) is a solution of
$$\sqrt{4\mathrm{\Lambda }^2v/m_\mathrm{b}^21}\mathrm{ln}v=2\mathrm{arctan}\left(1/\sqrt{4\mathrm{\Lambda }^2v/m_\mathrm{b}^21}\right).$$
(5)
$`m_\mathrm{b}`$ is a monotonic function of $`v`$, decreasing from $`m_\mathrm{b}2m_\varphi \left(1(Ng^2)^2/32\right)`$ for $`v1`$ to $`m_\mathrm{b}=0`$ for $`v=1`$. The mass of the $`\mathrm{O}(N1)`$ singlet soliton also goes to zero as $`v`$ goes to one because the two degenerate minima of the effective potential derived from $`S_{\mathrm{eff}}`$ merge at $`v=1`$. More generally, the $`\varphi `$-$`\varphi `$ bound state and the lightest solitonic state will become massless together on the critical hypersurface $`_{i=1}^{N1}v_i=1`$. Near this hyperboloid, the $`1/N`$ expansion has IR divergencies and is no longer reliable. These divergencies are due to the fact that we are near a critical point below the critical dimension. To describe the physics near the critical surface, we must go beyond the $`1/N`$ approximation and sum the most relevant (i.e. the most IR divergent) contributions. This also automatically resolves the difficulties associated with massless propagators in two dimensions. One can show that the low energy effective lagrangian on the critical surface is, in the large $`N`$ limit,
$$L_{\mathrm{eff}}=\frac{1}{2}_\mu \varphi _N_\mu \varphi _N+\frac{\pi \mathrm{\Lambda }^2}{_{i=1}^{N1}1/v_i}\varphi _N^4,$$
(6)
where the interaction, which is proportional to $`1/N`$, must be treated exactly, since the IR divergencies compensate for the $`1/N`$ factors. We thus obtain the Landau-Ginzburg description of an Ising critical point. The field $`\varphi _N`$ is the order operator and the massless soliton corresponds to the disorder operator; they are exchanged by Kramers-Wannier duality. Displacements on the critical surface are associated here with irrelevant operators in the IR.
It is important to note that the $`1/N`$ corrections to the equation of the critical surface itself only suffer from mild logarithmic IR divergencies that can be handled . The existence and form of this surface can thus be reliably studied in a $`1/N`$ expansion. We will see that the critical surface actually intersects with the hyperplanes $`v_i=0`$, and thus join with the other sheet of the surface existing for $`v_i<0`$. The relevant physics is discussed in the following section.
## d The quantum space of parameters
There are regions in the space of parameters where it is easy to see that the classical and quantum hypersurfaces of singularities coincide. When one of the $`v`$s is zero, say $`v_1=0`$, and all the other $`v`$s are large, the low energy theory is an $`\mathrm{O}(2)`$ sigma model of small effective coupling $`g_{\mathrm{eff}}`$. Both classically and quantum mechanically, this is a massless theory, and thus $`_{\mathrm{cl}}=_\mathrm{q}=\left\{v_1=0\right\}`$ in this region. However, the physics are not the same: the classical theory of a free massless boson is replaced in the quantum case by the non-trivial CFT of a boson compactified on a circle of radius $`R`$ with $`R^2=4\pi /g_{\mathrm{eff}}^2=\mathrm{ln}(_{i=2}^{N1}v_i)`$. If $`p`$ of the $`v_i`$s, $`i2`$, decrease, while we stay on the hyperplane $`v_1=0`$, the low energy theory will tend to become an $`\mathrm{O}(p+2)`$ non linear sigma model and thus develop a mass gap. In the large $`N`$ limit, it can be shown that this transition takes place on a surface $`𝔥_{\mathrm{q},1}`$ whose equation is $`_{i=2}^{N1}v_i=1`$, $`v_1=0`$, and that the low energy theory on $`𝔥_{\mathrm{q},1}`$ is an $`\mathrm{O}(2)`$ symmetric Ashkin-Teller model with Landau-Ginzburg potential
$$V_{\mathrm{eff}}=\frac{\pi \mathrm{\Lambda }^2}{_{i=2}^{N1}1/v_i}\left(\varphi _1^2+\varphi _N^2\right)^2.$$
(7)
This model is indeed equivalent to a compactified boson for the particular radius $`R=R_{\mathrm{KT}}=2\sqrt{2}`$, and the transition through $`𝔥_{\mathrm{q},1}`$ is nothing but a Kosterlitz-Thouless phase transition. A local analysis shows that $`_\mathrm{q}`$ joins the hyperplane $`v_1=0`$ orthogonally along $`𝔥_{\mathrm{q},1}`$, and that going from $`v_1=0`$ to $`v_1>0`$ on $`_\mathrm{q}`$ is equivalent to turning on a relevant operator which decouples one of the two Ising spins in the Ashkin-Teller model. The other spin would decouple in the $`v_1<0`$ region. We see that the Ashkin-Teller model on $`𝔥_{\mathrm{q},1}`$ is made of the coupling of the two Ising models (one for $`v_1>0`$, the other for $`v_1<0`$) found in the preceding section.
It is possible to find a simple equation for $`_\mathrm{q}`$ valid for all values of the $`v_i`$s, positive or negative, in the large $`N`$ limit:
$$_\mathrm{q}:\underset{i=1}{\overset{N}{}}\underset{ji}{}(rh_j)=\mathrm{\Lambda }^{2(N1)},$$
(8)
where $`r`$ is the largest real root of the polynomial $`_{i=1}^N(xh_i)\stackrel{~}{\mathrm{\Lambda }}^{2N}`$. $`\stackrel{~}{\mathrm{\Lambda }}^{2N}`$ can in principle be determined in terms of $`\mathrm{\Lambda }`$ by a non-perturbative calculation. The fact that the Kosterlitz-Thouless transition occurs at a strictly positive radius $`R_{\mathrm{KT}}`$ implies $`\stackrel{~}{\mathrm{\Lambda }}<\mathrm{\Lambda }`$. The shape of $`_\mathrm{q}`$ is likely to be qualitatively reproduced by (8) even for small values of $`N`$, and we have used it to draw Fig. 1, which gives a global picture of the quantum space of parameters of our model in the case $`N=4`$.
## e The double scaling limits
On the critical surface $`_\mathrm{q}`$, the sum of Feynman diagrams of a given topology (i.e. contributing at a fixed order in $`1/N`$) diverges. It is then natural to ask whether it is possible to approach $`_\mathrm{q}`$ and take the limit $`N\mathrm{}`$ in a correlated way in order to obtain a finite answer taking into account diagrams of all topologies. This is the idea of the double scaling limit , and for vector models the result can be interpreted as giving the partition function of a continuous theory of randomly branched polymers .
Let us consider for example the case of the $`\mathrm{O}(N1)`$ symmetric theory $`v_1=\mathrm{}=v_{N1}=v`$. Near the critical point $`v=1`$, the effective theory is just given by (6) with the relevant perturbation $`(1/2)\mathrm{\Lambda }^2(1v)\mathrm{\Phi }_N^2`$. The idea (used in ) is then to eliminate the explicit $`N`$ dependence in the interaction term by using a rescaled space-time variable $`y=x/\sqrt{N}`$,
$`S_{\mathrm{eff}}={\displaystyle d^2y\left\{\frac{1}{2}_\mu ^y\varphi _N_\mu ^y\varphi _N+\frac{1}{2}\mathrm{\Lambda }^2N(1v)\varphi _N^2+\pi \mathrm{\Lambda }^2\varphi _N^4\right\}}.`$ (9)
A consistent double scaling limit can be defined to be $`N\mathrm{}`$, $`v1`$, with $`N(v1)3\mathrm{ln}N`$ kept fixed. The logarithmic correction to the naive scaling comes from the fact that the large $`N`$ limit is also a large UV cutoff limit in the $`y`$ variable, and (9) needs to be renormalized, which is done by a simple normal ordering. A similar non-trivial double scaling limit can be defined near the Ashkin-Teller point. For example, for $`N\mathrm{}`$, $`v_10`$, $`v_2=\mathrm{}=v_{N1}=v1`$, the combinations that must be kept fixed are $`N(v1)4\mathrm{ln}N`$ and $`N(vv_11)4\mathrm{ln}N`$. Note that our theory in these limits is free of the inconsistencies found in the standard vector models in a similar context , as can be shown from a straightforward calculation of the 1PI effective action.
## f Other models
It is possible to study other mass terms or/and other target spaces along the lines of the present work. It could also be interesting to perform lattice calculations for this class of models. For example, a mass term $`m_{ij}=m_{ji}`$ of canonical dimension one allows to obtain higher critical points. The quantum space of parameters of a $`\mathrm{P}^N`$ model with mass terms, which has instantons, a $`\theta `$ angle, and exhibit confinement at strong coupling, is also likely to display a rich structure. Finally, it is natural to consider supersymmetric versions of our models. It turns out that the massive version of the supersymmetric $`\mathrm{P}^N`$ model shows quantitative similarities with $`N=2`$ super Yang-Mills . A discussion of these models along with details on the present work will be published elsewhere .
I would like to thank Princeton University for offering me matchless working conditions. |
warning/0003/cond-mat0003511.html | ar5iv | text | # Elastic property of single double-stranded DNA molecules: Theoretical study and comparison with experiments
## I Introduction
DNA molecule is the primary genetic material of most organisms. It is a double-helical biopolymer in which two chains of complementary nucleotides (the subunits whose sequence constitutes the genetic message) wind (usually right-handedly) around a common axis to form a double-helical structure . Because of this unique structure, the elastic property of DNA molecule influences its biological functions greatly. There are mainly three kinds of deformations in DNA double-helix: stretching and bending of the molecule, twisting of one nucleotide chain relative to its counterpart. All these deformations have vital biological significance. During DNA replication, hydrogen bonds between the complementary DNA bases should be broken and the two nucleotide chains be separated. This strand-separation process requires cooperative unwinding of the double-helix . In DNA recombination reaction, RecA proteins polymerize along DNA template and the DNA molecule is stretched to $`1.5`$ times its relaxed contour length . It is suspected that thermal fluctuations of DNA central axis might be very important for RecA polymerization . Another important example is the process of chromosome condensation during prophase of the cell cycle, where the long (circular) DNA chain wraps tightly onto histone proteins and is severely bent . Further more, in living cells DNA chain is usually closed, i.e., the two ends of the molecule is linked together by covalent bonds and the molecule becomes endless. With this chain-closing process, all those quantities characterizing the topological state of the chain are fixed and can only be changed externally by topoisomerases, which are capable of transiently cutting one or both DNA strands and making one strand pass through the other at the cutting point (in the case of type I topoisomerases ) or one segment of DNA pass through another (in the case of type II topoisomerases ). It is possible that this kind of enzyme-induced topology-changing processes are also closely related to the particular mechanic property of DNA molecule. For example, the frequency of collisions between two distant DNA segments in a circular DNA molecule is influenced by the different knot types and different linking numbers (for a definition of this quantity, see below and Sec. II C). A thorough investigation of the deformation and elasticity of DNA will enable us to gain better understanding on many important biological processes concerned with life and growth.
Detailed study on DNA elasticity now becomes possible with the recent experimental developments, including, e.g., optical tweezer methods, atomic force microscopy, fluorescence microscopy. These techniques make it possible to manipulate directly single polymeric molecules and to record their elastic responses with high precision. Experiments done on double-stranded DNA (dsDNA) have revealed that this molecule has very novel elastic property . When a torsionally relaxed DNA is pulled with a force less than $`10`$ picoNewton (pN), its elastic response can be quantitatively understood by regarding the chain as an inextensible thin string with certain bending rigidity (namely, the wormlike chain model ). However, if the external force is increased up to $`65`$ pN, DNA chain becomes highly extensible. At this force, the molecule transit to an over-stretched configuration termed S-DNA, which is $`1.6`$ times longer than the same molecule in its standard B-form structure . Besides external forces, it is also possible to apply torsional constraints to DNA double-helix by external torques. The linking number of DNA, i.e., the total topological turns one DNA strand winds around the other, can be fixed at a value larger (less) than the molecule’s relaxed value. In such cases we say the DNA molecule is positively (negatively) supercoiled. It is shown experimentally that when external force is less than a threshold value of about $`0.3`$ pN, the extension of DNA molecule decreases with increasing twisted stress and the elastic response of positively supercoiled DNA is similar to that of negatively supercoiled DNA, indicating the DNA chain might be regarded as achiral. However, if the external force is increased to be larger than this threshold, negatively and positively supercoiled DNA molecules behave quite differently. Under the condition of fixed external force between $`0.3`$ pN and $`3`$ pN, while positive twist stress keeps shrinking the DNA polymer, the extension of negatively supercoiled DNA is insensitive to supercoiling degree . In higher force region, it is suggested by some authors that positively supercoiled DNA may transit to a configuration called Pauling-like DNA (P-DNA) with exposed nucleotide bases , while negative torque may lead to strand-separation in DNA molecule (denaturation of DNA double-helix ). A very recent systematic observation performed by L$`\stackrel{´}{\mathrm{e}}`$ger et al. , on the other hand, suggested another possibility that negative supercoiling may result in left-handed Z-form configuration in DNA.
The above-mentioned complicated elastic property revealed by experiments may be directly related to the versatile roles played by DNA molecule in living organisms. Theoretically, to understand DNA elastic property is of current interest. Concerned with one or another aspect of DNA elasticity, models were proposed and valuable insights were obtained (see, for example, Refs. ), and now it is widely accepted that the competition between DNA bending and torsional deformations deserves to pay considerable attention in understanding the elastic property of DNA molecules. However, it is still a great challenge to understand systematically and quantitatively all aspects of DNA mechanical property based on the same unified framework. What is the intrinsic reason for DNA molecule’s entropic elasticity, highly extensibility as well as its supercoiling property? Is it possible for negative torque to stabilize left-handed DNA configurations? These are just some examples of unsolved questions.
In the present work, we have tried to obtain a comprehensive and quantitative understanding on DNA mechanical property. We have thought that the double-stranded nature of DNA structure should be extremely important to its elastic property, and therefore have constructed a general elastic model in which this characteristic is properly taken into account via the introduction of a new structural parameter, the folding angle $`\phi `$. The elastic property of long dsDNA molecules was then studied based on this model, where the base-stacking interactions between DNA adjacent nucleotide basepairs, their steric effects, and the electrostatic interactions along DNA backbones were all considered. Quantitative results were obtained by using path integral method, and excellent agreement between theory and the experimental observations of several groups were attained. It was revealed that, on one hand, the strong intensity of the base-stacking interactions ensures the structural stability of DNA molecule; while on the other hand, the short-ranged nature of such interactions makes externally-stimulated large structural fluctuations possible. The entropic elasticity, highly extensibility, and supercoiling property of dsDNA molecule are all closely related to this fact. The present work also revealed the possibility that negative torque can induce structural transitions in highly extended DNA from right-handed B-form configuration to left-handed Z-form-like configurations. Some discussions on this respect were performed and we suggested that a possible direct way to check the validity of this opinion is to measure the values of the critical torques under which such transitions are anticipated to take place by the present calculations.
This paper is organized as follows: In Sec. II we introduce the elastic model for dsDNA biopolymers. At the action of an external force, the elastic response of dsDNA molecules is investigated in Sec. III and compared with experimental observations of Smith et al. and Cluzel et al. . Section IV focuses on supercoiled dsDNA molecules, where the relationship between extension and supercoiling degree is obtained numerically and compared with the experiment of Strick et al. . From the calculated folding angle distribution, we infer that negative torque can cause structural transitions in dsDNA molecules from right-handed double-helix to left-handed ones. Section V is reserved for conclusion. Two appendices are also presented: in Appendix A we review some basic ideas on the application of path integral method in polymer physics; and in Appendix B we list the matrix elements of the operators in Eq. (18) and Eq. (27). Some parts of this work have been briefly reported in a previous letter.
## II Elastic model of double-stranded DNA molecule
As already stressed, DNA molecule is a double-stranded biopolymer. Its two complementary sugar-phosphate chains twist around each other to form a right-handed double-helix. Each chain is a linear polynucleotide consisting of the following four bases: two purines (A, G) and two pyrimidines (C, T) . The two chains are joined together by hydrogen bonds between pairs of nucleotides A-T and G-C. Hereafter, we refer to the two sugar-phosphate chains as the backbones and the hydrogen-bonded pairs of nucleotides as the basepairs. In this section we discuss the energetics of such an elastic system (see Fig. 1). First of all, the bending energy of the backbones of such double-stranded polymers will be considered; then we will discuss the interactions between DNA basepairs and energy terms related with external fields.
### A Bending and folding deformations
The backbones can be regarded as two inextensible wormlike chains characterized with a very small bending rigidity $`\kappa =k_BT\mathrm{}_p`$, where $`k_B`$ is Boltzmann’s constant and $`T`$ the environmental temperature, and $`\mathrm{}_p1.5`$ nm is the bending persistence length of single-stranded DNA (ssDNA) chains . The bending energy of each backbone is thus expressed as $`\kappa _0^L(d𝐭_i/ds)^2𝑑s`$, where $`𝐭_i(s)`$ $`(i=1,2)`$ is the unit tangent vector at arclength $`s`$ along the $`i`$-th backbone , and $`L`$ is the total contour length of each backbone. The position vectors of the two backbones are expressed as $`𝐫_i(s)=^s𝐭_i(s^{})𝑑s^{}`$.
Since there are many relatively rigid basepairs between the two backbones in many cases the lateral distance between the backbones can be regarded to be constant and equal to $`2R`$.<sup>*</sup><sup>*</sup>*In our present work, we have not taken into account the possible deformations of the nucleotide basepairs. In many cases this may be a reasonable assumption. However, under some extreme conditions such kind of deformations may turn to be important. For example, when DNA double-helix is stretched and at the same time is applied with a large positive torque, the nucleotide basepairs may collapse (see Ref. for a description). In this subsection we focus on the bending energy of the backbones, therefore, for the moment we regard each basepair as a rigid rod of length $`2R`$ linking between the two backbones and pointing along direction denoted by a unit vector $`𝐛`$ from $`𝐫_1`$ to $`𝐫_2`$ (Fig. 1). Then, $`𝐫_2(s)𝐫_1(s)=2R𝐛(s)`$. In B-form DNA the basepair plane is perpendicular to DNA axis, therefore, in our model relative sliding of the two backbones is not considered and the basepair rod is thought to be perpendicular to both backbones , with $`𝐛(s)𝐭_1(s)=𝐛(s)𝐭_2(s)0`$. The central axis of the double-stranded polymer can be defined as $`𝐫(s)=𝐫_1(s)+R𝐛(s)`$ $`[=𝐫_2(s)R𝐛(s)=(𝐫_1(s)+𝐫_2(s))/2]`$, and its tangent vector is denoted by $`𝐭`$. In consistence with actual DNA structures, the central axial tangent $`𝐭`$ is also perpendicular to $`𝐛`$, i.e., $`𝐛(s)𝐭(s)=0`$. (Notice that, however, $`𝐭(s)d𝐫/ds`$; in this paper, $`s`$ always refers to the arclength of the backbones.)
Since all the tangent vectors $`𝐭_1`$, $`𝐭_2`$, and $`𝐭`$ lie on the same plane perpendicular to $`𝐛`$, we can write that
$$\begin{array}{c}𝐭_1(s)=𝐭(s)\mathrm{cos}\phi (s)+𝐧(s)\mathrm{sin}\phi (s),\hfill \\ 𝐭_2(s)=𝐭(s)\mathrm{cos}\phi (s)𝐧(s)\mathrm{sin}\phi (s),\hfill \end{array}$$
(1)
where $`𝐧`$ is also a unit vector and $`𝐧=𝐛\times 𝐭`$, and $`\phi `$ is defined as half the rotational angle from $`𝐭_2`$ to $`𝐭_1`$, with $`𝐛`$ being the rotational axis (Fig. 1). We call $`\phi `$ the folding angle, and it can vary in the range $`(\pi /2,+\pi /2)`$, with $`\phi >0`$ corresponding to right-handed rotations and hence right-handed double-helical configurations and $`\phi <0`$ to left-handed ones. With the help of Eq. (1), we know that
$$\frac{d𝐛}{ds}=\frac{𝐭_2𝐭_1}{2R}=\frac{\mathrm{sin}\phi }{R}𝐧,$$
(2)
and
$$\frac{d𝐫}{ds}=\frac{1}{2}(𝐭_1+𝐭_2)=𝐭\mathrm{cos}\phi .$$
(3)
Equation (3) indicates that $`\mathrm{cos}\phi `$ measures the extent to which the backbones are $`\mathrm{`}\mathrm{`}`$folded” with respect to the central axis. Based on Eqs. (1-3), the total bending energy of the two backbones can be expressed in the following form:
$`E_b`$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}{\displaystyle _0^L}({\displaystyle \frac{d𝐭_1}{ds}})^2𝑑s+{\displaystyle \frac{\kappa }{2}}{\displaystyle _0^L}({\displaystyle \frac{d𝐭_2}{ds}})^2𝑑s`$ (4)
$`=`$ $`\kappa {\displaystyle _0^L}𝑑s\left[({\displaystyle \frac{d𝐭}{ds}})^2\mathrm{cos}^2\phi +\mathrm{sin}^2\phi ({\displaystyle \frac{d𝐧}{ds}})^2+({\displaystyle \frac{d\phi }{ds}})^2\right]`$ (5)
$`=`$ $`{\displaystyle _0^L}𝑑s\left[\kappa ({\displaystyle \frac{d𝐭}{ds}})^2+\kappa ({\displaystyle \frac{d\phi }{ds}})^2+{\displaystyle \frac{\kappa }{R^2}}\mathrm{sin}^4\phi \right].`$ (6)
The bending energy is thus decomposed into the bending energy of the central axis (the first term of Eq. (4)) plus the folding energy of the backbones (the second and third terms of Eq. (4)). The physical meanings of these two energy contributions are very clear, and Eq. (4) is very helpful for our following calculations. In Eq. (4), the bending energy of the central axis is very similar with that of a wormlike chain , both of which are related to the square of the changing rate of the axial tangent vectors. But there are two important differences: (a) in the derivative $`d𝐭/ds`$ of Eq. (4), the arclength parameter $`s`$ is measured along the backbone, not along the central axis; (b) in the wormlike chain model the central axis is inextensible, while here the central axis is extensible.
In deriving Eq. (4), the basepairs are models just as thin rigid rods of fixed length, i.e., the DNA molecule is viewed as a ladder-like structure (see Fig. 1). Actually, however, basepairs form disc-like structures and have finite volume. The steric effects caused by the finite volume of basepairs is anticipated to hinder considerably the bending deformation of the central axis, and hence will increases its bending rigidity greatly . Furthermore, dsDNA molecule is a strong polyelectrolyte, with negatively-charged groups distributed regularly along the chain’s surface. The electrostatic repulsion force between these negatively-charged groups will also considerably increase the bending rigidity of the dsDNA chain . To quantitatively take into account the above mentioned two kinds of effects is very difficult. Here we treat this problem phenomenologically by simply replacing the bending rigidity $`\kappa `$ in the first term of Eq. (4) with a quantity $`\kappa ^{}`$. It is required that $`\kappa ^{}>\kappa `$, and the precise value of $`\kappa ^{}`$ will then be determined self-consistently by the best fitting with experimental data as shown in Sec. III.
### B Base-stacking interactions between basepairs
In Sec. II A, we have discussed in detail the bending energy of dsDNA polymers, which is caused by bending of the backbones as well as steric effects and electrostatic interactions. In dsDNA molecule there is another kind of important interactions, namely the base-stacking interaction between adjacent nucleotide basepairs . The base-stacking interactions originate from the weak van der Waals attraction between the polar groups in adjacent nucleotide basepairs. Such interactions are short-ranged and their total effect is usually described by a potential energy of the Lennard-Jones form ($`6`$-$`12`$ potential ). Base-stacking interactions play significant role in stabilization of DNA double-helix. The main reason why DNA can but RNA can not form long double-helix is as follows: Because of the steric interference caused by the hydroxyl group attached to the $`2^{}`$ carbon of RNA riboses, the stacking interaction between adjacent RNA nucleotide basepairs is very weak and can not stabilize the formed double-helical structure; while in the DNA ribose, it is a hydrogen atom attached to its $`2^{}`$ carbon and serious steric interference is avoided (fortunately!).
In a continuum theory of elasticity, the summed total base-stacking potential energy is converted into the form of the following integration:
$$E_{LJ}=\underset{i=1}{\overset{N1}{}}U_{i,i+1}=_0^L\rho (\phi )𝑑s,$$
(7)
where $`U_{i,i+1}`$ is the base-stacking potential between the $`i`$-th and the $`(i+1)`$-th basepair, $`N`$ is the total number of basepairs, and the base-stacking energy density $`\rho `$ is expressed as
$$\rho (\phi )=\{\begin{array}{cccc}\frac{ϵ}{r_0}[(\frac{\mathrm{cos}\phi _0}{\mathrm{cos}\phi })^{12}\hfill & & 2(\frac{\mathrm{cos}\phi _0}{\mathrm{cos}\phi })^6]\hfill & (\mathrm{for}\phi 0),\hfill \\ \frac{ϵ}{r_0}[\mathrm{cos}^{12}\phi _0\hfill & & 2\mathrm{cos}^6\phi _0]\hfill & (\mathrm{for}\phi <0).\hfill \end{array}$$
(8)
In Eq. (8), the parameter $`r_0`$ is the backbone arclength between adjacent bases ($`r_0=L/N`$); $`\phi _0`$ is a parameter related to the equilibrium distance between a DNA dimer ($`r_0\mathrm{cos}\phi _03.4`$ $`\AA `$); and $`ϵ`$ is the base-stacking intensity which is generally base-sequence specific . In this paper we focus on macroscopic properties of long DNA chains composed of relatively random sequences, therefore we just consider $`ϵ`$ in the average sense and take it as a constant, with $`ϵ`$ $`14.0`$ $`k_BT`$ as averaged over quantum-mechanically calculated results on all the different DNA dimers .
The asymmetric base-stacking potential Eq. (8) ensures a relaxed DNA to take on a right-handed double-helix configuration (i.e., the B-form) with its folding angle $`\phi \phi _0`$. To deviate the local configuration of DNA considerably from its B-form generally requires a free energy of the order of $`ϵ`$ per basepair. Thus, DNA molecule will be very stable under normal physiological conditions and thermal energy can only make it fluctuate very slightly around its equilibrium configuration, since $`ϵkT`$. Nevertheless, although the stacking intensity $`ϵ`$ in dsDNA is very strong compared with thermal energy, the base-stacking interaction by its nature is short-ranged and hence sensitive to the distance between the adjacent basepairs. If dsDNA chain is stretched by large external forces, which cause the average inter-basepair distance to exceed some threshold value determined intrinsically by the molecule, the restoring force provided by the base-stacking interactions will no longer be able to offset the external forces. Consequently, it will be possible that the B-form configuration of dsDNA will collapse and the chain will turn to be highly extensible. Thus, on one hand, the strong base-stacking interaction ensures the standard B-form configuration to be very stable upon thermal fluctuations and small external forces (this is required for the biological functions of DNA molecule to be properly fulfilled ); but on the other hand, its short-rangedness gives it considerable latitude to change its configuration to adapt to possible severe environments (otherwise, the chain may be pulled break by external forces, for example, during DNA segregation ). This property of DNA base-stacking interactions is very important to dsDNA molecule. As we will see in Secs. III and IV, the mechanical property of DNA chain is indeed closely related to the above-mentioned insight.
### C External forces and torques
In the previous two subsections, we have described the intrinsic energy of DNA double-helix. Experimentally, to probe the elastic response of linear DNA molecule, the polymer chain is often pulled by external force fields and/or untwisted or overtwisted by external torques. To study the mechanic response of dsDNA molecule, we consider in this subsection the energy terms related to external forces and torques in our theoretical framework.
For the external force fields, here we constrained ourselves to the simplest situation where one terminal of DNA molecule is fixed and the other terminal is pulled with a force $`𝐅=f𝐳_0`$ along direction of unit vector $`𝐳_0`$ . (In fact, hydrodynamic fields or electric fields are also frequently used to stretch semiflexible polymers , but we will not discuss such cases in this paper.) The end-to-end vector of a DNA chain is expressed as $`_0^L𝐭(s)\mathrm{cos}\phi (s)𝑑s`$ according to Eq. (3). Then the total $`\mathrm{`}\mathrm{`}`$potential” energy of the chain in the external force field is
$$E_f=_0^L𝐭\mathrm{cos}\phi ds𝐅=_0^Lf𝐭𝐳_0\mathrm{cos}\phi ds.$$
(9)
In the experimental setup, external torques can be applied on linear dsDNA molecule by the following procedure: first, DNA ligases are used to ligate all the possible single-stranded nicks; then the two strands of dsDNA molecule at one end are fixed onto a template, while the two strands at the other end is attached tightly to a magnetic bead; afterwards, torques are introduced into DNA double-helix by rotating the magnet bead with an external magnetic field . Torque energy is then related to the topological turns caused by the external torque on DNA double-helix. The total number of topological turns one DNA strand winds around the other, which is usually termed the total linking number $`Lk`$ , is expressed as the sum of the twisting number, $`Tw(𝐫_1(𝐫_2),𝐫)`$ of backbone $`𝐫_1`$ (or $`𝐫_2`$) around the central axis $`𝐫`$ and the writhing number $`Wr(𝐫)`$ of the central axis; i.e., $`Lk=Tw+Wr`$. According to Refs. and Eq. (2), we obtain that
$$Tw(𝐫_1,𝐫)=\frac{1}{2\pi }_0^L𝐭\times 𝐛\frac{d𝐛}{ds}𝑑s=\frac{1}{2\pi }_0^L\frac{\mathrm{sin}\phi }{R}𝑑s.$$
(10)
The writhing number of the central axis is generally much more difficult to calculate. It is expressed as the following Gauss integral over the central axis :
$$Wr(𝐫)=\frac{1}{4\pi }\frac{d𝐫\times d𝐫^{}(𝐫𝐫^{})}{|𝐫𝐫^{}|^3}.$$
(11)
In the case of linear chains, provided that some fixed direction (for example the direction of the external force, $`𝐳_0`$) can be specified and that the tangent vector $`𝐭`$ never points to $`𝐳_0`$ (i.e., $`𝐭𝐳_01`$), it was proved by Fuller that the writhing number Eq. (11) can be calculated alternatively according to the following formula :
$$Wr(𝐫)=\frac{1}{2\pi }\underset{0}{\overset{L}{}}\frac{𝐳_0\times 𝐭d𝐭/ds}{1+𝐳_0𝐭}𝑑s.$$
(12)
The above equation can be further simplified for highly extended linear DNA chains whose tangent $`𝐭`$ fluctuates only slightly around $`𝐳_0`$. In this case, Eq. (12) leads to the approximate expression that
$$Wr(𝐫)\frac{1}{4\pi }_0^L\left[t_y\frac{dt_x}{ds}+t_x\frac{dt_y}{ds}\right]𝑑s,$$
(13)
where $`t_x`$ and $`t_y`$ are, respectively, the two components of $`𝐭`$ with respect to two arbitrarily chosen orthonormal directions ($`𝐱_0`$ and $`𝐲_0`$) on the plane perpendicular to $`𝐳_0`$.
The energy caused by the external torque of magnitude $`\mathrm{\Gamma }`$ is then equal to
$$E_t=2\pi \mathrm{\Gamma }Lk=2\pi \mathrm{\Gamma }(Tw+Wr).$$
(14)
To conclude this section, the total energy of a dsDNA molecule under the action of an external force and an external torque is expressed as
$`E`$ $`=`$ $`E_b+E_{LJ}+E_f+E_t`$ (15)
$`=`$ $`{\displaystyle _0^L}\left[\kappa ^{}({\displaystyle \frac{d𝐭}{ds}})^2+\kappa ({\displaystyle \frac{d\phi }{ds}})^2+{\displaystyle \frac{\kappa }{R^2}}\mathrm{sin}^4\phi +\rho (\phi )f𝐭𝐳_0\mathrm{cos}\phi {\displaystyle \frac{\mathrm{\Gamma }}{R}}\mathrm{sin}\phi \mathrm{\Gamma }{\displaystyle \frac{𝐳_0\times 𝐭d𝐭/ds}{1+𝐳_0𝐭}}\right]𝑑s`$ (16)
$``$ $`{\displaystyle _0^L}\left[\kappa ^{}({\displaystyle \frac{d𝐭}{ds}})^2+\kappa ({\displaystyle \frac{d\phi }{ds}})^2+{\displaystyle \frac{\kappa }{R^2}}\mathrm{sin}^4\phi +\rho (\phi )f𝐭𝐳_0\mathrm{cos}\phi {\displaystyle \frac{\mathrm{\Gamma }}{R}}\mathrm{sin}\phi +{\displaystyle \frac{\mathrm{\Gamma }}{2}}t_y{\displaystyle \frac{dt_x}{ds}}{\displaystyle \frac{\mathrm{\Gamma }}{2}}t_x{\displaystyle \frac{dt_y}{ds}}\right]𝑑s.`$ (17)
Notice that Eq. (17) can be applied only in the case of highly extended DNA. In the following two sections, we will study the mechanical property of single dsDNA molecules based on the model energy Eq. (16) and Eq. (17). The theoretical results will be compared with experimental observations and discussed.
## III Extensibility and entropic elasticity of DNA
In this section we investigate the elastic responses of single DNA molecules under the actions of external forces based on the model introduced in Sec. II. There is no external torque acted, thus $`\mathrm{\Gamma }=0`$ in Eq. (16). The particular form of the energy function Eq. (16) of the present model makes it convenient for us to study its statistical property by path integral method. In Appendix A a detailed description on the application of path integral method to polymer physics is given . Our calculations in this and the next sections are based on this method.
For a polymer whose energy is expressed in the form of Eq. (16) with $`\mathrm{\Gamma }=0`$, according to the technique outlined in Appendix A (see Eqs. (A1) and (A7)), the Green equation governing the evolution of the $`\mathrm{`}\mathrm{`}`$wave function” $`\mathrm{\Psi }(𝐭,\phi ;s)`$ of the system is obtained to be of the following form:
$$\frac{\mathrm{\Psi }(𝐭,\phi ;s)}{s}=\left[\frac{^2}{4\mathrm{}_p^{}𝐭^2}+\frac{^2}{4\mathrm{}_p\phi ^2}+\frac{f\mathrm{cos}\phi }{k_BT}𝐭𝐳_0\frac{\rho (\phi )}{k_BT}\frac{\mathrm{}_p}{R^2}\mathrm{sin}^4\phi \right]\mathrm{\Psi }(𝐭,\phi ;s),$$
(18)
where $`\mathrm{}_p^{}=\kappa ^{}/k_BT`$ and $`\mathrm{}=\kappa /k_BT`$. The spectrum of the above Green equation is discrete and, for a long dsDNA molecule according to Eqs. (A11) and (A15), its average extension can be obtained either by differentiation of the ground-state eigenvalue, $`g_0`$, of Eq. (18) with respect to $`f`$:
$$Z=_0^L𝐭𝐳_0\mathrm{cos}\phi 𝑑s=Lk_BT\frac{g_0}{f},$$
(19)
or by a direct integration with the normalized ground-state eigenfunction, $`\mathrm{\Phi }_0(𝐭,\phi )`$, of Eq. (18):
$$Z=L|\mathrm{\Phi }_0|^2𝐭𝐳_0\mathrm{cos}\phi d𝐭d\phi .$$
(20)
Both $`g_0`$ an $`\mathrm{\Phi }_0(𝐭,\phi )`$ can be obtained numerically through standard diagonalization methods and identical results are obtained by Eqs. (19) and (20). Here we just briefly outline the main procedures in converting Eq. (18) into the form of a matrix.
Firstly, for our convenience we perform the following transformation:
$$\phi =\stackrel{~}{\phi }\frac{\pi }{2},$$
(21)
hence the new argument $`\stackrel{~}{\phi }`$ can change in the range from $`0`$ to $`\pi `$. Then, we choose the combination of $`Y_{lm}(𝐭)`$ and $`f_n(\stackrel{~}{\phi })`$ as the base functions of the Green equation Eq. (18):
$$\mathrm{\Psi }(𝐭,\stackrel{~}{\phi };s)=\underset{lmn}{}C_{lmn}(s)Y_{lm}(𝐭)f_n(\stackrel{~}{\phi }).$$
(22)
In the above expression, $`Y_{lm}(𝐭)=Y_{lm}(\theta ,\varphi )`$ ($`l=0,1,2,\mathrm{};m=0,\pm 1,\mathrm{},\pm l)`$ are the spherical harmonics , where $`\theta `$ and $`\varphi `$ are the two directional angles of $`𝐭`$, i.e., $`𝐭=(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )`$; and
$$f_n(\stackrel{~}{\phi })=\sqrt{\frac{2}{a}}\mathrm{sin}(\frac{n\pi }{a}\stackrel{~}{\phi })(n=1,2,\mathrm{})$$
(23)
are the eigenfunctions of one-dimensional infinitely deep square potential well of width $`a`$. In the actual calculations, the right boundary of the square well is chosen to be slightly less than $`\pi `$ to avoid flush-off of computer memory caused by the divergence of the base-stacking potential Eq. (8) at $`\stackrel{~}{\phi }=\pi `$. We set $`a=0.95\pi `$ in this paper. However, we have checked that the results are almost identical for other values of $`a`$, provided that $`a165^{}`$.
With the wave function $`\mathrm{\Psi }`$ being expanded using the above mentioned base functions, the operator acting on $`\mathrm{\Psi }`$ (i.e., the expression in the square brackets of Eq. (18)) can also be written into matrix form under these base functions. This matrix, whose elements are listed in Appendix B, is then diagolized numerically to obtain its ground-state eigenvalue and eigenfunction. To simplify the calculation, we further notice that in the present case of Eq. (18), the ground-state is independent to $`\varphi `$, i.e., $`m`$ can be set to $`m=0`$ in Eq. (22).
The resulting force vs extension relation obtained from Eq. (19) or Eq. (20) is shown in Fig. 2 in the whole relevant force range and compared with the experimental observation of Cluzel et al. . The theoretical curve in this figure is obtained with just one adjustable parameter (see caption of Fig. 2); the agreement with experiment is excellent. Figure 2 demonstrates that the highly extensibility of DNA molecule under large external forces can be quantitatively explained by the present model.
To further understand the force-induced extensibility of DNA, in Fig. 3 the folding angle distribution of dsDNA molecule is shown, with the external force kept at different values. Here, according to Eq. (A15) the folding angle distribution $`P(\phi )`$ is calculated by the following formula:
$$P(\phi )=|\mathrm{\Phi }_0(𝐭,\phi )|^2𝑑𝐭.$$
(24)
Figures 2 and 3, taking together, demonstrate that the elastic behaviors of dsDNA molecule are radically different under the condition of low and large applied forces. In the following, we will discuss them separately.
The low-force region When external force is low ($`10`$ pN), the folding angle is distributed narrowly around the angle of $`\phi +57^{}`$, and there is no probability for the folding angle to take on values less than $`0^{}`$ (Fig. 3), indicating that DNA chain is completely in the right-handed B-form configuration with small axial fluctuations. This should be attributed to the strong base-stacking intensity, as pointed out in Sec. II B. Consequently, the elasticity of DNA is solely caused by thermal fluctuations in the axial tangent $`𝐭`$ (Fig. 2), and DNA molecule can be regarded as an inextensible chain. This is the physical reason why, in this force region, the elastic behavior of DNA can be well described by the wormlike chain model . Indeed, as shown in Fig. 4, at forces $`10`$ pN, the wormlike chain model and the present model give identical results. Thus, we can conclude with confidence that, when external fields are not strong, the wormlike chain model is a good approximation of the present model to describe the elastic property of dsDNA molecules; and the bending persistence length of the molecule is $`2\mathrm{}_p^{}\mathrm{cos}\phi `$, as indicated by Eq. (3) and Eq. (4).
The large-force region With the continuous increase of external pulling forces, the axial fluctuations becomes more and more significant. For example, at forces $`50`$ pN, although the folding angle distribution is still peaked at $`\phi 57^{}`$, there is also considerable probability for the folding angle to be distributed in the region $`\phi 0^{}`$ (Fig. 4). Therefore, at this force region, DNA polymer can no longer be regarded as inextensible. At $`f65`$ pN, another peak in the folding angle distribution begins to emerge at $`\phi 0^{}`$, marking the onset of cooperative transition from B-form DNA to overstretched S-form DNA . This is closely related to the short-ranged nature of the base-stacking interactions (see Sec. II B). At even higher forces ($`f80`$ pN) This threshold $`f_t`$ of over-stretch force is also consistent with a plain evaluation from base-stacking potential of $`ϵf_tr_0`$, i.e. $`f_t90`$ pN., the DNA molecule becomes completely into the overstretched form with its folding angle peaked at $`\phi =0^{}`$.
The force-induced axial fluctuations in DNA double-helix can be biologically significant. For example, it has been demonstrated that axial fluctuations in dsDNA enhance considerably the polymerization of RecA proteins along DNA chain . An quantitative study on the coupling between RecA polymerization and DNA axial fluctuation is anticipated to be helpful.
It seems that in the experiments the transition to S-DNA occurs even more cooperatively and abruptly than predicted by the present theory (see Fig. 2). This may be related to the existence of single-stranded breaks (nicks) in the dsDNA molecules used in the experiments. Nicks in DNA backbones can lead to strand-separation or relative sliding of backbones , and they can make the transition process more cooperative. However, the comprehensive agreement achieved in Figs. 2 and 4 indicates that such effects are only of limited significance. The elasticity of DNA is mainly determined by the competition between folding angle fluctuation and tangential fluctuation, which are governed, respectively, by the base-stacking interactions ($`ϵ`$) and the axial bending rigidity ($`\kappa ^{}`$) in Eq. (16).
## IV Elastic Property of supercoiled DNA
In the preceding section, we have discussed the elastic response of long DNA chains under the action of external forces. In the present section, we turn to study the elasticity of supercoiled DNA double-helix. For this purpose, in the experimental setup, all the possible nicks in the DNA nucleotide strands are ligated , and a torque as well as an external pulling force is acting on one terminal of the DNA double-helix, which typically untwists or overtwists the original B-form double-helix to some extent and makes its total linking number (refer to Sec. II C for the definition of the linking number) less or greater than the equilibrium value. We say that such DNA molecules with deficit (excess) linking number are negatively (positively) supercoiled, and define the degree of supercoiling as
$$\sigma =\frac{LkLk_0}{Lk_0},$$
(25)
where $`Lk_0`$ represents the linking number of a relaxed DNA of the same contour length. In living organisms, DNA molecules are often negatively supercoiled, with a linking number deficit of about $`\sigma =0.06`$. Thus, a detailed investigation on the mechanical property of supercoiled DNA molecules is not only of academic interest but can also help us to understand the possible biological advantages of negative supercoiling.
### A Relationship between extension and supercoiling degree
We focus on the property of highly extended DNA molecules whose tangent vectors fluctuate only slightly around the force direction $`𝐳_0`$. According to what we have mentioned in Sec. II C, in this case the approximate energy expression Eq. (17) can be used. The external stretching force is restricted to be greater than $`0.3`$ pN to make sure that the end-to-end distance of DNA chain approaches its contour length (see Fig. 5). Based on Eqs. (17) and (A7), the Green equation for highly stretched and supercoiled dsDNA is then obtained to be
$`{\displaystyle \frac{\mathrm{\Psi }(𝐭,\phi ;s)}{s}}`$ $`=`$ $`[{\displaystyle \frac{^2}{4\mathrm{}_p^{}𝐭^2}}+{\displaystyle \frac{^2}{4\mathrm{}_p\phi ^2}}+{\displaystyle \frac{f\mathrm{cos}\phi }{k_BT}}𝐭𝐳_0{\displaystyle \frac{\rho (\phi )}{k_BT}}{\displaystyle \frac{\mathrm{}_p}{R^2}}\mathrm{sin}^4\phi `$ (26)
$`+{\displaystyle \frac{\mathrm{\Gamma }}{Rk_BT}}\mathrm{sin}\phi {\displaystyle \frac{\mathrm{\Gamma }}{4k_BT\mathrm{}_p^{}}}{\displaystyle \frac{}{\varphi }}+{\displaystyle \frac{\mathrm{\Gamma }^2}{16\mathrm{}_p^{}(k_BT)^2}}\mathrm{sin}^2\theta ]\mathrm{\Psi }(𝐭,\phi ;s),`$ (27)
where $`(\theta ,\varphi )`$ are the two directional angles of $`𝐭`$ as mentioned before in Sec. III. Similar to what we have done in Sec. III, we can now express the above Green equation in matrix form using the combinations of spherical harmonics $`Y_{lm}(\theta ,\varphi )`$ and $`f_n(\phi )`$ as the base functions. The ground-state eigenvalue and eigenfunction of Eq. (27) can then be obtained numerically for given applied force and torque and the average extension be calculated through Eq. (19) or through the following formula:
$$Z=L\chi _0(𝐭,\phi )𝐭𝐳_0\mathrm{cos}\phi \mathrm{\Phi }_0(𝐭,\phi )𝑑𝐭𝑑\phi ,$$
(28)
where $`\chi _0(𝐭,\phi )`$ is the ground-state left-eigenfunction of Eq. (27). <sup>§</sup><sup>§</sup>§As remarked in Appendix A, because the operator in the square brackets of Eq. (27) acting on $`\mathrm{\Psi }(𝐭,\phi ;s)`$ is not Hermitian, the resulting matrix form of the operator may not be diagonalized by unitary matrices. Consequently, in general $`\chi _0(𝐭,\phi )\mathrm{\Phi }_0^{}(𝐭,\phi )`$. The writhing number Eq. (13) is calculated according to Eq. (A18) to be
$$Wr=L\frac{\mathrm{\Gamma }}{16\pi \mathrm{}_p^{}k_BT}\chi _0(𝐭,\phi )\mathrm{sin}^2\theta \mathrm{\Phi }_0(𝐭,\phi )𝑑𝐭𝑑\phi ,$$
(29)
and average linking number is then calculated to be
$$Lk=Tw+Wr=\frac{L}{2\pi R}\chi _0\mathrm{sin}\phi \mathrm{\Phi }_0d𝐭d\phi +\frac{L\mathrm{\Gamma }}{16\pi \mathrm{}_p^{}k_BT}\chi _0\mathrm{sin}^2\theta \mathrm{\Phi }_0d𝐭d\phi .$$
(30)
Thus, after we have obtained the ground-state eigenvalue as well as its left- and right-eigenfunction numerically, we can calculate numerically all the quantities of our interest, for example, the average extension, the average supercoiling degree, the folding angle distribution (see also Appendix A). The relation between extension and supercoiling degree can also be obtained by fixing the external force and changing the value of the applied torque. To calculate the ground-state eigenvalue and eigenfunctions of an asymmetric matrix turns out to be complicated and time-consuming. Fortunately, as we have calculated in Appendix B, each eigenfunction of Eq. (27) shares the same quantity $`m`$; the matrix for $`m=0`$ is still Hermitian and can be diagonalized by unitary matrices. The ground-state eigenvalue for $`m=0`$ is lower in several order than those for $`m0`$ in the whole relevant region of external torque $`\mathrm{\Gamma }`$ from $`5.0k_BT`$ to $`5.0k_BT`$. Thus, actually we only need to consider the case of $`m=0`$ and in this case we still have $`\chi _0(𝐭,\phi )=\mathrm{\Phi }_0^{}(𝐭,\phi )`$. The whole procedure we have performed in Sec. III can safely be repeated in this section, and the relationships between force and extension and between torque and linking number can be consequently calculated.
To make the calculation further easier and also to make sure that the above-mentioned calculation is indeed correct, we introduce here an approximate method which reduces the computational complexity considerably. It turns out that the calculated results using this method are in considerable agreement with the above-mentioned precise method. In the experiment of Ref. , the applied external forces change in the region of $`0.3`$ pN to $`10`$ pN. In this region, as demonstrated in Figs. 3 and 4, both the tangential ($`𝐭`$) and the folding angle ($`\phi `$) fluctuations of dsDNA molecules are small. Taking into account this fact, then in Eq. (17), the energy term $`\kappa ^{}(d𝐭/ds)^2`$ can be approximately calculated to be $`\kappa ^{}((dt_x/ds)^2+(dt_y/ds)^2)`$, and $`f𝐭𝐳_0\mathrm{cos}\phi f\mathrm{cos}\phi f\mathrm{cos}\phi (t_x^2+t_y^2)/2`$. Thus, Eq. (17) is decomposed into two $`\mathrm{`}\mathrm{`}`$independent” parts. The first part is only related to $`\phi `$. At each value of $`f`$ and $`\mathrm{\Gamma }`$, we can calculate the average quantities $`\mathrm{cos}\phi `$ and $`\mathrm{sin}\phi `$ based on this energy using the method of path integral. The second part is quadratic in $`t_x`$ and $`t_y`$, therefore the average values of $`𝐭𝐳_0`$ and $`Wr`$ (Eq. (13)) can be obtained analytically. Using this decomposition and preaveraging technique, the average extension and average supercoiling degree can both be calculated at each value of external force and torque, and the relation between extension and linking number at fixed forces can be then obtained.
The theoretical relationship between extension and supercoiling degree is shown in Fig. 5 and compared with the experiment of Strick et al. . In obtaining these curves, the values of the parameters are the same as those used in Fig. 2 and no adjustment has been done to fit the experimental data. We find that in the case of negatively supercoiled DNA, the theoretical and experimental results are in quantitative agreement, indicating that the present model is capable of explaining the elasticity of negatively supercoiled DNA; in the case of positively supercoiled DNA, the agreement between theory and experiment is not so good, especially when the external force is relatively large. In our present work, we have not considered the possible deformations of the nucleotide basepairs. While this assumption might be reasonable in the negatively supercoiled case, it may fail for positively supercoiled DNA chain, especially at large stretching forces. The work done by Allemand et al. suggested that positive supercoiled and highly extended DNA molecule can take on Pauling-like configurations with exposed bases. To better understand the elastic property of positively supercoiled DNA, it is certainly necessary for us to take into account the deformations of basepairs.
For negatively supercoiled DNA molecule, both theory and experiment reveal the following elastic aspects: (a) When external force is small, DNA molecule can shake off its torsional stress by writhing its central axis, which can lead to an increase in the negative writhing number and hence restore the local folding manner of DNA strands to that of B-form DNA: (b) However, writhing of the central axis causes shortening of DNA end-to-end extension, which becomes more and more unfavorable as the external force is increased. Therefore, at large forces, the torsional stress caused by negative torque (supercoiling degree) begins to unwind the B-form double-helix and triggers the transition of DNA internal structure, where a continuously increasing portion of DNA takes on some certain new configuration as supercoiling increases, while its total extension keeps almost invariant. Our Monte Carlo simulations have also confirmed the above insight .
What is the new configuration? According to the opinion of Ref. , such new configuration corresponds to denatured DNA segments, i.e., negative torque leads to breakage of hydrogen bonds between the complementary DNA bases and consequently to strand-separation. They have also done an elegant experiment in which short single-stranded homologous DNA segments are inserted into the experimental buffer . They found that, in confirmative with their insight, these homologous DNA probes indeed bind onto negatively supercoiled dsDNA molecules. Recently, L$`\stackrel{´}{\mathrm{e}}`$ger et al. have also done experiment on single dsDNA molecules. To explain qualitatively their experimental result, they found that left-handed Z-form DNA should be considered as a possible configuration for negatively supercoiled dsDNA chain, while the molecules need not be denatured. As seen in Fig. 5, although the present model has not taken into account the possibility of strand-separation, it can quantitatively explain the behavior of negatively supercoiled DNA. Therefore, it may be helpful for us to investigate the possibility of formation of left-handed configurations based on our present model. This effort is done in the next subsection, where the energetics of such configurations will also be discussed.
### B Possible left-handed DNA configurations
We have mentioned in Sec. II A that left-handed configurations correspond to $`\phi <0`$ in our present model (see also Fig. 1). Based on the present model, then information about the new configuration mentioned in Sec. IV A can be revealed by the folding angle distributions $`P(\phi )`$, as have been discussed in Sec. III. In the present case, $`P(\phi )`$ is calculated as follows:
$$P(\phi )=\chi _0(𝐭,\phi )\mathrm{\Phi }_0(𝐭,\phi )𝑑𝐭,$$
(31)
where, as mentioned in Sec. IV A, $`\mathrm{\Phi }_0(𝐭,\phi )`$ and $`\chi _0(𝐭,\phi )`$ are, respectively, the ground-state right- and left-eigenfunction of Eq. (27); and actually, $`\chi _0(𝐭,\phi )=\mathrm{\Phi }_0^{}(𝐭,\phi )`$.
The calculated folding angle distribution is shown in Fig. 6, which is radically different with that of torsionally relaxed dsDNA molecules shown in Fig. 3. It has the following aspects: When the torsional stress is small (with the supercoiling degree $`|\sigma |<0.025`$), the distribution has only one narrow and steep peak at $`\phi +57.0^{}`$, indicating that DNA is completely in B-form. With the increase of torsional stress, however, another peak appears at $`\phi 48.6^{}`$ and the total probability for the folding angle to be negatively-valued increases gradually with supercoiling. Since negative folding angles correspond to left-handed configurations, the present model suggests that, with the increasing of supercoiling, left-handed DNA conformation is nucleated and it then elongates along the DNA chain as B-DNA disappears gradually. The whole chain becomes completely left-handed at $`\sigma 1.85`$.
It is worth to be noticed that, (a) as the supercoiling degree changes, the positions of the two peaks of the folding angle distribution remain almost fixed and, (b) between these two peaks, there exists an extended region of folding angle from $`0`$ to $`\pi /6`$ which always has only extremely small probability of occurrence. Thus, a negatively supercoiled DNA can have two possible stable configurations, a right-handed B-form and a left-handed configuration with an average folding angle $`48.6^{}`$. A transition between these two structures for a DNA segment will generally lead to an abrupt and finite variation in the folding angle.
To obtain the energetics of such transitions, we have calculated how the sum of the base-stacking energy and torsional energy, $`(\kappa /R^2)\mathrm{sin}^4\phi +\rho (\phi )(\mathrm{\Gamma }/R)\mathrm{sin}\phi `$, changes with external torque . Figure 7a shows the numerical result, and Fig. 7b demonstrates the relation between supercoiling degree and external torque. (In both figures the external force is fixed at $`1.3`$ pN.) From these figures we can infer that, (a) for negative torque less than the critical value $`\mathrm{\Gamma }_c3.8`$ k<sub>B</sub>T, DNA can only stay in B-form state; (b) near this critical torque, DNA can either be right- or be left-handed and, as negative supercoiling increases (see Fig. 7b) more and more DNA segments will stay in the left-handed form, which is much lower in energy ($`2.0`$ k<sub>B</sub>T per basepair) but stable only when torque reaches $`\mathrm{\Gamma }_c`$; (c) for negative torque greater than $`\mathrm{\Gamma }_c`$ DNA is completely left-handed.
Nevertheless, we should emphasize that the above calculations are all based on our present model which has assumed that nucleotide basepairs do not break. Figure 5 indicates that for negatively supercoiled DNA chains, the extension vs supercoiling degree relation can be quantitatively explained by the present model; and Fig. 6 reveals the reason of the quantitative agreement is that the present model allows the possibility of occurrence of left-handed DNA configurations. At the present time, to say that negatively supercoiled DNA will prefer left-handed configurations rather than denaturation and strand-separation is premature. To clarify this question, the present model should be improved to consider the deformations of DNA basepairs. In the experimental side, it might also be helpful to measure precisely the critical torque at which the elastic behavior of negatively supercoiled DNA changes abruptly and compare the measured results with the value calculated in the present work. (In the earlier experiment of Allemand et al. , the critical torque is estimated to be $`2k_BT`$ by assuming that the torsional rigidity of dsDNA to be $`75`$ nm and that torsional stress builds up linearly along the DNA chain. This value, however, maybe not precise enough since the torsional rigidity of dsDNA molecule is not a precisely determined quantity and the values giving by different groups are scattered widely.Another possibility may be that we have over-estimated the value of the critical torque. It maybe possible that the transition from right- to left-handed configurations is initiated in the weaker AT-rich regions, whose value of $`ϵ`$ should be less than the average value taken by the present paper.
The structural parameters of the left-handed configuration suggested by Figs. 6 and 7 are listed in table I and compared with those of Z-form DNA . The strong similarity in these parameters suggests that the torque-induced left-handed configurations, if they really exist, belong to Z-form DNA .
## V conclusion
In this article, we have presented an elastic model for double-stranded biopolymers such as DNA molecules. The key progress is that the bending deformations of the backbones of DNA molecules and the base-stacking interactions existing between adjacent DNA basepairs are quantitatively considered in this model, with the introduction of a new structural parameter, the folding angle $`\phi `$. This model has also qualitatively taken into account the effects of the steric effects of DNA basepairs and electrostatic interactions along DNA chain. In calculation technique, the model is investigated using path integral method; and Green equations similar in form to the Schr$`\ddot{\mathrm{o}}`$dinger equation in quantum mechanics are derived and their ground-state eigenvalues and eigenfunctions obtained by precise numerical calculations. The force-extension relationship in torsionally relaxed and the extension-linking number relationship in torsionally constrained DNA chains are studied and compared with experimental results. This work demonstrated that DNA molecule’s entropic elasticity and highly extensibility, as well as the elastic property of negatively supercoiled DNA can all be quantitatively explained by the present theory. The comprehensive agreement between theory and experiments indicated that the short-ranged base-stacking interactions are very important in determining the elastic response of double-stranded DNA molecules. The present work showed that highly extended and negatively supercoiled DNA molecules can be left-handed, probably in the Z-form configuration. A possible way to check the validity of this opinion is to measure the critical external torque at which the transition between B-form DNA and the new configuration takes place.
The present work regarded DNA basepairs as rigid objects and did not consider their possible deformations and the possibility of strand-separation. The comprehensive agreement between theory and experiments indicates that this approach is well justified in many cases. However, as already mentioned in Sec. II A, under some extreme conditions, this assumption may not be appropriate. For example, recent experimental work of Allemand et al. showed that positively supercoiled DNA under high applied force can take on Pauling-like configurations with exposed bases. In this case the basepairing of DNA is severely distorted, and because the present model has not taken into account the possible deformations of the basepairs, the theoretical results on positively supercoiled DNA molecules were not in quantitative agreement with experiment (see Fig. 5). Furthermore, although the present work showed that left-handed DNA configurations can be stabilized by negative torques, much theoretical work is still needed to calculate the denaturation free energy and be compared with the free energy of left-handed DNA configurations.
## VI Acknowledgement
It gives us great pleasure to acknowledge the helps and valuable suggestions of our colleagues, especially Lou Ji-Zhong, Liu Lian-Shou, Yan Jie, Liu Quan-Hui, Zhou Xin, Zhou Jian-Jun, and Zhang Yong. The numerical calculations are performed partly at ITP-Net and partly at the State Key Lab. of Scientific and Engineering Computing.
## A Path integral method in polymer physics
In this appendix we review some basic ideas on the application of path integral method to the study of polymeric systems . Consider a polymeric string, and suppose its total $`\mathrm{`}\mathrm{`}`$arclength” is $`L`$, and along each arclength point $`s`$ one can define a n-dimensional $`\mathrm{`}\mathrm{`}`$vector” $`𝐫(s)`$ to describe the polymer’s local state at this point.For example, in the case of a flexible Gaussian chain, $`𝐫`$ is a three-dimensional position vector; in the case of a semiflexible chain such as the wormlike chain , $`𝐫`$ is the unit tangent vector of the polymer and is two-dimensional. We further assume that the energy density (per unit arclength) of the polymer can be written as the following general form:
$$\rho _e(𝐫,s)=\frac{m}{2}(\frac{d𝐫}{ds})^2+𝐀(𝐫)\frac{d𝐫}{ds}+V(𝐫),$$
(A1)
where $`V(𝐫)`$ is a scalar field and $`𝐀(𝐫)`$ is a vectorial field. The total partition function of the system is expressed by the following integration:
$$\mathrm{\Xi }(L)=𝑑𝐫_f\varphi _f(𝐫_f)G(𝐫_f,L;𝐫_i,0)\varphi _i(𝐫_i)𝑑𝐫_i,$$
(A2)
where $`\varphi _i(𝐫)`$ and $`\varphi _f(𝐫)`$ are, respectively, the probability distributions of the vector $`𝐫`$ at the initial ($`s=0`$) and final ($`s=L`$) arclength point; $`G(𝐫,s;𝐫^{},s^{})`$ is called the Green function, it is defined in the following way:
$$G(𝐫,s;𝐫^{},s^{})=_𝐫^{}^𝐫𝒟[𝐫^{\prime \prime }(s)]\mathrm{exp}[\beta _s^{}^s𝑑s^{\prime \prime }\rho _e(𝐫^{\prime \prime },s^{\prime \prime })],$$
(A3)
where integration is carried over all possible configurations of $`𝐫^{\prime \prime }`$, and $`\beta =1/k_BT`$ is the Boltzmann coefficient. It can be verified that the Green function defined above satisfies the following relation :
$$G(𝐫,s;𝐫^{},s^{})=𝑑𝐫^{\prime \prime }G(𝐫,s;𝐫^{\prime \prime },s^{\prime \prime })G(𝐫^{\prime \prime },s^{\prime \prime };𝐫^{},s^{}),(s^{}<s^{\prime \prime }<s).$$
(A4)
The total free energy of the system is then expressed as
$$=k_BT\mathrm{ln}\mathrm{\Xi }.$$
(A5)
To calculate the total partition function $`\mathrm{\Xi }`$, we define an auxiliary function $`\mathrm{\Psi }(𝐫,s)`$ and call it the wave function because of its similarity with the true wave function of quantum systems. Suppose the value of $`\mathrm{\Psi }`$ at arclength point $`s`$ is related to its value at $`s^{}`$ through the following formula such that<sup>\**</sup><sup>\**</sup>\**In fact, the choice of the wave function $`\mathrm{\Psi }(𝐫,s)`$ is not limited. Any function determined by an integration of the form of Eq. (A6) can be viewed as a wave function.
$$\mathrm{\Psi }(𝐫,s)=𝑑𝐫^{}G(𝐫,s;𝐫^{},s^{})\mathrm{\Psi }(𝐫^{},s^{}),(s>s^{})$$
(A6)
then, we can derive from Eqs. (A1), (A3), and (A6) that :
$$\frac{\mathrm{\Psi }(𝐫,s)}{s}=\left[\frac{_𝐫^2}{2m\beta }\beta V(𝐫)+\frac{𝐀(𝐫)_𝐫}{m}+\frac{_𝐫𝐀(𝐫)}{2m}+\frac{\beta 𝐀^2(𝐫)}{2m}\right]\mathrm{\Psi }(𝐫,s)=\widehat{H}\mathrm{\Psi }(𝐫,s).$$
(A7)
Equation (A7) is called the Green equation, it is very similar with the Schr$`\ddot{\mathrm{o}}`$dinger equation of quantum mechanics . However, there is an important difference. In the case of $`𝐀(𝐫)0`$, the operator $`\widehat{H}`$ in Eq. (A7) is not Hermitian. Therefore, in this case the matrix form of the operator $`\widehat{H}`$ may not be diagonalized by unitary matrix.
Denote the eigenvalues and the right-eigenfunctions of Eq. (A7) as $`g_i`$ and $`|i=\mathrm{\Phi }_i(𝐫)`$ ($`i=0,1,\mathrm{}`$), respectively. Then it is easy to know, from the approach of quantum mechanics, that
$$\mathrm{\Psi }(s)=\underset{i}{}e^{g_i(ss^{})}|i(s)i(s^{})|\mathrm{\Psi }(s^{}),$$
(A8)
where $`i|=\chi _i(𝐫)`$ $`(i=0,1,\mathrm{})`$ denote the left-eigenfunctions of Eq. (A7), which satisfy the following relation:
$$i|i^{}=𝑑𝐫\chi _i(𝐫)\mathrm{\Phi }_i^{}(𝐫)=\delta _i^i^{}.$$
In the case where $`\widehat{H}`$ is Hermitian (i.e., $`𝐀(𝐫)=0`$), then we can conclude that
$$\chi _i(𝐫)=\mathrm{\Phi }_i^{}(𝐫).$$
From Eqs. (A2), (A6), (A7), and (A8) we know that
$`\mathrm{\Xi }(L)`$ $`=`$ $`{\displaystyle 𝑑𝐫𝑑𝐫^{}G(𝐫,L;𝐫^{},s^{})\varphi _f(𝐫)\varphi _i(𝐫^{})}={\displaystyle \underset{i}{}}\varphi _f|ii|\varphi _ie^{g_iL}`$ (A9)
$`=`$ $`e^{g_0L}\varphi _f|00|\varphi _i(\mathrm{for}L1/(g_1g_0)).`$ (A10)
Consequently, for long polymer chains the total free energy density is just expressed as
$$/L=k_BTg_0,$$
(A11)
and any quantity of interest can then be calculated by differentiation of $``$. For example, the average extension of a polymer under external force field $`f`$ can be calculated as $`Z=/f=Lk_BTg_0/f`$.
We continue to discuss another very important quantity, the distribution probability of $`𝐫`$ at arclength $`s`$, $`P(𝐫,s)`$. This probability is calculated from the following expression:
$$P(𝐫,s)=\frac{𝑑𝐫_f𝑑𝐫_i\varphi _f(𝐫_f)G(𝐫_f,L;𝐫,s)G(𝐫,s;𝐫_i,0)\varphi _i(𝐫_i)}{𝑑𝐫_f𝑑𝐫_i\varphi _f(𝐫_f)G(𝐫_f,L;𝐫_i,0)\varphi _i(𝐫_i)}.$$
(A12)
Based on Eqs. (A6) and (A8) we can rewrite Eq. (A12) in the following form:
$`P(𝐫,s)`$ $`=`$ $`{\displaystyle \frac{𝑑𝐫_f𝑑𝐫_i𝑑𝐫^{}\varphi _f(𝐫_f)G(𝐫_f,L;𝐫^{},s)\delta (𝐫^{}𝐫)G(𝐫,s;𝐫_i,0)\varphi _i(𝐫_i)}{𝑑𝐫_f𝑑𝐫_i\varphi _f(𝐫_f)G(𝐫_f,L;𝐫_i,0)\varphi _i(𝐫_i)}}`$ (A13)
$`=`$ $`{\displaystyle \frac{\underset{m}{}\underset{n}{}\varphi _f|mn|\varphi _i\chi _m(𝐫)\mathrm{\Phi }_n(𝐫)\mathrm{exp}[g_m(Ls)g_ns]}{\underset{m}{}\varphi _f|mm|\varphi _i\mathrm{exp}(g_mL)}}`$ (A14)
For the most important case of $`0sL`$, Eq. (A13) then gives that the probability distribution of $`𝐫`$ is independent of arclength $`s`$, i.e.,
$$P(𝐫,s)=\chi _0(𝐫)\mathrm{\Phi }_0(𝐫)(\mathrm{for}\mathrm{\hspace{0.33em}\hspace{0.33em}0}sL).$$
(A15)
With the help of Eq. (A15), the average value of a quantity which is a function of $`𝐫`$ can be obtained. For example,
$$Q(s)=𝑑𝐫Q(𝐫)P(𝐫,s)=𝑑𝐫\chi _0(𝐫)Q(𝐫)\mathrm{\Phi }_0(𝐫)=0|Q|0,$$
(A16)
and
$$_0^LQ(𝐫(s))𝑑s=_0^LQ(s)𝑑s=L0|Q|0(\mathrm{for}L1/(g_1g_0)).$$
(A17)
Finally, we list the formula for calculating $`𝐁(𝐫)d𝐫/ds`$, here $`𝐁(𝐫)`$ is a given vectorial field. The formula reads:
$$𝐁(𝐫)\frac{d𝐫}{ds}=\frac{1}{2}0|\left[𝐫(\widehat{H}𝐁)𝐁(\widehat{H}𝐫)\right]|0(\mathrm{for}L1/(g_1g_0)).$$
(A18)
## B Matrix formalism of the Green Equation Eq. (23)
In Eq. (27) (and also Eq. (18)) the variables $`𝐭`$ and $`\phi `$ are coupled together. Denote the operator in the square brackets of Eq. (27) as $`\widehat{H}`$. We express this operator in matrix form. To this end we choose the base functions of this system to be the combinations of spherical harmonics $`Y_{lm}(\theta ,\varphi )`$ and $`f_n(\stackrel{~}{\phi })`$ (see Eq. (23)).
For $`m=0`$, the base functions are expressed to be
$$|i=|lN_\phi +n=|l;n=Y_{l0}(\theta ,\varphi )f_n(\stackrel{~}{\phi });$$
(B1)
and for $`m=1,2,\mathrm{}`$, the base functions are expressed to be
$$|i=|2(lm)N_\phi +2(n1)+k=|l;n;k=Y_{lm}^{(k)}(\theta ,\varphi )f_n(\stackrel{~}{\phi }).$$
(B2)
In the above two equations, $`n=1,2,\mathrm{},N_\phi `$ with $`N_\phi `$ set to be $`60`$ in our present calculations; $`l=0,1,\mathrm{},N_l1`$ for the case of $`m=0`$ (with $`N_l`$ set to $`30`$), or $`l=m,m+1,\mathrm{},m+N_l1`$ for the case of $`m0`$ (with $`N_l`$ set to $`15`$); $`k=1,2`$; and
$`Y_{l0}(\theta ,\varphi )`$ $`=`$ $`\sqrt{{\displaystyle \frac{2l+1}{4\pi }}}P_l(\mathrm{cos}\theta ),`$ (B3)
$`Y_{lm}^{(1)}(\theta ,\varphi )`$ $`=`$ $`(1)^m\sqrt{{\displaystyle \frac{2l+1}{2\pi }}{\displaystyle \frac{(lm)!}{(l+m)!}}}P_l^m(\mathrm{cos}\theta )\mathrm{cos}(m\varphi ),`$ (B4)
$`Y_{lm}^{(2)}(\theta ,\varphi )`$ $`=`$ $`(1)^m\sqrt{{\displaystyle \frac{2l+1}{2\pi }}{\displaystyle \frac{(lm)!}{(l+m)!}}}P_l^m(\mathrm{cos}\theta )\mathrm{sin}(m\varphi ).`$ (B5)
With the above definitions, we then obtain that, for $`m=0`$:
$`i_p|(\widehat{H})|i`$ $`=`$ $`l_p;n_p|(\widehat{H})|l;n`$ (B6)
$`=`$ $`\delta _l^{l_p}\delta _n^{n_p}\left[{\displaystyle \frac{l(l+1)}{4\mathrm{}_p^{}}}+{\displaystyle \frac{n^2\pi ^2}{4\mathrm{}_pa^2}}\right]`$ (B7)
$``$ $`{\displaystyle \frac{f}{k_BT}}\left[\delta _{l_p}^{l+1}a_{l,0}+\delta _{l_p}^{l1}a_{l1,0}\right]n_p|\mathrm{sin}\stackrel{~}{\phi }|n`$ (B8)
$`+`$ $`\delta _{l_p}^ln_p|\left[{\displaystyle \frac{\rho (\stackrel{~}{\phi })}{k_BT}}+{\displaystyle \frac{\mathrm{}_p}{R^2}}\mathrm{cos}^4\stackrel{~}{\phi }\right]|n`$ (B9)
$`+`$ $`\delta _{l_p}^l{\displaystyle \frac{\mathrm{\Gamma }}{Rk_BT}}n_p|\mathrm{cos}\stackrel{~}{\phi }|n`$ (B10)
$``$ $`\delta _{n_p}^n{\displaystyle \frac{\mathrm{\Gamma }^2}{16\mathrm{}_p^{}(k_BT)^2}}\left[\delta _{l_p}^l(1a_{l,0}^2a_{l1,0}^2)\delta _{l_p}^{l+2}a_{l,0}a_{l+1,0}\delta _{l_p}^{l2}a_{l1,0}a_{l2,0}\right];`$ (B11)
and for $`m0`$:
$`i_p|(\widehat{H})|i`$ $`=`$ $`l_p;n_p;k_p|(\widehat{H})|l;n;k`$ (B12)
$`=`$ $`\delta _{l_p}^l\delta _{n_p}^n\delta _{k_p}^k\left[{\displaystyle \frac{l(l+1)}{4\mathrm{}_p^{}}}+{\displaystyle \frac{n^2\pi ^2}{4\mathrm{}_pa^2}}\right]`$ (B13)
$``$ $`{\displaystyle \frac{f}{k_BT}}\delta _{k_p}^k\left[\delta _{l_p}^{l+1}a_{l,m}+\delta _{l_p}^{l1}a_{l1,m}\right]n_p|\mathrm{sin}\stackrel{~}{\phi }|n`$ (B14)
$`+`$ $`\delta _{l_p}^l\delta _{k_p}^kn_p|\left[{\displaystyle \frac{\rho (\stackrel{~}{\phi })}{k_BT}}+{\displaystyle \frac{\mathrm{}_p}{R^2}}\mathrm{cos}^4\stackrel{~}{\phi }\right]|n`$ (B15)
$`+`$ $`\delta _{l_p}^l\delta _{k_p}^k{\displaystyle \frac{\mathrm{\Gamma }}{Rk_BT}}n_p|\mathrm{cos}\stackrel{~}{\phi }|n+\delta _{l_p}^l\delta _{n_p}^n{\displaystyle \frac{\mathrm{\Gamma }}{4\mathrm{}_p^{}k_BT}}(kk_p)m`$ (B16)
$``$ $`\delta _{n_p}^n\delta _{k_p}^k{\displaystyle \frac{\mathrm{\Gamma }^2}{16\mathrm{}_p^{}(k_BT)^2}}\left[\delta _{l_p}^l(1a_{l,m}^2a_{l1,m}^2)\delta _{l_p}^{l+2}a_{l,m}a_{l+1,m}\delta _{l_p}^{l2}a_{l1,m}a_{l2,m}\right].`$ (B17)
In Eqs. (B11) and (B17), the following notations are used:
$`a_{l,m}=\sqrt{{\displaystyle \frac{(l+1)^2m^2}{(2l+1)(2l+3)}}},`$ (B18)
$`n_p|y(\stackrel{~}{\phi })|n={\displaystyle \frac{2}{a}}{\displaystyle _0^a}\mathrm{sin}({\displaystyle \frac{n_p\pi \stackrel{~}{\phi }}{a}})y(\stackrel{~}{\phi })\mathrm{sin}({\displaystyle \frac{n\pi \stackrel{~}{\phi }}{a}}),`$ (B19)
where $`y(\stackrel{~}{\phi })`$ is any function of $`\stackrel{~}{\phi }`$.
The ground-state eigenvalues of the matrices Eq. (B11) and Eq. (B17) have been calculated at force $`1.3`$ pN in the whole relevant region of $`\mathrm{\Gamma }`$ from $`5.0k_BT`$ to $`+5.0k_BT`$. We have found that the ground-state eigenvalues for the case of $`m=0`$ are of the order of $`10`$ nm<sup>-1</sup>; while those for the case of $`m0`$ are of the order of $`10^6`$ nm<sup>-1</sup> (data not shown). Thus, we can conclude with confidence that the $`\mathrm{`}\mathrm{`}`$low energy” eigenstates of the system all have the same $`m=0`$. This can greatly reduce the calculation tasks.
As a final point, here we demonstrate how to calculate the average values of the quantities of our interest. Suppose $`y(𝐭,\stackrel{~}{\phi })`$ is a quantity whose average value we are interested in. Denote $`𝒰`$ as the unitary matrix which can diagonalize the matrix Eq. (B11). (Because the matrix Eq. (B11) is real symmetric, $`𝒰`$ is actually a real orthogonal matrix.) Then its column vector $`𝐮(i)=_j𝒰(j,i)|j`$ $`(i=1,2,\mathrm{})`$ corresponds to the $`i`$-th eigenvector of Eq. (B11). Consequently, we can calculate based on Eq. (A16) that
$$y(𝐭,\stackrel{~}{\phi })=\underset{i,i_p}{}𝒰(i_p,1)𝒰(i,1)i_p|y(𝐭,\stackrel{~}{\phi })|i.$$
(B20) |
warning/0003/gr-qc0003081.html | ar5iv | text | # 1 The violation of the principle of equivalence
## 1 The violation of the principle of equivalence
The principle of equivalence which is the base of the general theory of relativity reflects the close connection between inertial coordinate system $`K`$ with uniform gravitational field and noninertial coordinate system $`K^{}`$ moving with uniform acceleration in empty space. According to this principle any physical process proceeds absolutely identically in the both systems. In other words, if we imagine the system $`K`$ as a closed laboratory, which is at rest on the Earth, and the system $`K^{}`$ as an identical laboratory, moving with acceleration g in a distance of gravitating masses, and the sizes of the laboratories are such chosen, that it will be possible to neglect the nonuniformity of the gravitational field in system $`K,`$ then the observer, who is in one of these laboratories and who does not have any connection with the external world, i. e. without having any possibility to look out of its limits, could not make any experiment inside the laboratory in order to find out in which of them he is, i. e. he could not define the character of his motion. That can be explained by the fact, that in system $`K^{}`$ during the uniform accelerating the field of inertial forces appears. The action of this field does not differ from the action of the gravitational field.
It is easy to make sure, that the mechanical processes will proceed identically in both laboratories. In fact, free bodies in every laboratory will move with acceleration $`g`$, identical pendulums will oscillate with equal periods, etc. Thus, systems $`K`$ and $`K^{}`$ are equivalent in respect of mechanical processes.
Einstein spread the equivalence of systems $`K`$ and $`K^{}`$ on all physical processes without exceptions, having formulated the principle of equivalence, according to which, not only mechanical, but also any physical processes have to proceed identically in systems $`K`$ and $`K^{}`$.
But we can point to the process which violates the principle of equivalence. It is the radiation of charges. It is easy to make sure, that with the help of charges it is possible to differ system $`K`$ from system $`K^{}`$. In fact, let us place the charges in both laboratories. The charge in the laboratory $`K^{}`$ has to radiate, so it is moving with acceleration. The observer who is in this laboratory can register this radiation having placed, for example, a charge into the water. Then a part of the radiating energy will be absorbed by the water and will warm it up. On measuring the temperature of the water the observer will be to register the radiation.
We won’t discuss the technical details how to carry out such an experiment. It is enough that it is a principal possibility to find out the radiation in immediate proximity to the charge.
In laboratory $`K`$ a motionless charge will not radiate. Thus, the observer in every laboratory very easily can define the character of his motion.
Therefore the principle of equivalence is violated.
Now we will consider uniformly rotating coordinate systems.
We can imagine one of these systems (we mark it as $`S^{}`$) as a disk, revolving with constant angular velocity on its axis. In accordance with the principle of equivalence the noninertial coordinate system $`S^{}`$, in which a field of centrifugal forces of inertia exists, can be considered as an inertial system $`K`$ with an uniform gravitational field.
Let us consider two charges: one of them is on the disk, i. e. rotates with this disk. The second one is at rest at system $`K`$. In the first case the charge has to radiate and in the second one it does not. It has been already told how to differ a radiating charge from a charge, which does not radiate. And what is more, a braking force of radiation friction $`f`$ must act on uniformly moving round a circle charge, which action one can observe at an as short as we want distance from the charge.
Therefore, in this case the principle of equivalence is violated.
So far as we have mentioned the radiation force of friction, it is necessary to point out another problem, which was first remarked by M. Born. By motion of the charge with uniform acceleration, force $`f`$ turns into zero. That leads to the violation of energy balance. Indeed $`f={\displaystyle \frac{2e^2}{3c^3}}\stackrel{\mathrm{}}{𝑟}`$, where $`\stackrel{\mathrm{}}{𝑟}`$ — the third derivative from a coordinate. If the motion firmly accelerates, $`\ddot{r}`$=const and $`\stackrel{\mathrm{}}{𝑟}`$=0. Therefore, an uniformly accelerating charge radiates without any losses of energy. That contradicts the law of conservation of energy.
We will return to laboratories $`K`$ and $`K^{}`$ again. We will change the character of their motion. Let laboratory $`K^{}`$, which is in empty space far from gravitating masses, move uniformly and straightforward and let $`K`$ freely move in a gravitational field. There will be state of weightlessness in both laboratories. According to the principle of equivalence, in this case systems $`K`$ and $`K^{}`$ are also equivalent. I. e. as in the last instance, the observer, who is in one of the laboratories and does not have any connection with the external world, could not make any experiment inside the laboratory in order to find out the character of his motion.
But we will place a charge in each of these laboratories again. In laboratory $`K`$ charge must radiate, for it moves with acceleration under the action of the gravitational field. The observer in this laboratory will be able to register the radiation. In laboratory $`K^{}`$ the charge will not radiate.
Therefore, in this case the principle of equivalence is not fulfilled.
## 2 The violation of laws of conservation in classical electrodynamics
By means of not complicated calculations it can be shown, that there is a violation of laws of conservation of energy, impulse and angular momentum in classical electrodynamics. And the violation of these laws arises when we consider the charges moving with uniform acceleration or uniformly moving round a circle — i. e. those charges, the electrodynamics of which violates the principle of equivalence. One can make sure in that, having calculated the interaction of such charges.
Let charge $`e`$ move arbitrarily. Its field $`\stackrel{}{E}`$ is determined with the following expression
$$\stackrel{}{E}=\frac{e\left(1\frac{v^2}{c^2}\right)}{\left(R\frac{\stackrel{}{R}\stackrel{}{v}}{c}\right)^3}\left(\stackrel{}{R}\frac{\stackrel{}{v}}{c}R\right)+\frac{e}{c^2\left(R\frac{\stackrel{}{R}\stackrel{}{v}}{c}\right)^3}\left[\stackrel{}{R}\left[\left(\stackrel{}{R}\frac{\stackrel{}{v}}{c}R\right)\stackrel{}{\stackrel{}{v}}\right]\right]$$
$`(1)`$
All the quantities in the right part are taken at the moment $`t^{}`$, which precedes the moment $`t`$ of the observation and which can be obtained from the equation
$$t^{}+\frac{R(t^{})}{c}=t$$
where $`R(t^{})`$ is a distance from the charge to the point of observation and $`\stackrel{}{v}`$ and $`\stackrel{}{\stackrel{}{v}}`$ are velocity and acceleration of the charge. The field consists of two parts. The first term describes Coulomb’s field of the charge, and the second one describes the field of radiation.
We will consider two nonrelative charges $`e`$ and $`e`$, connected with a pivot, which length is $`l`$. These charges are moving with acceleration $`\stackrel{}{\stackrel{}{v}}`$, directed along the pivot (fig. 1). Let us define Coulomb’s field $`\stackrel{}{E}_{k1}`$, which the charge $`e`$ creates in the point $`C`$, where charge $`e`$ is. At the moment $`t^{}`$ the charge $`e`$ is in point $`A`$. Vector $`\stackrel{}{R}(t^{})`$ is equal to the length of segment $`AC`$ by its module. Vector $`{\displaystyle \frac{\stackrel{}{v}}{c}}R`$ is equal to the length of segment $`AB`$ (fig. 1). Point $`B`$ is a position, in which charge $`e`$ could be at the moment $`t`$, as if it were moving from point A uniformly with velocity $`v(t^{})`$ during the period of time $`tt^{}={\displaystyle \frac{R}{c}}`$. But the charge $`e`$ is moving with acceleration $`\stackrel{}{\stackrel{}{v}}`$. That is why at the moment $`t`$ it will be in point $`D`$. And segment $`AD=v\tau +{\displaystyle \frac{\dot{v}\tau ^2}{2}}`$, where $`\tau =tt^{}`$.
Therefore the Coulomb’s field $`\stackrel{}{E}_{k1}`$ is equal to
$$E_{k1}=\frac{e}{\left(l\frac{\dot{v}\tau ^2}{2}\right)^2}.$$
For $`vc`$, then $`R(t^{})l`$, $`\tau {\displaystyle \frac{l}{c}}`$, then $`E_{k1}={\displaystyle \frac{e}{\left(l\frac{\dot{v}l^2}{2c^2}\right)^2}}.`$
Therefore, force $`\stackrel{}{F}_1`$, acting on the charge $`e`$, is equal to
$$F_1=\frac{e^2}{\left(l\frac{\dot{v}l^2}{2c^2}\right)^2}.$$
In the same way we will determine field $`\stackrel{}{E}_{k2}`$, which the charge $`e`$ creates in point $`D`$, where the charge $`e`$ is. At the moment $`t^{}`$ the charge $`e`$ will be in point $`K`$ (fig. 1). The segment $`KL=v\tau `$. The segment $`KC=v\tau +{\displaystyle \frac{\dot{v}\tau ^2}{2}}.`$
So field $`\stackrel{}{E}_{k2}`$ is equal to
$$E_{k2}=\frac{e}{\left(l+\frac{\dot{v}l^2}{2c^2}\right)^2}.$$
The force $`\stackrel{}{F}_2`$, acting on the charge $`e`$ is equal to
$$F_2=\frac{e^2}{\left(l+\frac{\dot{v}l^2}{2c^2}\right)^2}.$$
Thus, forces $`\stackrel{}{F}_1`$ and $`\stackrel{}{F}_2`$ are not equal by the quantily. Their resultant $`\mathrm{\Delta }\stackrel{}{F}=\stackrel{}{F}_1+\stackrel{}{F}_2`$ is
$$\mathrm{\Delta }F=\frac{2e^2\dot{v}}{c^2l}$$
and is directed along the vector $`\stackrel{}{\stackrel{}{v}}`$.
The forces of the interaction of the charges, which are specified by the second item in (1), are equal to zero in this case.
In the case, when vector $`\stackrel{}{\stackrel{}{v}}`$ and the axis of the pivot make angle $`\alpha `$, the forces $`\stackrel{}{F}_1`$ u $`\stackrel{}{F}_2`$ are not equal by the quantity and do not lie on the same straight line (vectors $`\stackrel{}{R}`$ and $`{\displaystyle \frac{\stackrel{}{v}}{c}}R`$ of the charge $`e`$ are shown in fig. 2).
The projection of their resultant $`\mathrm{\Delta }\stackrel{}{F}`$ on axis $`x`$ is equal to
$$\mathrm{\Delta }F_x=\frac{e^2\dot{v}}{c^2l}(2\mathrm{cos}^2\alpha \mathrm{sin}^2\alpha )$$
and on axis $`y`$
$$\mathrm{\Delta }F_y=\frac{3e^2\dot{v}}{2c^2l}\mathrm{sin}2\alpha .$$
The forces of the interaction of the charges, specified by the second item in (1), are not equal to zero in the ammount in this case either. The projection of its resultant $`\mathrm{\Delta }F^{}`$ on axis $`x`$ and axis $`y`$ are equal
$$\mathrm{\Delta }F_x^{}=\frac{2e^2\dot{v}}{c^2l}\mathrm{sin}^2\alpha ;$$
$$\mathrm{\Delta }F_y^{}=\frac{e^2\dot{v}}{c^2l}\mathrm{sin}2\alpha .$$
It is quite evidently, that the presence of the resultant $`\mathrm{\Delta }\stackrel{}{F}`$, depending on the orientation of the pivot and acting on the electrically neutral system, contradicts the laws of conservation of energy and impulse. Within the framework of classical electrodynamics it is impossible to compensate this force.
We will consider the following instance. Let cylinder of radius $`r`$ the forming $`d`$ revolve on its axis. We should choose quantaties $`r`$ and $`d`$, so as 2$`\pi r=0,01`$ $`d`$. At the ends of the forming we place charges $`e`$ and $`e`$ (fig. 3). Let linear velocity of the charges be equal to $`v=0,01c`$. In order to calculate Coulomb’s field $`\stackrel{}{E}_k`$, which the charge $`e`$ creates in point $`A`$, where charge $`e`$ is, it is necessary to find out the preceding position of charge $`e`$ at moment $`t^{}`$ — as in the previous instance. It is evidently, that for this case the period of time $`tt^{}`$ is equal to the period of revolution of the cylinder. So, at the moment $`t^{}`$ the charge $`e`$ will be in point B, where it is also at the moment $`t`$. Vectors $`\stackrel{}{R},{\displaystyle \frac{\stackrel{}{v}}{c}}R`$ and $`\stackrel{}{R}{\displaystyle \frac{\stackrel{}{v}}{c}}R`$ will be directed as it is shown in fig. 3 ($`BP=2\pi r`$). The field $`\stackrel{}{E}_k`$ in point $`A`$ will be equal $`E_k{\displaystyle \frac{e}{d^2}}`$ and will direct along the segment $`PA`$.
That is why Coulomb’s force, acting on the charge $`e`$, has got the component on the tangent to the point $`A`$. The quantity of this component is equal to $`F=0,01`$ F<sub>k</sub> where $`F_k`$ is Coulomb’s force.
The force, acting on the charge $`e`$, specified by the second item in(1), will be directed perpendicularly to the velocity of the charge $`e`$.
In the same way we will calculate the effect of charge $`e`$ on charge $`e`$. There fore on the cylinder under review, which is a closed system, the rotational moment $`M=0,02`$ F$`{}_{k}{}^{}r`$ is acting. That contradicts the principles of conservation of energy and angular momentum.
It is easy to show, that if the velocity of the rotation of this cylinder is diminished twice, Coulomb’s interaction of charges $`e`$ and $`e`$ will lead to the braking of the cylinder.
## 3 Rutherford’s atom
One of the most important problems of classical electrodynamics is the problem of the stability of Rutherford’s atom. As it is known, the planetary model of atom, offered by Rutherford, has a principal shortcoming — it was unstable. According to classical electrodynamics, an electron, moving round a circular orbit, has to radiate and as a resultat it will fall down on the nucleus. Bohr’s postulates, which had appeared soon after that, did not make clear this problem — the existance of allowed orbits also contradicted classical electrodynamics.
The quantum mechanics does not explain either, why electron in a stationary orbit does not radiate. The radiation of the electron, when it is going from one allowed orbit to another one, is not connected with the radiation of an accelerated electron anyhow. Atom can be in an excited state rather long, but an accelerated electron must radiate constantly. And what is more, in quantum mechanics the concept of movement with acceleration is not considered at all. The reference to in applicability of classical electrodynamics to atomic objects is quite unfounded. In fact, electron is held by forces, which submit to principles of electrodynamics. Then, why don’t these principles spread on motion of electrons?
Thus, the problem of the stability of atom remains to be unsolved up to now.
## 4 Do inertial coordinate systems exist?
In classical electrodynamics motion of charges is considered in inertial coordinate systems. These systems are defined in the following way: if any forces do not act on the body, or the sum of acting forces is equal to zero, then there are such coordinate systems, relative to which the body is moving without any acceleration. These systems are named inertial (The first Newton’s law). Inertial systems are moving uniformly and straigtforward relatively to each other, so the acceleration of the body has identical value in any of these systems, i. e. acceleration is absolute quantity.
But how can we realise inertial coordinate system in practice? Obviously, in order to realise it, it is necessary to take a body, on which any force does not act at all, or their sum must be equal to zero, and then to find a coordinate system, relatively to which this body will be moving without any acceleration. That system will be inertial.
But it is impossible to isolate a body from any action of external forces. Such isolation would mean, that there is only one body in the Universe. Thus, we can only try to bring the sum of external actions to zero. But how can we find out, that the sum of forces, acting on a body, is equal to zero? We can do it only in the following way — to define the acceleration on the body relatively to inertial coordinate system. If the acceleration is equal to zero, then the sum of the forces is equal to zero.
So we have got an exclusive circle; in order to determine an inertial coordinate system we have to choose a body so as the sum of acting forces on it is equal to zero; but we need to have an inertial coordinate system in order to determine equality of this sum of external forces to zero.
Thus, we have got a real difficulty when we try to realize the inertial coordinate system in practice.
In principle there is a possibility in classical electrodynamics to realize the inertial coordinate system. In fact, according to electrodynamics an accelerated charge must radiate. The intensity of radiation (in nonrelativistic case) is defined with an expression
$$I=\frac{2e^2w^2}{3c^3}$$
$`(2)`$
where e is a quantity of the charge, and w is its acceleration. Equations of electrodynamisc are true in inertial coordinate systems.
That is why, if a charge moves with acceleration, it radiates, if its acceleration is equal to zero, there is no any radiation. Radiation is an absolute quantity. It is impossible to create or to destroy it by any choose of any coordinate system. So the acceleration, which is undoubtedly connected with radiation, has to be an absolute quantity. Therefore, if some charge does not radiate, the system, which is connected with it, will be inertial.
In classical physics it is supposed, that equations of electrodynamics are true only in inertial coordinate systems. It is easy to make sure, that such a point of view is not founded enough. Indeed, the equations of Maxwell were received as a result of the generalization of experimentals. So far as there already were distinguished inertial coordinate systems, the equations of Maxwell were brought to these systems. But we do not have any grounds to affirm, that the experimentals, in which electromagnetic phenomens are studied and which are carried out in inertial coordinate systems and the results of the same experiments, which are carried out in an uniformly accelerated coordinate system, will be different. And if the results of these experiments in these systems are identical, in will mean, that the equations of Maxwell will be also true in uniformly accelerated systems.
Will Coulomb’s law be true in uniformly accelerated systems? In other words, if one can measure the interaction of motionless charges in a laboratory, which moves with uniform acceleration, then can this interaction be described by Coulomb’s law? In particular, will the forces of the interaction of the charges be on the same straight line? Obviously, that we can answer this question only on the grounds of the experimental. According to classical physics, Coulomb’s law is not true in uniformly accelerated coordinate system, for the charges, which are at rest in this system, move with acceleration in inertial coordinate systems, and the electromagnetic field of such charges is described by means of lagging potentials. That is why, Coulomb’s forces, with which the charges, which are at rest in the accelerated laboratory, interact, do not lie on the same straight line. But in the first place, this conclusion is not confirmed with experimentals, and in the second place, as we have seen, Coulomb’s interaction of the charges, moving with uniform acceleration and calculated within framework of classical electrodynamics leads to the violation of laws of conservation of energy and impulse.
So the fact, that there is not any radiation of the charge, is not a reason to suppose, that the coordinate system, which is connected with the charge, is inertial. That fact only indicates, that all electromagnetic processes in this system will proceed in the same way as in the system, which is connected with the Earth.
But if we are not able to realize the inertial coordinate system, then all the systems become equal. In this case the acceleration loses its absolute sense and becomes a quantity, which is as relative as the velocity. So as to determine the intensity of the radiation, the quantity of the acceleration in expression (2) has to be chosen so as to take additional terms into account.
It is remarkable, that the concept of the equality of all coordinate systems has been completely realised in general theory of relativity. And what is more, if we managed to realise the inertial coordinate system, the general theory of relativity would turn out groundless.
## 5 The electrodynamics and the Principle of equivalence
We will try to remove the contradictions, which have been noticed before. First we will consider a charge moving freely in a gravitational field. So as the principle of equivalence will be true, this charge must not radiate. This conclusion seems to be rather logical. Indeed, it is supposed in classical physics, that there is a straight line, transpiercing through all the Universe and a charge, which is moving along it with constant velocity and does not radiate. Any deflection of such movement must be accompanied by the radiation. In classical physics only a ray of light must be a straight line. But light declines in a gravitational field. So one can get a straight line only when there is no any gravitational fields, i. e. in a limited spheres of space.
In the general theory of relativity a concept of geodesic line is introduced. It is a generalized notion of a straight line. It is a trajectory of the motion of a body, on which any forces do not act, except the gravitational ones. Such movement of a body in the general theory of relativity is inertial. So we can expect, that the free movement of a charge in the gravitational field will not be accompanied by radiation. That means, that in the coordinate system, freely moving in a gravitational field, the equations of Maxwell will have the same form, as in the inertial system.
Let us consider the charge $`q`$, which is at rest in the inertial coordinate system with a constant gravitational field. Evidently, vector $`\stackrel{}{𝐸}`$ of such charge in any coordinate system will lie on a straight line, connected the charge and the point of observation. Indeed, let a charge $`e`$ lie at a distance of charge $`q`$. The force, acting on the charge $`e`$ from the side of the charge $`q`$, will be on a straight line, connected these charges. It is quite evidently, that the observer can move in any way, i. e. he can be in any coordinate system, but the direction of the force, acting on the charge $`q`$, will not change. Therefore field $`\stackrel{}{𝐸}`$ of the charge $`q`$ has to be described with expression
$$\stackrel{}{𝐸}=\frac{q\stackrel{}{𝑅}}{R^3}\frac{1\frac{v^2}{c^2}}{(1\frac{v^2}{c^2}\mathrm{sin}^2\theta )^{3/2}},$$
$`(3)`$
which is to be true for this charge in any coordinate system. Here $`R`$ is a distance from charge $`q`$ to the point of the observation, $`\theta `$ is an angle between the direction of the motion of the charge and radius-vector $`\stackrel{}{𝑅}`$, $`v`$ is momentary velocity of the charge relative to the observer. Expression (3) describes the field of a uniformly moving charge. Its applicability in this case evident. Indeed, let there be an accelerated system and an inertial system, accompanying it. It is clear, that field $`\stackrel{}{𝐸}`$ in the both systems will be identical at any moment for the observers.
For the charge $`q`$, which is at rest in a constant gravitational field, does not radiate then so as to fulfil the principle of equivalence, the charges, moving with uniform acceleration or uniformly moving round a circle, must not radiate either. Field $`\stackrel{}{𝐸}`$ of such charges must also be described with expression (3), which is true for these charges in any coordinate system. (We consider only those charges, the quantity of the acceleration of which does not change during the period of observation, i. e. we except transitional processes). That means, that in the coordinate systems, which moves with an uniform acceleration or uniformly moves round a circle, Maxwell’s equations has the same form as in the inertial systems. Under this assertion we mean the following: in the systems, moving with an uniform acceleration or uniformly moving round a circle, all electromagnetic processes will proceed in the same way as in the inertial systems with constant gravitational field.
But if the field $`\stackrel{}{𝐸}`$ of the charges, moving with uniform acceleration or uniformly moving round a circle, is described with expression (3), there is no any violation of the laws of conservation by calculation of its interaction. Indeed, in this case the forces, which the charges, which were considered before, interact with, will be equal by the quantity and will be contrary directed.
The lack of the force of radiation friction during the motion with uniform acceleration becomes clear. For such charge does not radiate, so there will not be any violation of the balanse of energy.
How do the changes, which have been brought into electrodynamics, accord with the experimentals? It is the first question. And the second one — if the charge radiates, what value of the acceleration should we put down into the expression(2) to define intensity of the radiation.
Let us answer the second question first. We will consider a charge, oscillating in accordance with harmonic law relatively to inertial system $`K`$. Under inertial system $`K`$ we mean here coordinate system, connected with the Earth. By such movement of the charge contradictions do not appear. So the charges, the motion of which can be brought to harmonic oscillations in system $`K`$, we will describe within the framework of classical electrodynamics. Now let the charge move straightforward with acceleration $`w(t)`$ in system $`K`$. If function $`w(t)`$ is periodic, we can expand it in trigonometric row
$$w(t)=w_0+\underset{n=1}{\overset{\mathrm{}}{\Sigma }}(b_n\mathrm{cos}\omega nt+c_n\mathrm{sin}\omega nt),$$
$`(4)`$
where $`w_0`$ — constant of function $`w(t)`$, which is accorded with the movement of the charge with uniform acceleration and $`b_n`$ and $`c_n`$ — amplitudes of harmonic $`n`$. I. e. the function $`w(t)`$ can be expanded in a form of a sum of two items $`w(t)=w_0+w_1(t)`$. Through $`w_1(t)`$ the sum of harmonics is designated in (4). For we think, that the uniformly accelerated charge does not radiate, then in the expression (2) for intensity of the radiation only value $`w_1(t)`$ has to be put in as an acceleration. If function $`w(t)`$ is not periodic, it has to be expanded into Fourier’s integral. In this case the steady component $`w_0`$ is equal to zero. So we will put the complete value of the acceleration $`w(t)`$ into expression (2), i. e. the radiation of the charge will be defined in the same way as in classical dynamics. For in practice a charge can move straightforward only during a short period of time, its acceleration cannot be a periodic function. The radiation of such a charge therefore will completely correspond to parameters, which have been received within the framework of electrodynamics.
During the movement of a charge in a closed trajectory with a constant period of revolution, the value of its acceleration $`w(t)`$ can be also expanded in a trigonometric row (4). The steady component $`w_0`$ corresponds to the acceleration of the uniform moving round a circle here. So when we calculate the intensity of the radiation in expression (2), as an acceleration we will put in only the value $`w_1(t)`$. If function $`w(t)`$ is not periodic, we will put in the complete value of the acceleration.
The fact, that the charge, uniformly moving round a circle, does not have any radiation, does not contradict the existence of the synchrotron radiation. Indeed, the quantity of the acceleration of electrons in cyclical accelerators is not a periodic function. If we consider the movement of electron in the accelerator as periodic, we will receive the following resultat. The intensity of radiation I<sub>1</sub> on frequency of circulation $`w_1`$ (the first harmonic), calculated within the framework of classical electrodynamics, is proportionate to
$$w_0^2+\frac{1}{2}(b_1^2+c_1^2).$$
The intensity of the $`nth`$ harmonic is proportionate to $`{\displaystyle \frac{1}{2}}(b_n^2+c_n^2).`$ If we suppose, that the charge, uniformly moving round a circle, does not radiate, it will only bring us to the fact, that only the intensity of the first harmonic will decrease. In this case it must by proportionate to $`{\displaystyle \frac{1}{2}}(b_1^2+c_1^2`$). The intensity of the rest harmonics will not change. As it is known from the electrodynamics, the great part of the radiation is concentrated in the range of frequencies $`ww_1\left({\displaystyle \frac{\epsilon }{mc^2}}\right)^3`$, where $`\epsilon `$ — is an energy of electron. On frequency $`w_1`$ a small part of energy is radiated. So when $`\epsilon =50`$ Mev $`I_1/I10^8`$, where $`I`$ is a complete intensity. It is necessary to take notice of the fact, that because of the influence of the conductive surfaces (the sides of the vacuum chamber of an accelerator) the intensity of the low frequency part of the radiation will decrease approximately $`\left({\displaystyle \frac{r}{d}}\right)^2`$ times, where $`r`$ is a radius of the orbit, and $`d`$ is a distance from the electron beam to the conductive surface.
So it is clear, that it is practically impossible to notice such deflection in synchrotron radiation. We need a special experiment, which could help us to test the intensity of the radiation on frequency of revolution. When the charge uniformly moves round a circle, there must not be any radiation at all. In such way electron in atom can move in an stationary orbit. And as it is known, such electron does not radiate. This fact confirms the results we have received. On the other hand, we have received the explanations of the stability of Rutherford’s atom, i. e. why electron does not radiate in a circular orbit.
## 6 Supplement
As it is known, the equations of physics, which are written in general covariance form, automatically satisfy the principle of equivalence. Therefore the general covariance expression for the intensity is not to violate the principle of equivalence. Let us receive such an expression. The charge is at rest in the inertial coordinate system. Its energy, which was radiated during period $`dt`$, is equal to
$$d\epsilon =\frac{2e^2w^2}{3c^2}dt.$$
$`(5)`$
Complete radiated impulse in this coordinate system is equal to zero
$$d\stackrel{}{P}=0.$$
$`(6)`$
In order to pass to an arbitary inertial coordinate system we will write down expressions (5) and (6) in four-dimensional form
$$dP^i=\frac{2e^2}{3c}\frac{du^k}{ds}\frac{du_k}{ds}u^ids.$$
$`(7)`$
So as to pass to general covariance expression for $`dP^i`$ in (7) we will substitute usual differentiation for covariance. Then we will receive
$$dP^i=\frac{2e^2}{3c}\left(\frac{du^k}{ds}+\mathrm{\Gamma }_{lm}^ku^lu^m\right)\left(\frac{du_k}{ds}\mathrm{\Gamma }_{kl}^iu_iu^m\right)u^ids.$$
$`(8)`$
Here $`\mathrm{\Gamma }_{lm}^k`$ are symbols of Christoffel, which are
$$\mathrm{\Gamma }_{il}^k=\frac{1}{2}g^{kn}\left(\frac{g_{nl}}{x^m}+\frac{g_{nm}}{x^l}\frac{g_{lm}}{x^n}\right),$$
$`(9)`$
where $`g^{lm}`$ is a fundamental tensor. Indexes take values 0, 1, 2, 3. Expression (8) is true in any coordinate system. For a body, moving freely in the gravitational field, it is
$$\frac{du^k}{ds}+\mathrm{\Gamma }_{lm}^ku^lu^m=0.$$
Therefore, the charge, moving in such way, does not radiate. We have received this conclusion before, when we were according electrodynamics with the principle of equivalence.
Now we will consider the charge, which is motionless in inertial coordinate system $`K`$ with constant gravitational field. The acceleration of such charge in system $`K`$ is zero, and the components of the 4-velocity are
$$u_0=1,u_1=u_2=u_3=0.$$
So it follows from (8), that the intensity of radiation of the charge in this case is
$$\frac{d\epsilon }{dt}=\frac{2e^2c}{3}\mathrm{\Gamma }_{\mathrm{0\hspace{0.17em}0}}^i\mathrm{\Gamma }_{i\mathrm{\hspace{0.17em}0}}^0.$$
For a weak gravitational field components of tensor $`g_{em}`$ are
$$g_{\mathrm{1\hspace{0.17em}1}}g_{\mathrm{2\hspace{0.17em}2}}g_{\mathrm{3\hspace{0.17em}3}}1,g_{\mathrm{0\hspace{0.17em}0}}=1+\frac{2\phi }{c^2},g_{ik}0(ik),$$
where $`\phi `$ is a gravitational potential. Taking the statical nature of the gravitational field into consideration, i. e. $`{\displaystyle \frac{g_{kl}}{x^0}}=0`$, and that the components of tensor $`g^{ik}`$ in this case have approximately the same values, as the components of tensor $`g_{ik}`$, we will receive, when $`\alpha `$=1, 2, 3.
$$\mathrm{\Gamma }_{00}^\alpha =\mathrm{\Gamma }_{\alpha 0}^0=\frac{1}{2}\frac{g_{00}}{x^\alpha }.$$
Therefore, the intensity of radiation is equal to
$$\frac{d\epsilon }{dt}=\frac{2e^2}{3c^3}\frac{\phi }{x^\alpha }\frac{\phi }{x_\alpha }.$$
$`(9)`$
For $`\dot{\stackrel{}{v}}=\phi `$, and acceleration $`\dot{\stackrel{}{v}}`$ in this case is the accelaration due to gravity $`g`$, then (9) will have the following form
$$\frac{d\epsilon }{dt}=\frac{2e^2}{3c^3}g^2.$$
Thus, the charge, which is immovable on the Earth, must radiate. The intensity of its radiation is identical to the intensity of the charge, moving with acceleration $`g`$ in empty space at a great distance away from any gravitating bodies. So the principle of equivalence is formally fulfilled. It is naturally for covariance expression (8), though it is clear, the charge, which is at rest in a gravitational field, does not charge. |
warning/0003/quant-ph0003094.html | ar5iv | text | # Quantum Communication with Correlated Nonclassical States
## I Introduction
Principal motivations for the investigation of manifestly quantum or nonclassical states of the electromagnetic field have been their possible exploitation for optical communication$`^{\text{[3, 4, 5]}}`$ and for enhanced measurement sensitivity.$`^{\text{[6]}}`$ For example, relative to a coherent state, the reduced quantum fluctuations associated with squeezed and number states offer potential for improving channel capacity in the transmission of information.$`^{\text{[5]}}`$ Squeezed states of light have been widely employed to achieve measurement sensitivity beyond the standard quantum limits in applications such as precision interferometry,$`^{\text{[7]}}`$ the detection of directly encoded amplitude modulation,$`^{\text{[8]}}`$ atomic spectroscopy,$`^{\text{[9]}}`$ and quantum noise reduction in optical amplification.$`^{\text{[10]}}`$ Likewise, nonclassical correlations for the amplitudes of spatially separated beams have been exploited in diverse situations, including demonstrations of the EPR paradox for continuous variables,$`^{\text{[11]}}`$ of quantum nondemolition detection (QND),$`^{\text{[12, 13]}}`$ and of a quantum-optical tap.$`^{\text{[14]}}`$
Within the broader setting of quantum information science (QIS), there has been growing interest and important progress concerning the prospects for quantum information processing with continuous quantum variables, including universal quantum computation,$`^{\text{[15]}}`$ quantum error correction,$`^{\text{[16, 17, 18]}}`$ and entanglement purification.$`^{\text{[19, 20]}}`$ Theories for quantum teleportation of continuous quantum variables in an infinite dimensional Hilbert space have been developed,$`^{\text{[21, 22, 23, 24]}}`$ including for broad bandwidth teleportation$`^{\text{[25]}}`$ and for teleportation of atomic wavepackets.$`^{\text{[26]}}`$ This formalism has also been applied to super-dense quantum coding.$`^{\text{[27]}}`$ On an experimental front, these developments in QIS led to the first bona fide demonstration of quantum teleportation, which was carried out by exploiting nonclassical states of light in conjunction with continuous quantum variables.$`^{\text{[28, 29]}}`$
Against this backdrop the focus of attention in this article is optical communication in two channels with quantum correlated light fields and the associated quadrature amplitudes.$`^{\text{[30]}}`$ The goal is to explore the extension of quantum cryptography from the usual setting of discrete variables as pioneered by C. Bennett and colleagues$`^{\text{[31]}}`$ (e.g., photon polarization as in the experiments of Refs.) into the realm of continuous quantum variables (e.g., the complex amplitude of the electromagnetic field). Apart from our work, several related schemes for quantum cryptography based upon continuous variables have recently been analyzed, including a single-beam scheme with squeezed light$`^{\text{[36]}}`$ as well dual-beam schemes with shared entanglement.$`^{\text{[37, 38]}}`$ However, we stress at the outset that neither for our scheme nor for any of these other protocols, can any claim about absolute security be made. Rather, we suggest that these protocols (and suitable extensions thereof) are worthy candidates for more detailed analyses. Such an undertaking would involve various important matters of principle as well as practice for continuous quantum variables, and might hopefully lead to security proofs such as have recently emerged in the case of discrete variables.$`^{\text{[39, 40, 41]}}`$
As illustrated in Figure 1, the basic idea in our scheme is to construct a “transmitter” which combines a coherent signal of amplitude $`ϵ/t`$ (the “message”) with the large fluctuating fields generated in nondegenerate optical parametric amplification (the “noise”).$`^{\text{[30]}}`$ The “message” and the “noise” are superimposed at mirror $`M`$ with transmission coefficient $`t`$ $`<<`$ 1. Note that although each of the two transmitted beams along channels ($`A,B`$) has large phase insensitive fluctuations that are individually indistinguishable from a thermal source,$`^{\text{[42]}}`$ the quadrature-phase amplitudes of the two-beams can be quantum copies of one another,$`^{\text{[11, 43]}}`$ and in fact form an entangled EPR state.$`^{\text{[44, 45]}}`$ Hence proper subtraction of the photocurrents at the “receiver” can result in the faithful reconstruction of the encoded “message” even though the signal-to-noise ratios $`R_j`$ ($`j`$=$`A,B`$) during transmission are individually much less than one. Indeed, in a lossless system with large parametric gain, the signal-to-noise ratio of the reconstructed message $`R_t`$ can approach the signal-to-noise ratio $`R_0`$ of the original message ($`ϵ^2/t^2`$), which was written in the transmitter as a coherent state before the mirror $`M`$ in Figure 1. Note that the individual Channels $`(A,B)`$ have a high degree of immunity to unauthorized interception since the signal-to-noise ratios $`R_{A,B}`$ in these channels are each very small. Furthermore, any attempt to extract information from the $`(A,B)`$ channels will reveal itself either by a decrease in R<sub>t</sub> (classical extraction) or by an increase in the fluctuations of the orthogonal quadrature amplitude (quantum extraction).
In addition to achieving a faithful reconstruction of the message transmitted through M to the receiver, note that the scheme of Figure 1 also preserves to a high degree the signal-to-noise ratio for the original message beam that reflects from M. More specifically, for high gain and for losses dominated by the transmission coefficient of M, the signal-to-noise ratio R<sub>r</sub> for the reflected beam can approach R<sub>0</sub> for the original message. In this limit, we then have that the information transfer coefficient T $``$ (R<sub>r</sub>+R$`{}_{t}{}^{})`$ / R$`{}_{0}{}^{}`$ 2, where $`0`$ T $``$ 1 for classical devices and 1 $`<`$ T $``$ 2 for manifestly quantum or nonclassical situations.$`^{\text{[14]}}`$ Hence the scheme depicted in Fig.1 acts as a quantum optical tap in the fashion originally discussed by Shapiro.$`^{\text{[46]}}`$ It provides a received (or “tapped”) message with a signal-to-noise ratio equal to that of the input ($`R_t/R_0)1`$, while simultaneously transmitting an output field with signal-to-noise ratio equal to that of the input $`(R_r/R_0)1`$.
Of course similar schemes for two-channel communication can be implemented with correlated classical noise sources (i.e., thermal light), each with large fluctuations which “hide” the message $`ϵ`$ during transmission. However, with classical sources of whatever type, only excess fluctuations can be subtracted; the quantum fluctuations at the vacuum-state level will remain unchanged and will enforce a noise floor for information transmission and extraction. For the case illustrated in Figure 1, this noise floor for the message at the receiver is given by the sum of independent vacuum fluctuations from fields in channels ($`A,B`$) and sets a fundamental noise level of “2” (with “1” as the individual vacuum-state limits for the two channels). Here we adopt the usual convention for the demarcation between classical and nonclassical correlations in terms of the behavior the Glauber-Sudarshan phase-space function.$`^{\text{[43]}}`$ Hence for the case illustrated in Figure 1 but with classical input fields, the signal-to-noise ratio $`R_t^{}`$ for the detected message at the receiver is given by $`R_t^{}|ϵ|^2R_t`$. In fact, for classical inputs, we have that $`T^{}R_t^{}+R_r^{}1`$, and the system no longer functions as a quantum optical tap. Furthermore, the individual channels ($`A,B`$) are not protected from unauthorized eavesdropping, since information can be extracted from these channels with impunity for classical noise much greater than the vacuum-state limit.
Apart from these considerations related to secure communication and quantum optical tapping, the configuration of Figure 1 can also be viewed as a means to realize super-dense quantum coding$`^{\text{[47]}}`$ for continuous quantum variables.$`^{\text{[27]}}`$ Here, the message $`ϵ/t`$ is again encoded at the mirror $`M`$, but now in a single channel corresponding to one component of the entangled EPR state (e.g., channel $`A`$). This combination of the message and the fluctuations from one component of the NOPA are transmitted to the receiving station where they are combined with the second component of the entangled output of the NOPA that has been independently transmitted (e.g., along channel $`B`$). The signal is then decoded by combining the outputs of the two channels in a fashion similar to that shown in Figure 1 as discussed in more detail in Ref.. The principal distinctions between this dense coding scheme and the aforementioned dual channel arrangement are (1) the message is encoded in a single component of the entangled EPR beam instead of symmetrically in both and (2) the received beams from paths $`(A,B)`$ must be physically recombined, with the phases of the local oscillators $`(A,B)`$ at the receiving station offset by $`\frac{\pi }{2}`$. Recall that for dense coding in its canonical form,$`^{\text{[47]}}`$ no signal modulation is applied to the second (i.e., channel $`B`$) component of the entangled state, so that it carries no information by itself.
In subsequent sections of this paper, we describe in more detail the implementation of this general discussion about quantum communication with correlated nonclassical fields. In our experiment, we have been able to demonstrate an improvement in signal-to-noise ratio by a factor of 2.1 over that possible with any classical source (that is, 10 $`log[R_t/R_t^{}]=3.2`$ dB) and have succeeded in suppressing the noise of the difference photocurrent $`i_{}i_Ai_B`$ below that associated with the vacuum fluctuations of even a single beam, thus making possible transmission with $`|ϵ|^2`$ $`<1`$. Quantum dense coding would thereby be enabled with the aforementioned changes in the overall experimental protocol. We conclude with a discussion of possible extensions for enhanced security against unauthorized eavesdropping.
## II Implementation by Nondegenerate Parametric Amplification
As illustrated in Figure 1, correlated nonclassical states for our work are generated by a nondegenerate optical parametric amplifier (NOPA) that produces orthogonally polarized but frequency degenerate signal and idler beams for channels ($`A,B`$). We emphasize that these beams represent a realization of the entangled state originally discussed by Einstein, Podolsky, and Rosen.$`^{\text{[44, 45]}}`$ For the original EPR state, there exist perfect correlations both in position and momentum for two massive particles. In the optical case, the quadrature amplitudes of the electromagnetic field play the roles of position and momentum with a finite degree of correlation for finite NOPA gain, as has been experimentally demonstrated$`^{\text{[11]}}`$ and exploited to realize quantum teleportation.$`^{\text{[28]}}`$
A coherent-state “message” of total amplitude $`ϵ/t`$ is encoded in equal measure onto these entangled EPR beams by orienting its polarization at 45 with respect to the signal and idler polarizations at the mirror $`M`$ of Figure 1. To obtain a quantitative statement of the performance of this system, we must include the finite gain of the amplifier as well as various passive losses, which together limit the degree of correlation that can be exploited for communication. Following the analysis of Ref., we find that the SNR $`R_j(\mathrm{\Omega })`$ for the individual signal and idler photocurrents for propagation and detection in the presence of overall channel efficiency $`\xi `$ is given by $`R_j(\mathrm{\Omega })=\xi ϵ^2/2G_q(\mathrm{\Omega })`$, where $`G_q(\mathrm{\Omega })`$ is the detected quantum-noise gain of the amplifier which can be determined experimentally from measurements of the spectral densities $`\mathrm{\Psi }_{A,B}(\mathrm{\Omega })`$ for the fluctuations of photocurrents for signal and idler beams alone at either detector. Relative to the frequency of the optical carrier determined by the down-conversion process in the NOPA, the frequency $`\mathrm{\Omega }`$ specifies the Fourier components of the quadrature-phase amplitudes of signal and idler fields as well as of the coherent field $`ϵ`$.$`^{\text{[43]}}`$ Note that $`\xi `$ ($`0\xi 1`$) incorporates the cavity escape efficiency for our NOPA, the propagation efficiency from the NOPA to the detectors, and the homodyne and quantum efficiencies of the balanced detectors themselves.$`^{\text{[11]}}`$
Although the individual fluctuations for channels ($`A,B`$) give rise to a level $`G_q(\mathrm{\Omega })>1`$, (that is, greater than the vacuum-state limit of either beam alone), these large fluctuations are correlated in a nonclassical manner and hence can be eliminated by proper choice of the quadrature amplitudes detected at ($`A,B`$). As shown in Ref., there is a continuous set of such amplitudes with minimum variance for their difference requiring only that the quadrature-phase angles ($`\theta _A,\theta _B)`$ satisfy $`\theta _A+\theta _B=2p\pi `$ ($`p`$ = integer). Denoting one such pair by ($`X_A,X_B`$), we have that
$`(X_A(\mathrm{\Omega })X_B(\mathrm{\Omega }))(X_A(\mathrm{\Omega }^{})X_B(\mathrm{\Omega }^{}))`$
$$=V_{}(\mathrm{\Omega })\delta (\mathrm{\Omega }+\mathrm{\Omega }^{}),$$
(1)
where $`V_{}(\mathrm{\Omega })`$ is a variance which quantifies the degree of correlation between $`(X_A,X_B`$). Explicit expressions for both $`V_{}(\mathrm{\Omega })`$ and $`G_q(\mathrm{\Omega })`$ are given in Ref.. For propagation and detection in the presence of loss, we introduce the quantities ($`V_{}^d,G_q^d`$) which refer to the variance and quantum noise gain for fictitious fields having propagated with total loss $`(1\xi )`$, where the spectral density of the photocurrent fluctuations $`\mathrm{\Phi }_{}(\mathrm{\Omega })`$ is proportional to $`V_{}^d(\mathrm{\Omega })`$. Hence, the SNR $`R_d`$ for detection of the message via $`i_{}`$ is given by $`R_d=2\eta ϵ^2/V_{}^d(\mathrm{\Omega })`$, where $`\eta `$ accounts for the propagation and detection efficiency for the message from the mirror M to the photocurrent $`i_{A,B}`$. Without discussing the general case, here we note simply that for efficient propagation and detection with $`(1\xi )<<1`$ and for near threshold operation with (analysis frequency $`\mathrm{\Omega })<<`$ (cavity linewidth $`\mathrm{\Gamma })`$, then $`G_q^d(\mathrm{\Omega })1+\frac{1}{2}(\mathrm{\Gamma }/\mathrm{\Omega })^2\xi `$, while $`V_{}^d(\mathrm{\Omega })2(1\xi )<<1`$, so that $`R_d(\mathrm{\Omega })\eta ϵ^2/(1\xi )`$. Hence in the ideal case with $`\eta `$1 and $`\xi (1|t|^2)`$, with t as the amplitude transmission coefficient of mirror M, we find that the reconstructed message is recovered with the same SNR with which it was originally encoded (namely $`R_dϵ^2/t^2`$), while the fluctuations in the individual channels become arbitrarily large ($`R_{A,B}\mathrm{\Omega }^2ϵ^2/\mathrm{\Gamma }^20`$ for $`\mathrm{\Omega }/\mathrm{\Gamma }0`$).
As for the performance as an optical tap, note that the transfer coefficient associated with the detected message at the receiver and with the reflected output field is given by $`T_d=(R_d+R_r)/R_0`$, where $`R_d`$ is related to $`R_t`$ by way of the propagation and detection efficiency $`\eta `$ from M to the photocurrents at the receivers. In the present case, we have that
$$T_d=\frac{|r|^2}{U_{}^r(\mathrm{\Omega })}+\frac{2\eta |t|^2}{V_{}^d(\mathrm{\Omega })},$$
(2)
with $`|r|^2`$ as the reflectivity of mirror M $`(|r|^2+|t|^2`$=1) and $`U_{}^r(\mathrm{\Omega })`$ as the variance of the reflected field. Hence in the ideal case with $`\xi (1|t|^2)`$, with V$`{}_{}{}^{d}(\mathrm{\Omega })2(1\xi )`$ and with $`U_{}^r(\eta )=|r|^2`$, we have that $`R_dR_t`$ and $`T_d2`$. Thus in addition to providing large quantum fluctuations for secure transmission, the system also acts as a quantum optical tap with a nearly ideal transfer coefficient $`T`$.
In fact the system can be considered as a realization of the scheme for quantum tapping that was originally suggested by Shapiro.$`^{\text{[46]}}`$ To see this more clearly, recall that the projection of signal and idler fields along the $`45^{}`$ polarization direction of the message beam results in a squeezed field.$`^{\text{[43]}}`$ Hence, from the perspective of Ref., we are “tapping” the original message field by injecting squeezed light into the normally open (or vacuum) port of mirror M. The use of the output of a nondegenerate parametric amplifier allows us subsequently to decompose this squeezed plus coherent field into individually noisy signal and idler fields at polarizer P for transmission.
## III Experimental Setup and Results
The general scheme for our experimental implementation of these ideas is shown in Figure 1, where frequency degenerate but orthogonally polarized signal and idler beams are generated by Type II down-conversion in a subthreshold optical parametric oscillator formed by a folded cavity containing an $`a`$-cut crystal of potassium titanyl phosphate (KTP) that provides noncritical phase matching at 1.08$`\mu `$m. The crystal is 10mm long, is anti-reflection coated for both 1.08$`\mu `$m and 0.54$`\mu `$m, and has a measured harmonic conversion efficiency of $`6\times 10^4`$/W (single-pass) for this geometry. The total intracavity passive losses at 1.08$`\mu `$m are 0.3% and the transmission coefficient of mirror M1 is 3%. The amplifier is pumped by green light at 0.54 $`\mu `$m generated by external frequency doubling of a frequency-stabilized, TEM<sub>00</sub>-mode Nd:YAP laser.$`^{\text{[48]}}`$ The subthreshold oscillator acts as a narrow-band amplifier (NOPA) which is locked to the original laser frequency with a weak counter-propagating beam. Simultaneous resonance for the orthogonally polarized signal and idler fields is achieved by adjusting the temperature of the KTP crystal around 60C with milliKelvin precision. The pump field at 0.54 $`\mu `$m is itself resonant in a separate and independently locked build-up cavity (enhancement $``$ 5x).
As we have demonstrated in our previous experiments, $`^{\text{[11]}}`$ the orthogonally polarized signal and idler fields generated by the NOPA individually are fields of zero mean values and exhibit large phase insensitive fluctuations. It is in the midst of this noise that we now hide a “message”, with this coherent field being combined with the signal and idler fields at the highly reflecting mirror $`M`$ shown in Figure 1 ($`(1t^2)0.99`$). The coherent beam is injected at 45 with respect to signal and idler polarizations and is frequency shifted by $`\mathrm{\Omega }_0/2\pi =1.1`$MHz (single-side band) from the primary laser frequency with the help of a pair of acoustooptic modulators, which are gated “on” and “off” to provide information encoded for transmission. The noisy but correlated signal and idler beams together with the coherent information are then separated by a polarizer $`P`$, transmitted independently over the two channels ($`A,B`$), and then directed to two separate balanced homodyne detectors for measurements of their individual quadrature-phase amplitudes and their mutual correlations. The local oscillators for the two balanced homodyne detectors originate from the laser at 1.08$`\mu `$m; their phases can be independently controlled by mirrors mounted on piezoelectric transducers. The spectral densities of the photocurrents for the two channels ($`A=`$ signal, $`B=`$ idler) are defined by
$$\mathrm{\Psi }_{A,B}(\mathrm{\Omega })=<i_{A,B}(t)i_{A,B}(t+\tau )>e^{i\mathrm{\Omega }\tau }𝑑\tau $$
(3)
and are recorded by a RF spectral analyzer, as is the spectral density
$$\mathrm{\Phi }_{}(\mathrm{\Omega })=<i_{}(t)i_{}(t+\tau )>e^{i\mathrm{\Omega }\tau }𝑑\tau $$
(4)
for the difference photocurrent $`i_{}i_Ai_B`$.
In Figures 2 and 3 we present results from a series of measurements of these various spectral densities. First of all, in Figure 2a, trace i gives the spectral density $`\mathrm{\Psi }_A`$ for channel A alone with an injected “message” and with the amplifier turned on to generate large $`(`$ 7 dB) phase insensitive noise above the vacuum-state level $`\mathrm{\Psi }_{0A}`$ (indicated by a dashed line in Figure 2) for the signal beam. A similar trace is obtained for the spectral density $`\mathrm{\Psi }_B`$. By contrast, trace ii in Figure 2a gives the spectral density $`\mathrm{\Phi }_{}`$ for the difference photocurrent $`i_{}`$, with the phases of the local oscillators adjusted for minimum noise and maximum coherent signal. In this trace, the coherent message that was completely obscured in trace i emerges with high signal-to-noise ratio. Note that in trace ii the correlated quantum fluctuations for signal and idler fields are subtracted to approximately 0.4 dB below the vacuum-noise level $`\mathrm{\Psi }_{0A}`$ of the signal beam alone (and likewise for the idler), indicating an improvement in SNR over a conventional single-channel communication scheme with a classical light source.
To complete the discussion, we present in Figure 2b results obtained with the amplifier turned off (that is, uncorrelated vacuum-state inputs for signal and idler fields which are combined with the coherent “message” information at mirror $`M`$). Trace i shows the result for the signal beam alone ($`\mathrm{\Psi }_A`$), where again the noise floor $`\mathrm{\Psi }_{0A}`$ is from the vacuum fluctuations of the signal beam; a similar trace is obtained for the idler beam $`(\mathrm{\Psi }_B)`$. Trace ii gives the corresponding result for $`\mathrm{\Phi }_{}`$ for the combined signal and idler photocurrents when the amplifier is off. Note that this trace represents the best possible SNR with which the encoded information can be recovered when correlated classical noise sources are employed since here the (uncorrelated) vacuum fluctuations of signal and idler beams set an ultimate noise floor 3 dB above $`\mathrm{\Psi }_{0A}`$ (that is, $`\mathrm{\Phi }_0=2\mathrm{\Psi }_{0A}`$).$`^{\text{[49]}}`$ On comparing traces ii in Figures 2a and 2b, we see that the correlated quantum fluctuations of signal and idler fields brought about by parametric amplification result in an improvement in SNR of 3.2 dB relative to that possible with classical noise sources.
The improvement in SNR with correlated quantum fields over classical fields in our two-channel communication scheme can be of utility especially when the message is so weak that the SNR is poor for transmission with correlated classical sources (that is, for the case where vacuum noise dominates the encoded message). This situation is illustrated in Figure 3, where we plot $`\mathrm{\Phi }_{}`$ for the two cases without (trace i) and with (trace ii) correlated quantum fields.$`^{\text{[49]}}`$ Relative to Figure 2, here the coherent beam has been attenuated resulting in a smaller SNR for the “message”. Indeed in trace i, this information is “buried” by the vacuum noise $`\mathrm{\Phi }_0`$ associated with independent vacuum fluctuations in channels $`A`$ and $`B`$; recovery of the encoded information is poor. On the other hand, as shown in trace ii, when correlated quantum fields are employed, there is a reduction in the noise floor by more than 3 dB which makes possible improved recovery of the encoded information, with the recovery here limited by losses in propagation and detection.$`^{\text{[11]}}`$
As for the actual performance with respect to optical tapping, our system falls far short of the projected possibilities discussed in the preceding section because of an unfortunate mismatch between the transmissivity $`|t|^2`$ for mirror M and the overall system efficiency $`\xi `$. In quantitative terms, recall that the transfer coefficient $`T`$ for encoding information from the input beam to the reflected and transmitted beams at M is given by $`T(R_r+R_t)/R_0`$ whereas the transfer coefficient for the detected message photocurrent and the reflected signal field is $`T_d(R_r+R_d)/R_0`$ as given explicitly in Eq. (2). For the propagation and detection efficiencies in our experiment $`(\xi `$ 0.65 and $`\eta `$ 0.75), these transfer coefficients are optimized for mirror transmission $`|t|^2`$ 0.5 for M. In our arrangement we have instead $`|t|^2=0.01`$, with the inferred result that $`T_d1.02`$, which is only marginally in the quantum domain.
In the experiment described here, the receiver uses a local oscillator (LO) that originates from the fundamental frequency of the same laser that generated the pump beam for the NOPA. This LO is necessary for proper detection of the quadrature amplitudes of the nonclassical beams and of the message, since it provides a phase reference that follows phase fluctuations of the NOPA’s pump beam. In practice, as the stability of the available lasers improve, one should consider schemes for which the measurement is carried out with nominally independent lasers for the LO and for the source. For example, one might employ a stabilized laser diode as a reference to phase lock lasers both at the sender and at the receiver, where the laser diode could be widely distributed through optical fibers. Alternatively, Ralph has analyzed a scheme in which the local oscillators are transmitted and recovered as part of the overall protocol.$`^{\text{[37]}}`$
## IV Comparison with Other Dual Beam Schemes
It is perhaps obvious that the degree of immunity to interception for a two channel scheme such as we have discussed is related to the degree of excess fluctuations for each individual beam. For the demonstration in Ref., the excess noise used to “hide” the encoded information in each beam comes from some artificial unrelated source. Unfortunately such uncorrelated excess fluctuations also add noise to the coincidence signal in the recovery of the “message,” even though the added noise scales differently as a function of photon number for single-beam measurements (linearly) and for dual beam measurements (quadratically). Hence larger background noise which better “hides” the encoded information also brings larger added noise in the extraction of the “message.” Because of the quadratic dependence on the total photon number for the extra noise added in coincidence detection, this scheme is best suited to low light level transmission, as demonstrated in the pioneering experiment by Hong et al.$`^{\text{[50]}}`$
The situation is quite different for the quadrature-phase amplitudes of the correlated signal and idler fields generated by the NOPA. As the NOPA is pumped harder and the threshold for parametric oscillation is approached, the gain of the amplifier increases, as do the excess fluctuations of the signal and idler fields. However, the correlation between the fluctuations of the signal and idler beams also improves, giving rise to even better SNR for the recovered signal. The key point is that the large fluctuations in the signal and idler beams needed for immunity to interception are intrinsic and do not add extra noise to the recovered signal but, on the contrary, serve to reduce the noise in $`i_{}`$ as the gain of the amplifier increases. In the end, the SNR for the recovered message is arbitrated by the imperfect correlation resulting from finite gain and from passive losses in propagation and detection. On the other hand, this dependence provides a powerful means to detect eavesdropping because unauthorized extraction of signal or idler fields from channels $`A`$ or $`B`$ results in a reduction of the detected correlation and hence an increase in the noise floor of the recovered message. Note that unauthorized extraction of information from both channels by way of a quantum optical tap$`^{\text{[46]}}`$ or a quantum nondemolition measurement$`^{\text{[13]}}`$ can likewise be detected because of the unavoidable increase of fluctuations for the orthogonal quadrature-phase amplitudes ($`\theta _{A,B}+\pi /2`$) of the two channels. Furthermore, these quantum eavesdropping schemes can be defeated in large measure by random switching of the phases of the message, signal, and idler beams as discussed below.
Our system also offers advantages with respect to the (classical) digital Vernan cipher, where a message is decomposed in two correlated random signals and transmitted over two one-way channels. Although this system seems to be similar to ours in the sense that is also secure provided the eavesdropper has access to one channel only, the situation is different if the eavesdropper can split a small fraction of both channels since in the classical case, this can be done without the knowledge of the receiver. However, in our system the eavesdropper cannot choose arbitrarily the reflectivity of any “beamsplitter” used for extraction from the two channels since in the quantum case, the fraction of the beams extracted should be big enough so that the signal-to-noise ratio for the intercepted message is greater than one. But if this is the case, then unavoidable extra “noise” added to the transmitted beams by the open port of the “beamsplitter” degrades the signal-to-noise ratio of the message at the legitimate receiver, thus revealing the unauthorized intervention during transmission.
One might attempt to circumvent this difficulty by employing a quantum extraction procedure, such as quantum nondemolition detection$`^{\text{[13]}}`$ of the quadrature amplitudes in Channels $`(A,B)`$. Although the signal-to-noise ratio $`R_d`$ at the receiver would not in this case be degraded by an ideal eavesdropper, the unauthorized intervention could nonetheless be discovered because of the injection of large fluctuations (“back-action” noise) in the quadrature orthogonal to that in which signal information is stored, as previously noted.
## V Extensions via Random Phase Switching
One way an eavesdropper Eve could access the signal and idler beams without the knowledge of the legitimate receiver is if she can intercept both channels completely, detect in the same manner as does the legitimate receiver (i.e., Eve should also have access to a local oscillator phase stable with respect to that of sender and receiver) and retransmit the beams in the same way as the legitimate sender. Because of this possibility, our protocol as described is certainly not secure, in contrast to the protocols for discrete variables.$`^{\text{[39, 40, 41]}}`$ However, we suggest that simple extensions of our protocol might lead to significant enhancements in security.
If the goal were to achieve quantum key distribution, one idea is to make straightforward adaptations of the protocols introduced by Bennett and colleagues for the discrete case, as in Ref.. Here, we propose that the sending station (Alice) and receiving station (Bob) make random choices for the set of phases of the coherent message beam, as well as for the signal and idler beams. Recall that the variance $`V_{}(\mathrm{\Omega })`$ of Eq.1 is the minimum possible and applies only for the choice of quadrature-phase angles $`(\theta _A,\theta _B)`$ for the signal and idler beams that satisfy $`\theta _A+\theta _B=2p\pi `$ ($`p`$ = integer). For definiteness, assume the following two choices.
1. $`(\theta _A^0,\theta _B^0)`$, with $`\theta _A^0+\theta _B^0=0`$ and corresponding quadrature amplitudes $`(X_A,X_B)`$.
2. $`(\theta _A^{\pi /2}=\theta _A^0+\frac{\pi }{2},\theta _B^{\pi /2}=\theta _B^0+\frac{\pi }{2})`$ and corresponding quadrature amplitudes $`(Y_A,Y_B)`$.
In the first case, the minimum variance $`V_{}(\mathrm{\Omega })`$ results for the combination $`(X_AX_B)`$, while in the second case, the combination $`(Y_A+Y_B)`$ has minimum variance. This is because $`Y_BY_B`$ is equivalent to the shift $`\theta _B^{\pi /2}\theta _B^{\pi /2}+\pi `$, so that $`\theta _A^{\pi /2}+\theta _B^{\pi /2}+\pi =2\pi `$.
With these definitions, Alice at the sending station (randomly) makes one of two choices.
1. Phase $`0`$ – Set the quadrature-phase angles $`(\theta _A,\theta _B)`$ to $`(\theta _A^0,\theta _B^0)`$ and the phases $`\beta _{A,B}=\beta _{A,B}^0`$ for the coherent message beam $`|\alpha =|\alpha |\mathrm{exp}[i\beta ]`$ corresponding to the $`X`$ quadratures of $`(A,B)`$.
2. Phase $`\frac{\pi }{2}`$ – Set $`(\theta _A,\theta _B)`$ to $`(\theta _A^{\pi /2},\theta _B^{\pi /2})`$ and $`\beta _{A,B}=\beta _{A,B}^{\pi /2}=\beta _{A,B}^0\pm \frac{\pi }{2}`$ corresponding to the $`Y`$ quadratures.
The encoded message (which could consist of $`|\alpha |=[a_0,a_1]`$ for a binary transmission) is sent to Bob’s receiving station precisely as in Figure 1. Bob must then choose the appropriate phases $`(\varphi _A,\varphi _B)`$ for his local oscillators $`(LO_A,LO_B)`$ to detect quadrature amplitudes such that the spectral density $`\mathrm{\Phi }_{}(\mathrm{\Omega })`$ for the difference photocurrent $`i_{}i_Ai_B`$ is minimized and the signal maximized. In the case Phase $`0`$, denote the local oscillator settings as $`(\varphi _A^0,\varphi _B^0)`$, in correspondence to the detection of $`(X_A,X_B)`$ with minimum variance $`V_{}(\mathrm{\Omega })`$. On the other hand, for the case Phase $`\frac{\pi }{2}`$, the local oscillator phases $`(\varphi _A,\varphi _B)(\varphi _A^{\pi /2},\varphi _B^{\pi /2})=(\varphi _A^0+\frac{\pi }{2},\varphi _B^0+\frac{3\pi }{2})`$, in correspondence to the detection of $`(Y_A,Y_B)`$ with minimum variance. In both cases, the encoded message would be recovered with maximum signal-to-noise ratio. Note that precisely such a switching protocol was implemented in our prior experiment of Ref. with results as stated for the variances.
Of course, Bob does not know in advance which choice $`[0,\frac{\pi }{2}]`$ Alice will have made for any given transmission. Hence, he makes a random selection between the alternatives $`(\varphi _A^0,\varphi _B^0)`$ and $`(\varphi _A^{\pi /2},\varphi _B^{\pi /2})`$, recovering the message in some cases but not others. After a series of transmissions, Alice and Bob communicate publicly about their choice of bases, keeping measurement results only when their choices coincide.
Now, if an eavesdropper Eve attempts to intervene (either by a strategy of partial tapping or by one of complete interception and re-broadcast), she will necessarily increase the noise level and error rate at Bob’s receiving station. The random switching of the phases $`(\theta _A,\theta _B)`$ by Alice forces Eve to make a guess as to the correct quadratures $`(\delta _A,\delta _B)`$ to be detected. Having made a choice, information about the orthogonal quadrature is lost. Of course, rather than homodyne detection, she could choose to employ heterodyne detection to gain information about the full complex amplitude. However, relative to homodyne detection, heterodyne detection brings a well-known penalty of a $`3`$dB reduction in signal-to-noise ratio.$`^{\text{[51]}}`$
While it is beyond the scope of the current paper to make any claims about the quantitative limits to the information that Eve might access or about the absolute ability of Alice and Bob to detect her presence, we do suggest that these would be interesting questions to investigate. There are certainly intervention strategies beyond those that we have mentioned that a cunning Eve would want to consider, such as an adaptive strategy for adjusting the phases $`(\delta _A,\delta _B)`$ during the duration of the transmission of any given message.$`^{\text{[52]}}`$ Likewise, in any real-world setting, overcoming the deleterious effects of losses in propagation from Alice to Bob will be a overriding consideration. The question of preserving the entanglement of the initial EPR state in the face of such losses is a fascinating one for continuous quantum variables. Although initial attempts have been made to develop error correcting quantum codes for continuous variables,$`^{\text{[16, 17, 18]}}`$ no adequate solution seems to yet have been found. Finally, it would be of interest to analyze the case where only one of the two correlated beams is sent to Bob, with then Alice retaining the other.
###### Acknowledgements.
We gratefully acknowledge the comments of J.H. Shapiro who pointed out the connection of our experiment to Ref., of S. L. Braunstein and H. Mabuchi for critical discussions, and of one of the referees who brought to our attention the Vernon cipher. This work was supported by the Office of Naval Research, by the National Science Foundation, and by DARPA via the QUIC administered by the Army Research Office. |
warning/0003/hep-ph0003253.html | ar5iv | text | # Contact interaction probes at the Linear Collider with polarized electron and positron beams11footnote 1Partially supported by the Research Council of Norway, and by MURST (Italian Ministry of University, Scientific Research and Technology).
## 1 Introduction
The possibility of constructing high energy polarized electron and positron beams is considered with great interest with regard to the physics programme at the Linear Collider (LC). Indeed, one of the most important advantages of initial beam polarization is that one can measure spin-dependent observables, which represent the most direct probes of the fermion helicity dependence of the electroweak interactions. Consequently, one would expect a substantial gain in the sensitivity to the features of possible non-standard interactions and, in particular, stringent constraints on the individual new coupling constants could be derived from the data analysis by looking for deviations of cross sections from the Standard Model (SM) predictions.
Here, we will consider the process of fermion pair production ($`fe`$, $`t`$)
$$e^++e^{}f+\overline{f}$$
(1)
at a future Linear Collider with longitudinally polarized electron and positron beams, and discuss the sensitivity of the measurable helicity cross sections to the $`SU(3)\times SU(2)\times U(1)`$ symmetric $`eeff`$ contact-interaction Lagrangian with helicity-conserving and flavor-diagonal fermion currents :
$$=\underset{\alpha \beta }{}\frac{g_{\mathrm{eff}}^2}{\mathrm{\Lambda }_{\alpha \beta }^2}\eta _{\alpha \beta }\left(\overline{e}_\alpha \gamma _\mu e_\alpha \right)\left(\overline{f}_\beta \gamma ^\mu f_\beta \right).$$
(2)
In Eq. (2), generation and color indices have been suppressed, $`\alpha ,\beta =\mathrm{L},\mathrm{R}`$ indicate left- or right-handed helicities, and the parameters $`\eta _{\alpha \beta }=\pm 1,0`$ specify the chiral structure of the individual interactions. Conventionally, one takes $`g_{\mathrm{eff}}^2=4\pi `$ as a reminder that the new interaction, originally proposed for compositeness, would become strong at $`\sqrt{s}\mathrm{\Lambda }_{\alpha \beta }`$. Actually, in a more general sense, $``$ should be considered as an effective Lagrangian which represents the leading, lowest dimensional, parameterization at the ‘low-energy’ $`E`$ at which we make measurements, of some non-standard interaction acting at a much larger energy scale $`\mathrm{\Lambda }E`$. For example, in addition to the remnant compositeness binding force, this is the case of a variety of interactions generated by the exchange of very heavy objects with masses much larger than the Mandelstam variables of the considered process (1), such as the exchanges of a $`Z^{}`$ with a few TeV mass and of a heavy leptoquark . In this effective framework, therefore, with the assumed conventional values of $`\eta `$’s and $`g_{\mathrm{eff}}^2`$, the scales $`\mathrm{\Lambda }_{\alpha \beta }`$ in Eq. (2) define a standard to compare the reach of different new-physics searches in the process (1).
Clearly, $``$ should manifest itself by deviations of observables from the SM theoretical predictions. The sensitivity of measurements to the new coupling constants, or, equivalently, the experimentally attainable reach in the free mass scales $`\mathrm{\Lambda }_{\alpha \beta }`$, can be assessed by the numerical comparison of such deviations to the expected experimental accuracies.
For a given flavor $`f`$, Eq. (2) defines eight individual, independent models corresponding to the combinations of the four chiralities $`\alpha ,\beta `$ with the $`\pm `$ signs of the $`\eta `$’s. However, in general, an observed contact interaction could be any linear combination of these models, and this leads to the complicated situation in which the aforementioned deviations of observables from the SM predictions simultaneously depend on all four-fermion effective couplings. A simplified, and commonly adopted, procedure is to assume a non-zero value for only one parameter at a time and constrain it by essentially a $`\chi ^2`$ fit analysis, keeping the remaining parameters set equal to zero. In this way, tests of the individual models are obtained.
On the other hand, a general, model-independent, analysis must simultaneously include all terms of Eq. (2) as free parameters and, at the same time, must allow to disentangle their contributions to the basic observables so as to avoid potential cancellations between different contributions. Such cancellations can make the constraints considerably weaker or even spoil them. For this purpose, the longitudinal polarization of initial beams offers the possibility of experimentally separating from the data the individual helicity cross sections of process (1), each one being directly related to a single $`eeff`$ contact term and, therefore, depending on the minimal set of free independent parameters. The approach we adopt here uses as basic observables two particular, polarized, integrated cross sections that allow to reconstruct the four helicity amplitudes via linear combinations of measurements at different beam polarizations.<sup>1</sup><sup>1</sup>1Integrated observables should be of advantage in the case of limited experimental statistics. Moreover, in the definition of such integrated observables, optimal kinematical regions can be chosen to maximize the sensitivity to the individual four-fermion contact interactions.
This kind of analysis, and the determination of the corresponding reach on $`\mathrm{\Lambda }_{\alpha \beta }`$, was applied in Ref. for the LC with $`\sqrt{s}=0.5\mathrm{TeV}`$ and only the electron beam longitudinally polarized, making standard assumptions on the luminosity and on the expected systematic uncertainties on the cross section of process (1) for the different flavors. Indeed, longitudinal polarization of one beam is by itself already sufficient to disentangle the helicity cross sections from the data, if at least two values of the polarization are available, e.g., $`\pm |P_e|`$. In what follows, we extend the analysis of Ref. and discuss the case where also positron beam longitudinal polarization is available at the LC with the same c.m. energy. Specifically, after giving the main definitions and briefly reviewing the procedure and findings for the sensitivity on $`\mathrm{\Lambda }_{\alpha \beta }`$ obtained in , we start by considering the effect of the uncertainty on the electron beam polarization that was disregarded there. We then consider the case of both electron and positron longitudinal polarizations, including in the analysis also the uncertainty on these polarizations.
## 2 Separation of the helicity cross sections
In Eq. (1) we limit ourselves to the cases $`fe,t`$ and make the approximation of negligible fermion mass with respect to the c.m. energy $`\sqrt{s}`$. Then, the amplitude for $`e^+e^{}f\overline{f}`$ is determined by the Born, $`s`$-channel, $`\gamma `$ and $`Z`$ exchanges plus the contact-interaction term of Eq. (2). With $`P_e`$ and $`P_{\overline{e}}`$ the longitudinal polarizations of the beams, and $`\theta `$ the angle between the incoming electron and the outgoing fermion in the c.m. frame, the differential cross section reads :
$$\frac{\text{d}\sigma }{\text{d}\mathrm{cos}\theta }=\frac{3}{8}\left[(1+\mathrm{cos}\theta )^2\sigma _++(1\mathrm{cos}\theta )^2\sigma _{}\right].$$
(3)
In terms of helicity cross sections $`\sigma _{\alpha \beta }`$ (with $`\alpha ,\beta =\mathrm{L},\mathrm{R}`$):
$`\sigma _+`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[(1P_e)(1+P_{\overline{e}})\sigma _{\mathrm{LL}}+(1+P_e)(1P_{\overline{e}})\sigma _{\mathrm{RR}}\right]`$ (4)
$`=`$ $`{\displaystyle \frac{D}{4}}\left[(1P_{\mathrm{eff}})\sigma _{\mathrm{LL}}+(1+P_{\mathrm{eff}})\sigma _{\mathrm{RR}}\right],`$
$`\sigma _{}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[(1P_e)(1+P_{\overline{e}})\sigma _{\mathrm{LR}}+(1+P_e)(1P_{\overline{e}})\sigma _{\mathrm{RL}}\right]`$ (5)
$`=`$ $`{\displaystyle \frac{D}{4}}\left[(1P_{\mathrm{eff}})\sigma _{\mathrm{LR}}+(1+P_{\mathrm{eff}})\sigma _{\mathrm{RL}}\right],`$
where
$$P_{\mathrm{eff}}=\frac{P_eP_{\overline{e}}}{1P_eP_{\overline{e}}}$$
(6)
is the effective polarization , $`|P_{\mathrm{eff}}|1`$, and $`D=1P_eP_{\overline{e}}`$. Obviously, for unpolarized positrons $`P_{\mathrm{eff}}P_e`$ and $`D1`$. It should be noted that with $`P_{\overline{e}}0`$, $`|P_{\mathrm{eff}}|`$ can be larger than $`|P_e|`$. Moreover, in Eqs. (4) and (5):
$$\sigma _{\alpha \beta }=N_C\sigma _{\mathrm{pt}}|A_{\alpha \beta }|^2,$$
(7)
where $`N_C3(1+\alpha _s/\pi )`$ for quarks and $`N_C=1`$ for leptons, respectively, and $`\sigma _{\mathrm{pt}}\sigma (e^+e^{}\gamma ^{}l^+l^{})=(4\pi \alpha ^2)/(3s)`$. The helicity amplitudes $`A_{\alpha \beta }`$ can be written as
$$A_{\alpha \beta }=Q_eQ_f+g_\alpha ^eg_\beta ^f\chi _Z+\frac{s\eta _{\alpha \beta }}{\alpha \mathrm{\Lambda }_{\alpha \beta }^2},$$
(8)
where $`\chi _Z=s/(sM_Z^2+iM_Z\mathrm{\Gamma }_Z)`$ is the gauge boson propagator, $`g_\mathrm{L}^f=(I_{3L}^fQ_fs_W^2)/s_Wc_W`$ and $`g_\mathrm{R}^f=Q_fs_W^2/s_Wc_W`$ are the SM left- and right-handed fermion couplings of the $`Z`$ with $`s_W^2=1c_W^2\mathrm{sin}^2\theta _W`$ and $`Q_f`$ the fermion electric charge.
Our analysis focuses on the helicity cross sections that, as the above relations clearly show, directly relate to the individual contact interactions in Eq. (2) with definite chiralities and, accordingly, lead to a model-independent analysis where all terms in this equation are taken into account as completely free parameters with no danger of accidental compensations. To disentangle the various contributions in Eqs. (4) and (5), one simply has to make measurements at two different values of the polarizations (a minimum of four measurements is needed). For example, two convenient sets of values for the polarizations, that we will use in the sequel, would be $`P_e=\pm P_1`$ and $`P_{\overline{e}}=P_2`$ ($`P_{1,2}>0`$) or, alternatively, $`P_{\mathrm{eff}}=\pm P`$ and $`D`$ fixed. The corresponding solutions of Eqs. (4) and (5) read:
$`\sigma _{\mathrm{LL}}`$ $`=`$ $`{\displaystyle \frac{1}{D}}\left[{\displaystyle \frac{1P}{P}}\sigma _+(P)+{\displaystyle \frac{1+P}{P}}\sigma _+(P)\right],`$ (9)
$`\sigma _{\mathrm{RR}}`$ $`=`$ $`{\displaystyle \frac{1}{D}}\left[{\displaystyle \frac{1+P}{P}}\sigma _+(P){\displaystyle \frac{1P}{P}}\sigma _+(P)\right],`$ (10)
with $`\sigma _{\mathrm{LR}}`$ and $`\sigma _{\mathrm{RL}}`$ obtained from $`\sigma _{\mathrm{LL}}`$ and $`\sigma _{\mathrm{RR}}`$, respectively, replacing $`\sigma _+`$ by $`\sigma _{}`$.
Actually, for the purpose of optimizing the resulting bounds on $`\mathrm{\Lambda }_{\alpha \beta }`$, one can more generally define the polarized cross sections integrated over the a priori arbitrary kinematical ranges ($`1,z^{}`$) and ($`z^{},1`$) :
$`\sigma _1(z^{},P,D)`$ $``$ $`{\displaystyle _z^{}^1}{\displaystyle \frac{\text{d}\sigma }{\text{d}\mathrm{cos}\theta }}\text{d}\mathrm{cos}\theta ={\displaystyle \frac{1}{8}}\left\{\left[8(1+z^{})^3\right]\sigma _++(1z^{})^3\sigma _{}\right\},`$ (11)
$`\sigma _2(z^{},P,D)`$ $``$ $`{\displaystyle _1^z^{}}{\displaystyle \frac{\text{d}\sigma }{\text{d}\mathrm{cos}\theta }}\text{d}\mathrm{cos}\theta ={\displaystyle \frac{1}{8}}\left\{(1+z^{})^3\sigma _++\left[8(1z^{})^3\right]\sigma _{}\right\}.`$ (12)
For simplicity of notations, the polarization dependence of $`\sigma _\pm `$ on the right-hand sides of Eqs. (11) and (12) has been suppressed. As abbreviations, we introduce
$$a(z^{})=\frac{8(1z^{})^3}{6(1z_{}^{}{}_{}{}^{2})},b(z^{})=\frac{(1z^{})^3}{6(1z_{}^{}{}_{}{}^{2})}.$$
(13)
By solving Eqs. (11) and (12) one obtains $`\sigma _+`$ and $`\sigma _{}`$ from the measurement of $`\sigma _1`$ and $`\sigma _2`$:
$`\sigma _+`$ $`=`$ $`\left[a(z^{})\sigma _1(z^{},P,D)+b(z^{})\sigma _2(z^{},P,D)\right],`$ (14)
$`\sigma _{}`$ $`=`$ $`\left[b(z^{})\sigma _1(z^{},P,D)+a(z^{})\sigma _2(z^{},P,D)\right].`$ (15)
Thus, according to this procedure, $`\sigma _{1,2}(z^{},P,D)`$ play the role of a basic set of integrated polarized observables to be measured. As a second step, the corresponding cross sections $`\sigma _\pm `$ are constructed using the relations (14) and (15) and the experimental values of the helicity cross sections $`\sigma _{\alpha \beta }`$ are finally determined from the linear system of equations (9)–(10). Moreover, the value of $`z^{}`$ is taken as an input parameter related to given experimental conditions, that can be tuned in order to get maximal sensitivity of the helicity cross sections $`\sigma _{\alpha \beta }`$ to the mass scales $`\mathrm{\Lambda }_{\alpha \beta }`$ we want to constrain.
For comparison, we recall the conventional observables, the total cross section $`\sigma `$ and the various asymmetries. These are generally given, according to Eqs. (3)–(6), by
$$\sigma =\sigma _++\sigma _{}=\frac{D}{4}\left[(1P_{\mathrm{eff}})(\sigma _{\mathrm{LL}}+\sigma _{\mathrm{LR}})+(1+P_{\mathrm{eff}})(\sigma _{\mathrm{RR}}+\sigma _{\mathrm{RL}})\right];$$
(16)
and
$`\sigma A_{\mathrm{FB}}`$ $``$ $`\sigma _\mathrm{F}\sigma _\mathrm{B}={\displaystyle \frac{3}{4}}\left(\sigma _+\sigma _{}\right)`$ (17)
$`=`$ $`{\displaystyle \frac{3}{16}}D\left[(1P_{\mathrm{eff}})(\sigma _{\mathrm{LL}}\sigma _{\mathrm{LR}})+(1+P_{\mathrm{eff}})(\sigma _{\mathrm{RR}}\sigma _{\mathrm{RL}})\right];`$
with
$$\sigma _\mathrm{F}=\sigma _1(z^{}=0)=_0^1(\text{d}\sigma /\text{d}\mathrm{cos}\theta )\text{d}\mathrm{cos}\theta ;\sigma _\mathrm{B}=\sigma _2(z^{}=0)=_1^0(\text{d}\sigma /\text{d}\mathrm{cos}\theta )\text{d}\mathrm{cos}\theta ,$$
(18)
and $`P_{\mathrm{eff}}0`$, $`D1`$ for unpolarized beams. For the case of polarized beams, one has also the left-right asymmetry
$$A_{\mathrm{LR}}=\frac{\sigma _\mathrm{L}\sigma _\mathrm{R}}{\sigma _\mathrm{L}+\sigma _\mathrm{R}}=\frac{(\sigma _{\mathrm{LL}}+\sigma _{\mathrm{LR}})(\sigma _{\mathrm{RL}}+\sigma _{\mathrm{RR}})}{\sigma _{\mathrm{LL}}+\sigma _{\mathrm{LR}}+\sigma _{\mathrm{RL}}+\sigma _{\mathrm{RR}}},$$
(19)
and the combined left-right forward-backward asymmetry
$$A_{\mathrm{LR},\mathrm{FB}}=\frac{(\sigma _\mathrm{L}^\mathrm{F}\sigma _\mathrm{R}^\mathrm{F})(\sigma _\mathrm{L}^\mathrm{B}\sigma _\mathrm{R}^\mathrm{B})}{(\sigma _\mathrm{L}^\mathrm{F}+\sigma _\mathrm{R}^\mathrm{F})+(\sigma _\mathrm{L}^\mathrm{B}+\sigma _\mathrm{R}^\mathrm{B})}=\frac{3}{4}\frac{\sigma _{\mathrm{LL}}\sigma _{\mathrm{RR}}+\sigma _{\mathrm{RL}}\sigma _{\mathrm{LR}}}{\sigma _{\mathrm{LL}}+\sigma _{\mathrm{RR}}+\sigma _{\mathrm{RL}}+\sigma _{\mathrm{LR}}},$$
(20)
where $`\sigma _\mathrm{L}`$ and $`\sigma _\mathrm{R}`$ denote the cross sections with left-handed and right-handed electrons and unpolarized positrons.
In the numerical analysis, radiative corrections including initial- and final-state radiation are taken into account by means of the program ZFITTER , which has to be used along with ZEFIT, adapted to the present discussion, with $`m_{\mathrm{top}}=175`$ GeV and $`m_H=100`$ GeV. One-loop SM electroweak corrections are accounted for by improved Born amplitudes , such that the form of the previous formulae remains the same. Concerning initial-state radiation, a cut on the energy of the emitted photon $`\mathrm{\Delta }=E_\gamma /E_{\mathrm{beam}}=0.9`$ is applied for $`\sqrt{s}=0.5\mathrm{TeV}`$ in order to avoid the radiative return to the $`Z`$ peak, and increase the signal originating from the contact interaction .
## 3 Sensitivity of observables and their optimization
Given the current bounds on $`\mathrm{\Lambda }_{\alpha \beta }`$, of the order of several TeV , at the LC c.m. energy $`\sqrt{s}=0.5`$ TeV the characteristic suppression factor $`s/\mathrm{\Lambda }^2`$ in Eq. (8) is such that we can only look at indirect manifestations of the contact interaction (2) as deviations from the SM predictions. In this case, we can assess the sensitivity of process (1) to the couplings in (2), that determines the corresponding reach on $`\mathrm{\Lambda }_{\alpha \beta }`$, on the basis of the foreseen experimental accuracy on the helicity cross sections $`\sigma _{\alpha \beta }`$. As stressed previously, the knowledge of the latter allows a model-independent analysis, where all the contact interactions are disentangled and therefore can be taken into account as free parameters simultaneously.
Specifically, we define the ‘significance’ of each helicity cross section by the ratio
$$𝒮(\sigma _{\alpha \beta })=\frac{|\mathrm{\Delta }\sigma _{\alpha \beta }|}{\delta \sigma _{\alpha \beta }},$$
(21)
where $`\mathrm{\Delta }\sigma _{\alpha \beta }`$ are the deviations from the SM prediction due to (2), dominated for $`\sqrt{s}\mathrm{\Lambda }_{\alpha \beta }`$ by the linear interference term
$$\mathrm{\Delta }\sigma _{\alpha \beta }\sigma _{\alpha \beta }\sigma _{\alpha \beta }^{\mathrm{SM}}2N_C\sigma _{\mathrm{pt}}\left(Q_eQ_f+g_\alpha ^eg_\beta ^f\chi _Z\right)\frac{s\eta _{\alpha \beta }}{\alpha \mathrm{\Lambda }_{\alpha \beta }^2},$$
(22)
and $`\delta \sigma _{\alpha \beta }`$ denotes the expected experimental uncertainty on $`\sigma _{\alpha \beta }`$, combining statistical and systematic uncertainties.
In the procedure of determining helicity amplitudes via the integrated polarized cross sections $`\sigma _{1,2}`$ outlined in the previous section (see Eqs. (9), (10), (14) and (15)), adding all uncertainties in quadrature and neglecting for the moment the systematic uncertainty on the electron and positron polarizations, one can write:
$`\left(\delta \sigma _{\mathrm{LL}}\right)^2`$ $`=`$ $`a^2(z^{})\left[\left({\displaystyle \frac{1P}{PD}}\right)^2(\delta \sigma _1(z^{},P))^2+\left({\displaystyle \frac{1+P}{PD}}\right)^2(\delta \sigma _1(z^{},P))^2\right]`$ (23)
$`+`$ $`b^2(z^{})\left[\left({\displaystyle \frac{1P}{PD}}\right)^2(\delta \sigma _2(z^{},P))^2+\left({\displaystyle \frac{1+P}{PD}}\right)^2(\delta \sigma _2(z^{},P))^2\right],`$
$`\left(\delta \sigma _{\mathrm{LR}}\right)^2`$ $`=`$ $`b^2(z^{})\left[\left({\displaystyle \frac{1P}{PD}}\right)^2(\delta \sigma _1(z^{},P))^2+\left({\displaystyle \frac{1+P}{PD}}\right)^2(\delta \sigma _1(z^{},P))^2\right]`$ (24)
$`+`$ $`a^2(z^{})\left[\left({\displaystyle \frac{1P}{PD}}\right)^2(\delta \sigma _2(z^{},P))^2+\left({\displaystyle \frac{1+P}{PD}}\right)^2(\delta \sigma _2(z^{},P))^2\right],`$
where $`a`$ and $`b`$ are given by Eq. (13). For simplicity of notations, the dependence of $`\delta \sigma _{1,2}`$ on $`D`$ has not been explicitly indicated. One can derive explicit expressions for $`\delta \sigma _{\mathrm{RR}}`$ and $`\delta \sigma _{\mathrm{RL}}`$ from $`\delta \sigma _{\mathrm{LL}}`$ and $`\delta \sigma _{\mathrm{LR}}`$, respectively, by the replacement in the above equations of $`\pm PP`$ in $`\delta \sigma _i(z^{},\pm P))`$ but not in the corresponding prefactors.
Combining in quadrature statistical and systematic uncertainties on $`\sigma _{1,2}`$, one finds:
$$(\delta \sigma _i)^2(\delta \sigma _i^{\mathrm{SM}})^2=\frac{\sigma _i^{\mathrm{SM}}}{ϵ_{\mathrm{int}}}+\left(\delta ^{\mathrm{sys}}\sigma _i^{\mathrm{SM}}\right)^2,i=1,2.$$
(25)
For our numerical analysis we shall assume the commonly used reference values of the identification efficiencies, $`ϵ`$, and the systematic uncertainties, $`\delta ^{\mathrm{sys}}`$, for the various fermionic channels : $`ϵ=95\%`$ and $`\delta ^{\mathrm{sys}}=0.5\%`$ for $`l^+l^{}`$; $`ϵ=60\%`$ and $`\delta ^{\mathrm{sys}}=1\%`$ for $`b\overline{b}`$; $`ϵ=35\%`$ and $`\delta ^{\mathrm{sys}}=1.5\%`$ for $`c\overline{c}`$. Notice that, as a simplification, we take the same $`\delta ^{\mathrm{sys}}`$ for both $`i=1`$ and 2, and independent of $`z^{}`$ in the relevant angular range. Concerning the statistical uncertainty, we consider the LC with $`\sqrt{s}=0.5`$ TeV, $`_{\mathrm{int}}=50\text{fb}^1`$ and $`_{\mathrm{int}}=500\text{fb}^1`$ (half for each polarization orientation), and a fiducial experimental angular range $`|\mathrm{cos}\theta |0.99`$.
Finally, as regards optimization of the bounds on contact-interaction couplings, which corresponds to the maximum value of the ‘significance’ defined in Eq. (21), one may notice from the equations above that the uncertainties $`\delta \sigma _{\alpha \beta }`$ depend on the, a priori free, kinematical parameter $`z^{}`$ in the definition of the polarized cross sections $`\sigma _i`$. Conversely, by definition, the deviations from the SM $`\mathrm{\Delta }\sigma _{\alpha \beta }`$ in Eq. (22) are $`z^{}`$-independent. Therefore, optimization can be achieved by choosing $`z^{}=z_{\mathrm{opt}}^{}`$ where $`\delta \sigma _{\alpha \beta }`$ becomes minimum, so that the corresponding sensitivity has a maximum and determines the highest bound on the corresponding mass scale $`\mathrm{\Lambda }_{\alpha \beta }`$. The $`z^{}`$ dependence of the statistical uncertainties $`\delta \sigma _{\alpha \beta }^{\mathrm{stat}}`$ in the right-hand side of (25) can be approximated by that corresponding to the known SM cross sections for the process (1) and the value of $`_{\mathrm{int}}`$. In the case of low luminosity where the statistical uncertainty dominates, this SM-determined $`z^{}`$ behaviour can be used for a simple, first determination of $`z_{\mathrm{opt}}^{}`$ for the various helicity amplitudes . In the general case where statistical and systematic uncertainties are comparable, the optimal $`z^{}`$ must be determined by a more complex numerical analysis taking into account the relevant experimental details.
## 4 Polarization uncertainty and two polarized beams
In order to assess the effects on the $`\delta \sigma _{\alpha \beta }`$ due to the systematic uncertainties $`\delta P_e`$ and $`\delta P_{\overline{e}}`$ on the $`e^{}`$ and $`e^+`$ polarizations respectively, we must supplement by appropriate terms Eqs. (23) and (24) and the similar ones for the remaining helicity amplitudes. From the formulae in Sec. 2, one can see that finite values of $`\delta P_e`$ and $`\delta P_{\overline{e}}`$ will influence the extraction of the helicity cross sections $`\sigma _{\alpha \beta }`$ through the prefactors of Eqs. (9), (10), (14) and (15), as well as through the dependence of $`\sigma _{1,2}`$ on $`P`$ and $`D`$. Clearly, a complete assessment of the latter effect would require detailed knowledge of the structure of the overall systematic uncertainty in terms of the different, individual sources, that is not available at present. For the sake of simplicity, we model the systematic uncertainty by assuming that such an effect can be considered as already included in the systematic uncertainties $`\delta \sigma _i^{\mathrm{sys}}`$ introduced in Eq. (25), regardless of the values of $`\delta P_e`$ and $`\delta P_{\overline{e}}`$ (and $`P_e`$ and $`P_{\overline{e}}`$) considered in our discussion. Then, we treat $`\sigma _{1,2}`$, $`P_e`$ and $`P_{\overline{e}}`$ in Eqs. (9) and (10) as if they were independent measurables, and in this spirit, we combine the additional contribution to the uncertainty, $`\delta \sigma _{\alpha \beta }^{\mathrm{pol}}`$, again in quadrature with the $`\delta \sigma _{\alpha \beta }`$ determined from the expressions (23) and (24). Thus:
$$(\delta \sigma _{\alpha \beta })^2(\delta \sigma _{\alpha \beta })^2+\left(\delta \sigma _{\alpha \beta }^{\mathrm{pol}}\right)^2.$$
(26)
Under the above assumptions, we obtain
$`\left(\delta \sigma _{\mathrm{LL}}^{\mathrm{pol}}\right)^2`$ $`=`$ $`[f(z^{},P)(1+P_{\overline{e}}P^2)f(z^{},P)(1P_{\overline{e}}P^2)]^2\left({\displaystyle \frac{\delta P_e}{D^2P^2}}\right)^2`$
$`+`$ $`[f(z^{},P)(1P_eP^2)f(z^{},P)(1+P_eP^2)]^2\left({\displaystyle \frac{\delta P_{\overline{e}}}{D^2P^2}}\right)^2,`$
$`\left(\delta \sigma _{\mathrm{RR}}^{\mathrm{pol}}\right)^2`$ $`=`$ $`[f(z^{},P)(1P_{\overline{e}}P^2)f(z^{},P)(1+P_{\overline{e}}P^2)]^2\left({\displaystyle \frac{\delta P_e}{D^2P^2}}\right)^2`$ (27)
$`+`$ $`[f(z^{},P)(1+P_eP^2)f(z^{},P)(1P_eP^2)]^2\left({\displaystyle \frac{\delta P_{\overline{e}}}{D^2P^2}}\right)^2,`$
with
$$f(z^{},P)=a(z^{})\sigma _1(z^{},P)+b(z^{})\sigma _2(z^{},P).$$
(28)
Furthermore, $`\delta \sigma _{\mathrm{LR}}^{\mathrm{pol}}`$ and $`\delta \sigma _{\mathrm{RL}}^{\mathrm{pol}}`$ are obtained from $`\delta \sigma _{\mathrm{LL}}^p`$ and $`\delta \sigma _{\mathrm{RR}}^p`$, respectively, by substituting $`a(z^{})b(z^{})`$. Numerically, for explicit assessments of the reach on $`\mathrm{\Lambda }_{\alpha \beta }`$, we shall work out the example of $`|P_e|=0.9`$ with $`\delta P_e/P_e=0.5\%`$ as currently attainable at the SLC , and $`|P_{\overline{e}}|=0.6`$ . This corresponds to the effective polarization $`P_{\mathrm{eff}}=P=0.974`$ and $`D=1.54`$. Clearly, introducing positron polarization may amount to a sort of “noise”, unless its magnitude is known with some precision. Since, at present, information on the achievable precision on the positron polarization is unknown, in our numerical analysis we shall vary $`\delta P_{\overline{e}}/P_{\overline{e}}`$ in a range up to a few tens of percent.
We start by considering, as a first example, electrons that are polarized, but unpolarized positrons, ($`|P_e|,|P_{\overline{e}}|`$) = ($`0.9,0.0`$), and then we discuss the case of both initial beams polarized with ($`|P_e|,|P_{\overline{e}}|`$) = ($`0.9,0.6`$). Also, as anticipated, we assume half the total integrated luminosity quoted above for each value of the effective polarization, $`P_{\mathrm{eff}}=\pm P`$. We focus on the impact of finite polarization uncertainties on the sensitivity of the helicity cross sections $`\sigma _{\alpha \beta }`$ to the contact interaction (2) and the corresponding reach on the mass scales $`\mathrm{\Lambda }_{\alpha \beta }`$ that, as discussed in the previous section, is determined by the uncertainties $`\delta \sigma _{\alpha \beta }`$ via Eqs. (21), (22) and (26).
In the starting example, with polarized electrons and unpolarized positrons, we compare the relative deviations $`\delta \sigma _{\alpha \beta }/\sigma _{\alpha \beta }\delta \sigma _{\alpha \beta }^{\mathrm{SM}}/\sigma _{\alpha \beta }^{\mathrm{SM}}`$ for finite $`\delta P_e`$ with the case of the same $`P_e`$, but $`\delta P_e=0`$, studied in . The ratio of the sensitivity (21) in the two cases, determining the effect of the electron polarization uncertainty introduced via Eq. (26), is shown in Fig. 1, for $`_{\mathrm{int}}=50\mathrm{fb}^1`$. This figure is obtained using the optimization procedure, and the determination of the relevant $`z_{\mathrm{opt}}^{}`$, outlined in the previous section. The sensitivity, via its square root, determines the reach in $`\mathrm{\Lambda }_{\alpha \beta }`$. For the $`\mu ^+\mu ^{}`$ final state, and LL and RR helicity configurations, the effect of $`\delta P_e`$ determining $`\delta \sigma _{\alpha \beta }^{\mathrm{pol}}`$ in (26) is found to change $`\delta \sigma _{\mathrm{LL}}`$ and $`\delta \sigma _{\mathrm{RR}}`$ as given by (23)–(25) and the stated input values by a really modest amount, of the order of a fraction of a %, unless $`\delta P_e/P_e`$ exceeds 3–4%, whereas for $`\delta \sigma _{\mathrm{LR}}`$ and $`\delta \sigma _{\mathrm{RL}}`$ there is no change at all. The reason for this can be found in Eqs. (4) and (28). Indeed, within the set of assumptions leading to those equations, one has numerically:
$`\left(\delta \sigma _{\mathrm{LL},\mathrm{RR}}^{\mathrm{pol}}\right)^2`$ $``$ $`[(\sigma _{\mathrm{LL}}\sigma _{\mathrm{RR}})P_{\overline{e}}P(\sigma _{\mathrm{LL}}+\sigma _{\mathrm{RR}})]^2(\delta P_e)^2`$ (29)
$`+`$ $`[(\sigma _{\mathrm{LL}}\sigma _{\mathrm{RR}})\pm P_eP(\sigma _{\mathrm{LL}}+\sigma _{\mathrm{RR}})]^2(\delta P_{\overline{e}})^2,`$
independent of $`z^{}`$, and similar expressions for $`\delta \sigma _{\mathrm{LR},\mathrm{RL}}^{\mathrm{pol}}`$ with the substitutions $`\text{LL,RR}\text{LR,RL}`$. Thus, for $`P_{\overline{e}}=\delta P_{\overline{e}}=0`$, $`\delta \sigma _{\mathrm{LL}}^{\mathrm{pol}}\delta \sigma _{\mathrm{RR}}^{\mathrm{pol}}\sigma _{\mathrm{LL}}^{\mathrm{SM}}\sigma _{\mathrm{RR}}^{\mathrm{SM}}`$, which for final-state muons vanishes in the limit of $`\mathrm{sin}^2\theta _\mathrm{W}0.25`$, whereas $`\delta \sigma _{\mathrm{LR}}^{\mathrm{pol}}\delta \sigma _{\mathrm{RL}}^{\mathrm{pol}}\sigma _{\mathrm{LR}}^{\mathrm{SM}}\sigma _{\mathrm{RL}}^{\mathrm{SM}}=0`$. It should be stressed that this lack of sensitivity to $`\delta P_e`$ depends on having no positron polarization, $`P_{\overline{e}}=0`$. For quarks, the corresponding differences of helicity cross sections do not vanish, and the effect of $`\delta P_e`$ is to yield a non-zero $`\delta \sigma _{\alpha \beta }^{\mathrm{pol}}`$. The contribution of $`\delta P_e`$ to the helicity cross section uncertainty, $`\delta \sigma _{\alpha \beta }^{\mathrm{pol}}`$, is still quite small with respect to the total uncertainty, as long as $`\delta P_e/P_e`$ remains less than 2–3%, except for the LL and RR cases of $`b\overline{b}`$ final states.
For higher luminosity, the curves become steeper, i.e., the sensitivity deteriorates faster with loss of polarization accuracy, see Fig. 2.
Turning to the case of both positron and electron longitudinal polarization, and referring to Eqs. (4) and (5), in the chosen helicity configuration where $`P_eP_{\overline{e}}<0`$, one has $`D>1`$ and $`|P_{\mathrm{eff}}|>\mathrm{max}(|P_e|,|P_{\overline{e}}|)`$, and in principle one could expect on statistical grounds an increase of the sensitivity due to the polarization of positrons, provided the luminosity remains the same. However, this improvement from positron polarization is obtained up to a maximum value of $`\delta P_{\overline{e}}/P_{\overline{e}}`$, above which there is no benefit, but, actually, a worsening of the sensitivity.
Indeed, it is instructive to compare the sensitivity of the helicity cross sections to four-fermion contact interactions for both beams polarized with that obtained with just one beam polarized. This comparison is expressed in terms of ratios of sensitivities as a function of the positron polarization uncertainty, $`\delta P_{\overline{e}}/P_{\overline{e}}`$, for $`_{\mathrm{int}}=50\mathrm{fb}^1`$ in Fig. 3 for lepton and quark final states. It is seen that if $`|\delta P_{\overline{e}}/P_{\overline{e}}||\delta P_e/P_e|=0.5\%`$, the advantage of positron polarization manifests itself in an increase in sensitivity by 10–40% depending on the helicity configuration and the final state. However, this ratio drops with increasing $`\delta P_{\overline{e}}/P_{\overline{e}}`$, and at those positron polarization uncertainties where it becomes less than unity, the advantage of positron polarization disappears. This useful region of the precision $`\delta P_{\overline{e}}/P_{\overline{e}}`$ ranges from 2% up to beyond 20% depending on the reaction and helicity combination.
This dependence on $`\delta P_{\overline{e}}`$ can be qualitatively understood from Eq. (29). In the case of muons (as opposed to quarks), the first term (proportional to $`(\delta P_e)^2`$) is relatively small (since we consider $`P_{\overline{e}}`$ considerably less than $`P_e`$), and the second term involving $`(\delta P_{\overline{e}})^2`$ becomes important already at small values of $`\delta P_{\overline{e}}`$. This explains why the curves (see Fig. 3) are rather steep. Other properties of Fig. 3 are also seen to follow from Eq. (29): (i) for $`e^+e^{}\mu ^+\mu ^{}`$, the dependence on $`\delta P_{\overline{e}}`$ is the same for the LR and RL cross sections, as well as for the LL and RR ones; (ii) for $`e^+e^{}b\overline{b}`$, the dependence on $`\delta P_{\overline{e}}`$ is relatively weak for the LR and RL cross sections since these cross sections are small; (iii) for $`e^+e^{}c\overline{c}`$, the dependence on $`\delta P_{\overline{e}}`$ is much weaker for the RL than for the LR cross section since $`\sigma _{\mathrm{LR}}`$ is bigger than $`\sigma _{\mathrm{RL}}`$, leading to a cancellation in one case and not in the other.
At higher luminosity, all the curves become more steep, since the uncertainty due to the polarization becomes more important w.r.t. the statistical uncertainty. For example, at $`_{\mathrm{int}}=500\mathrm{fb}^1`$, as Fig. 4 shows, for muon final states the positron polarization (at a value $`P_{\overline{e}}=0.6`$) stops being useful for the RR and LL cross sections at $`\delta P_{\overline{e}}/P_{\overline{e}}=0.8\%`$ and 0.6%, respectively.
In the next section we are going to conclude our numerical discussion by explicitly deriving the reach on the mass scales $`\mathrm{\Lambda }_{\alpha \beta }`$ obtainable in the case where the uncertainty on the electron and positron longitudinal polarizations are, respectively, 0.5% and 1%, for the two values $`_{\mathrm{int}}=50\text{fb}^1`$ and $`_{\mathrm{int}}=500\text{fb}^1`$.
## 5 Bounds on $`\mathrm{\Lambda }_{\alpha \beta }`$
As a preliminary step in the derivation of the constraints on $`\mathrm{\Lambda }_{\alpha \beta }`$, we show in Fig. 5 the relative uncertainties $`\delta \sigma _{\alpha \beta }/\sigma _{\alpha \beta }`$ as functions of $`z^{}`$, for the lower option for the luminosity.
The optimal values of $`z^{}`$ where the sensitivity is maximum can be easily read off from these figures, and in Table 1 we report such $`z_{\mathrm{opt}}^{}`$ for the two different values of the luminosity.
Numerical constraints on the four-fermion contact interactions of Eq. (2) are obtained from a $`\chi ^2`$ analysis of data on each helicity cross section, with (see Eq. (21)):
$$\chi ^2=\left(\frac{\mathrm{\Delta }\sigma _{\alpha \beta }}{\delta \sigma _{\alpha \beta }}\right)^2.$$
(30)
Bounds on the allowed values of the contact interaction parameters from the non-observation of the corresponding deviations within the expected uncertainty $`\delta \sigma _{\alpha \beta }`$ are derived by imposing $`\chi ^2<\chi _{\mathrm{CL}}^2`$, where the actual value of $`\chi _{\mathrm{CL}}^2`$ specifies the desired ‘confidence’ level. As Eq. (22) shows, the deviations $`\mathrm{\Delta }\sigma _{\alpha \beta }`$ depend on a single ‘effective’ contact-interaction free parameter, and therefore in such a $`\chi ^2`$ analysis we take $`\chi _{\mathrm{CL}}^2=3.84`$ for 95% C.L. as consistent with a one-parameter fit.
The results for the bounds on $`\mathrm{\Lambda }_{\alpha \beta }`$ are reported in Table 1. The table shows that the helicity cross sections $`\sigma _{\alpha \beta }`$ are quite sensitive to contact interactions, with discovery limits that, at the highest considered luminosity 500 $`\text{fb}^1`$, can range from 75 up to 150 times the c.m. energy, depending on the considered final fermion state. Indeed, the best sensitivity is achieved for the $`\mu ^+\mu ^{}`$ and $`b\overline{b}`$ final states, while the worst one corresponds to the $`c\overline{c}`$ channel. A direct comparison with the sensitivity achieved using ‘conventional’ observables, Eqs. (16)–(20)<sup>2</sup><sup>2</sup>2See, for example, the results obtained in in the context of specific contact-interaction models., is quite difficult and might be unclear, because it depends on the assumed model of new physics involved and the kind of parameterization adopted for the uncertainty. In this regard, as repeatedly stressed, we point out that the separation of the helicity cross sections performed here (and the corresponding values in Table 1) has the qualitative advantage of providing, by definition, unambiguous and model-independent information on the non-standard parameters of Eq. (2).
For a sort of contact to the conventional observables (16)–(20), we have reported in Table 1 also the limits on $`\mathrm{\Lambda }_{\alpha \beta }`$ obtainable at $`z^{}=0`$ instead of $`z^{}=z_{\mathrm{opt}}^{}`$. The results show that, at $`z^{}=0`$, the sensitivity to $`\mathrm{\Lambda }_{\mathrm{LR}}`$ and $`\mathrm{\Lambda }_{\mathrm{RL}}`$ would be considerably smaller. As one can see from the table, the ‘optimal’ choice $`z^{}=z_{\mathrm{opt}}^{}`$ allows to substantially increase the bounds for the LR and RL cases, to the level of the LL and RR ones, for which the improvement is really modest. This relates to the $`z^{}`$ behavior of the relative uncertainties on $`\sigma _{\alpha \beta }`$, that, as is seen in Fig. 5, is flat in the latter case and varies more rapidly around $`z_{\mathrm{opt}}^{}`$ in the former one.
As discussed previously, and illustrated in Fig. 3, the benefit of positron polarization depends on it being known with some precision. We show in Figs. 6 and 7 the effect of the positron polarization uncertainty on the reach in $`\mathrm{\Lambda }_{\alpha \beta }`$, for the two luminosities considered.
We see from these figures that if the positron polarization is known with high precision, an amount $`P_{\overline{e}}=0.6`$ can increase the reach in $`\mathrm{\Lambda }`$ by typically 5–25%. The critical level of precision, by which the positron polarization should be known, in order to be beneficial for contact-interaction searches, depends very much on the channel considered, as well as the luminosity. At low luminosities, less polarization precision is required for the positron polarization to be useful.
While one polarized beam is a necessity in order to be able to extract the helicity cross sections, the benefit of both beams being polarized is less clear. For some combinations of final state and helicity channels, the increased reach in $`\mathrm{\Lambda }`$ can be considerable, although half luminosity (and correspondingly reduced number of events) has been assumed for the two configurations of electron and positron beam polarizations. However, due to the limiting effect of the polarization uncertainties on the sensitivity (21), such improvements do not seem dramatic. More luminosity might easily lead to the same gain, especially if the positron polarization is only known with a moderate accuracy.
Actually, for full completeness in this regard, the dependence on the actual value of the uncertainty $`\delta ^{\mathrm{sys}}`$ in (25) should be considered simultaneously with that from $`\delta \sigma _{\alpha \beta }^{\mathrm{pol}}`$, as suggested by the combination in Eq. (26). Clearly, we should expect reduction of the $`\mathrm{\Lambda }`$ reach for increasing $`\delta ^{\mathrm{sys}}`$. As an indication, by doubling the values of $`\delta ^{\mathrm{sys}}`$ with respect to those listed below Eq. (25), and adopted for the explicit numerical example presented here, at $`_{\mathrm{int}}=50\text{fb}^1`$ the typical effect amounts to a few percent for the LR and RL cases, but can be as large as 20% for the LL and RR combinations. This indicates that the latter helicity cross sections are much more sensitive to systematic uncertainties than the former ones.
Clearly, although these considerations are numerically drawn from a specific example using as inputs some particular, hypothetical, values of initial beam polarization and corresponding uncertainties, and from assumptions on the values and properties of the uncertainties $`\delta ^{\mathrm{sys}}`$ of Eq. (25), such conclusions should hold in general. For a definite, quantitative statement about the relative roles of statistical and systematic uncertainties (including $`\delta P_e`$ and $`\delta P_{\overline{e}}`$) in the determination of the accuracy on $`\sigma _{\alpha \beta }`$ in a realistic experimental situation, we must wait for more detailed information on the expected experimental errors. |
warning/0003/astro-ph0003290.html | ar5iv | text | # The evolution of the stellar hosts of radio galaxies
## 1 Introduction
At low redshifts, FRI and FRII radio sources with radio luminosities at 151 MHz of $`L_{R(151)}\stackrel{>}{_{}}10^{24}\mathrm{WHz}^1`$ are associated almost exclusively with giant elliptical host galaxies. If this continues to be the case out to high redshift, then radio galaxies can give us a unique insight into the formation and evolution of a single class of massive galaxy. This is particularly exciting in the light of submillimetre detections of $`z4`$ radio galaxies (e.g. Archibald et al. 2000) which may indicate that we can see these objects during their major bursts of star formation \[although see also Willott, Rawlings & Jarvis (2000)\]. Furthermore, the similarity of radio galaxy hosts, in contrast to the wide range in luminosity of radio-quiet quasar hosts (McLure et al. 1999; Ridgway et al. 2000), suggests that one of the conditions necessary for producing powerful radio jets is the presence of a massive spheroidal component, and therefore a supermassive ($`\stackrel{>}{_{}}10^9M_{})`$ black hole (Lacy, Ridgway & Trentham 2000).
Studies of high redshift radio galaxy hosts have traditionally concentrated on the $`Kz`$ Hubble Diagram. Work on the 3C and 1 Jy samples by Lilly (1989 and refs. therein) initially pointed to a passively-evolving stellar host formed at high redshift which evolved into the giant elliptical radio galaxy hosts seen today, but this was challenged when some high redshift radio galaxies were found to have significant emission line contributions to their $`K`$-band light (Eales & Rawlings 1993, 1996). Only by finding low AGN-luminosity radio galaxies at high redshift could the controversy be resolved. Eales et al. (1997) used the 6C sample, a factor of five fainter in radio flux than the 3C sample, to show that there did seem to be a radio luminosity dependence of host galaxy magnitude. Further work by Roche, Eales & Rawlings (1998) indicated that the hosts of 6C radio galaxies were not only significantly fainter than their 3C counterparts, but also had smaller scale sizes.
We have used the 7C-iii radio galaxy redshift survey of Lacy et al. (1999b) to select a complete sample of $`z>0.8`$ radio galaxies. The 7C redshift surveys are a factor of 4–5 lower still in luminosity at a given redshift than the 6C sample of Eales et al. (1997), and thus allow the study of high redshift radio galaxy hosts over a wide range in radio luminosity (Willott et al. 1999). The 7C-iii sources were imaged in the near-infrared on the 3-m NASA Infrared Telescope Facility (IRTF) and the 3-m Shane Telescope at Lick Observatory. These data have allowed us to further investigate the radio luminosity dependence of host properties at $`z1`$ and, because our 7C objects at $`z2`$ have similar radio luminosities to $`z1`$ 6C radio galaxies and $`z0.3`$ 3C radio galaxies, we can also investigate host galaxy evolution over a wide range in redshift.
We assume a cosmology with $`\mathrm{\Omega }_\mathrm{M}=1,\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`H_0=50\mathrm{kms}^1\mathrm{Mpc}^1`$ except where otherwise stated.
## 2 Observations
Most sources were observed on the IRTF with NSFCAM, a near infrared imaging camera employing a $`256\times 256`$ InSb array, on the nights of 1999 July 28 – 29 UT. The 0.3 arcsec/pixel scale was used for all observations. Details of the observations are given in Table 1. One of the $`J`$, $`H`$ or $`K^{^{}}`$ filters was chosen for each object so as to cover as far as possible the rest-frame $`R`$-band, thus minimizing the $`k`$-correction and the associated uncertainties. In practice this meant that objects with $`0.8<z<1.2`$ were observed in $`J`$-band, those with $`1.2<z<1.8`$ in $`H`$-band and those with $`z>1.8`$ in $`K`$-band (using a $`K^{^{}}`$ filter). Exceptions to this were two objects with uncertain redshifts, 7C 1756+6520 and 7C 1804+6313, which were both observed in $`K^{^{}}`$ for ease of estimating photometric redshifts. The sky was clear throughout the run, and most objects were observed in conditions of sub-arcsecond seeing. Two more objects with uncertain redshifts were observed using the Gemini instrument (McLean et al. 1993, 1994) on the Shane Telescope at Lick Observatory on the nights of 1999 June 1 and 4 UT. Gemini was used with a dichroic beamsplitter which enabled us to image in $`J`$ and $`K^{^{}}`$ simultaneously. The detector on the short wavelength arm was a $`256\times 256`$ HgCdTe array, and that on the long wavelength arm a $`256\times 256`$ InSb array. The image scale in both arms was 0.68 arcsec/pixel. 7C 1814+6704 was observed in both runs, but the better seeing of the IRTF data allowed more accurate photometry in the moderately crowded field so only the IRTF data is presented. In addition a $`K^{^{}}`$ image of 7C 1745+6624 was presented in Lacy et al. (1999b).
The data were reduced using the dimsum package in iraf, and flat-fielded using dome flats. The final images were magnified by a factor of two before combination to improve the sampling of the final image. The reduced images are shown in Fig. 1 and the photometric properties detailed in Table 2.
## 3 Analysis
### 3.1 Magnitudes and $`k`$-corrections
Aperture magnitudes were measured as listed in Table 2 and corrected to a standard 63.9 kpc metric aperture according to the prescription of Eales et al. (1997), which assumes a power-law curve of growth of the form $`I(<r)r^{0.35}`$ for $`z>0.6`$ objects. Our metric aperture corresponds to about 8 arcsec at $`z\stackrel{>}{_{}}1`$ in our assumed cosmology, and this angular size is not strongly dependent on the choice of $`\mathrm{\Omega }_\mathrm{M}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. For many of our objects we have been able to measure the magnitude in a large aperture, but for some we have had to restrict the aperture to $`3`$ arcsec to avoid contamination from neighbouring objects or excessive noise in the case of faint objects. The aperture correction applied in these cases is $`0.4`$ magnitudes. The disadvantage of this technique is that the aperture magnitude and in particular the correction are scale-size dependent. An alternative would have been to use total magnitudes (cf. Roche et al. 1998). However, we decided against this as for low signal-to-noise detections the curve of growth is ill-defined and can lead to large errors in the magnitudes (Stevens 1999). The mean half-light radius of our objects is 0.9 arcsec, so the typical amount of flux missed in an 8 arcsec aperture will be very small in any case. We thus believe that the use of these large aperture magnitudes will not seriously affect the results of this paper.
For objects with spectroscopic redshifts, $`k`$-corrections to rest-frame $`R`$-band were made using 1Gyr-burst models generated using the pegase code (Fioc & Rocca-Volmerange 1997), integrating over the appropriate filter profiles. A model galaxy with a total age of 2 Gyr was used for objects with $`2<z3`$, and 3 Gyr for objects with $`0.8<z2`$. The $`z>3`$ objects used a 1.4 Gyr-old model galaxy. For most objects the $`k`$-corrections are small ($`<0.2`$ magnitudes), the exceptions to this being the $`z>3`$ objects for which $`K^{^{}}`$ corresponds to rest-frame wavelengths significantly shorter than $`R`$-band (but still above 4000Å) and 7C 1756+6520 which was observed in $`K^{^{}}`$ to check its redshift photometrically. The biggest correction was $``$0.61 mag. for 7C 1814+6704 at a probable redshift of 4.05.
### 3.2 Photometric redshifts for objects in the 7C-iii sample
Several of the objects observed had uncertain redshifts \[grade $`\gamma `$ in Lacy et al. (1999b)\] or no redshift information at all. Of the objects without spectroscopic redshifts, 7C 1748+6703 has a $`K^{^{}}`$ magnitude and colours consistent with $`2.4<z<4`$ whereas 7C 1753+6311 and 7C 1804+6313 both have magnitudes and colours consistent with them having redshifts in the range $`1.2<z<1.8`$ (Willott et al. 2000). In addition Lacy et al. (1999b) estimated a photometric redshift of 1.7 for 7C 1743+6341. The identification of 7C 1753+6311 in Lacy et al. (1999b), which had a provisional redshift of 1.96, actually corresponds to a blue galaxy along the radio axis. Our Lick image showed that the true identification is a very red ($`RK^{^{}}\stackrel{>}{_{}}5`$) galaxy 4.7 arcsec away which is significantly closer to the likely radio central component. Nevertheless, the original identification may well be associated with the radio galaxy, and the magnitude and colours are consistent with $`z\stackrel{<}{_{}}2`$, so we have not revised our redshift estimate. As we aligned the slit with the radio axis, though, we might have expected to see Ly$`\alpha `$ in the spectrum if the redshift were indeed $`1.96`$. Thus we note that the redshift may well be lower. 7C 1756+6520 had a provisional redshift of 1.48, and its $`K^{^{}}`$ magnitude is just consistent with this, although somewhat fainter than the mean $`Kz`$ relation at this redshift. 7C 1814+6702 with a provisional redshift of 4.05 has a $`K^{^{}}`$ magnitude in both the IRTF and Lick data consistent with both $`z4`$ and $`1.2<z<1.8`$, although its diffuse and possibly aligned structure is typical of $`z>3`$ objects (van Breugel et al. 1998) so has been kept at 4.05 pending a deeper spectrum. 7C 1816+6710 was observed in $`H`$-band by Lacy et al. (1999a), its magnitude is consistent with its provisional redshift of 0.92. Only one object, 7C 1820+6657, was undetected in our study; its faint magnitude is consistent with its provisional redshift of 2.98.
In summary, we now have photometry on all nine of the high redshift radio galaxies with either no redshift information or uncertain redshifts in the 7C-iii sample of 54 objects. All six objects with uncertain redshifts have near-infrared magnitudes consistent with their tentative spectroscopic redshifts. Of the three objects with no redshift information, two have photometry consistent with them having redshifts $`1.5`$, and the one remaining object (7C 1748+6703) is probably at $`z3`$.
### 3.3 Emission line contamination of broad-band magnitudes
Line contamination of continuum magnitudes is expected to be small in nearly all the 7C objects, as the correlation of emission line and radio luminosities means that, on average, emission line fluxes will be much less than those from objects in radio-brighter samples. Whereas the $`z\stackrel{>}{_{}}2`$ objects studied by Eales & Rawlings (1996) typically have line contaminations of a few tens of percent, the 7C objects, which are selected at radio fluxes about four to five times lower, should have line contaminations $`\stackrel{<}{_{}}10`$ %. We have listed in Table 2 the estimated percentage contribution to the observed flux for all objects with bright emission lines in the observed band. These are based on the spectroscopy of Lacy et al. (1999b), using the median Ly$`\alpha /`$H$`\alpha `$+\[Nii\] ratio of 0.7 in Eales & Rawlings (1996), the line ratios in the composite radio galaxy spectrum of McCarthy (1993) and the H$`\alpha `$/H$`\beta `$ ratios of Koski (1978) to estimate the strengths of lines falling in the observed near-infrared. As expected, most objects have emission line contributions of $`<10`$ % . For only one object, 7C 1802+6456, is the emission line contamination expected to be $`>`$20 % .
### 3.4 Estimation of half-light radii
The seeing during most of our IRTF observations was very good, and it was clear from looking at the images that there was a wide range in the scale sizes of the hosts. We therefore selected a subsample consisting of the 15 sources with grade $`\alpha `$ or $`\beta `$ redshifts in the range $`0.8<z<2.7`$ in Lacy et al. (1999b) for which we had IRTF images. Although not complete, this sample should be representative of the 7C objects in this redshift range. To estimate the scale sizes we convolved model elliptical (de Vaucouleurs) galaxies having a range in half-light radius $`r_{\mathrm{hl}}`$ from 0.1 to 6.7 arcsec and zero ellipticity with a PSF from a nearby star in each IRTF image. The model was subtracted from the data and the minimum in the sum of the squares of the residuals found, along with the approximate $`\pm 1\sigma `$ range. An example of this process is shown in Fig. 2. The results were also checked by eye to see that the fits were reasonable. Nearby companion objects and discrete sub-components were subtracted prior to the fits. Although this is far from a full host galaxy model (in particular we made no attempt to discriminate between disc and de Vaucoulours fits), we feel that it should give a fair estimate of the scale size, and was all that the signal:noise in the images typically justified.
To obtain a major axis scale size to compare to other work, we need to assume a mean ellipticity, $`e`$, and correct our scale sizes with this. Roche et al. and Govoni et al. (2000) derive very similar mean values of $`e0.2`$ for the 6C radio galaxies at $`z1`$ and nearby radio galaxy hosts respectively, so this value has been assumed to correct the mean scale sizes quoted in Section 4 by multiplying the scale size by the square root of the axial ratio corresponding to $`e=0.2`$, namely $`1.06`$.
For two objects we have been able to compare the scale sizes with those measured on images taken with the Hubble Space Telescope (HST). 7C 1758+6719 ($`z=2.70`$) was observed with WFPC2 through the F702W filter; details of the observations will be given in a future paper (Lacy et al. in preparation). For this rest-frame UV image an exponential disk with $`r_{\mathrm{hl}}=0.21`$ arcsec was found to be a good fit to the radial profile. This can be compared with the estimate of $`r_{\mathrm{hl}}`$ from the IRTF image of 0.35 arcsec with a $`1\sigma `$ range of 0.12 – 1.0 arcsec, which we consider fair agreement. 7C 1754+6420 $`(z=1.09)`$ was observed through the F675W filter (Ridgway & Lacy in preparation). Again, the rest-frame UV emission was better fit by a disk than a de Vaucouleurs profile. In this case the scale size came out significantly smaller than the estimate from the IRTF image, 0.7 arcsec in the HST image compared to a best fit of 4.5 arcsec and a range of 1.6 – 13 arcsec for the IRTF image. This was despite the subtraction of a nearby companion to the south which was successfully removed from both the HST and IRTF images. The cause of this discrepancy is not clear. This image had the worst seeing of any of our IRTF images (1.1 arcsec), but the scale size measured in the IRTF image is much larger than the seeing HWHM. The radial profile of the IRTF image is, however, not at all well fit by the disk model from the HST image. Therefore perhaps the most likely cause of this discrepancy is that the galaxy has a UV-bright disk component embedded in a much larger scale-size elliptical.
### 3.5 Other data from the literature
To increase the number of high redshift objects in our study we have added the results of photometry of $`z25`$ radio galaxies by van Breugel et al. (1998). We also estimated scale sizes from their figure 2 for objects which appeared to be dynamically-relaxed ellipticals (seven out of the eight of their radio galaxies with $`z<3`$). We applied $`k`$\- and aperture corrections in a consistent manner to that for the 7C-iii radio sources. We have also added those $`1.6<z<2.4`$ objects without a strong point source contribution from the NICMOS study of radio galaxies by Pentericci (1999; and Pentericci et al. 2000), including the estimates of scale sizes for the five out of the nine objects in this redshift range for which Pentericci considers it possible to fit meaningful scale sizes. We also include the low luminosity $`z=4.42`$ radio galaxy VLA 123642+621331 discovered in the flanks of the Hubble Deep Field (Waddington et al. 1999) and 53W002 at $`z=2.239`$ (Windhorst, Mathis & Keel 1992).
The resulting high redshift (HZ) sample is detailed in Table 3. It mostly contains objects of radio luminosity comparable to the 3C radio galaxies at $`z1`$. As these objects (particularly those at $`z>3`$) were generally observed outside of rest-frame $`R`$-band, and are mostly not from complete samples, the results including the HZ sample will not be as reliable, though as we shall see the same trends seem to exist with or without this sample.
## 4 Discussion
### 4.1 Radio-luminosity and redshift dependence of absolute magnitudes
In Fig. 3 we plot the rest-frame $`R`$-band magnitudes against redshift for several samples of radio galaxies. We plot the $`z>0.8`$ 7C-iii objects with spectroscopic redshifts and infrared imaging, $`z>0.8`$ 6C radio galaxies with photometry from Eales et al. (1997), $`z>0.8`$ 3C radio galaxies in the Laing, Riley & Longair (1983; LRL) complete sample with photometry from Best, Longair & Röttgering (1998) \[apart from 3C22 which is a lightly-reddened quasar (Rawlings et al. 1995) and therefore excluded from the sample, 3C 175.1 ($`z=0.92`$) which has photometry from Ridgway & Stockton (1997) and 3C 263.1 ($`z=0.824`$) which has photometry from Eales (personal communication)\]. We have also added local radio galaxies in LRL from Owen & Laing (1989) and the HZ sample of galaxies of Table 3. All these objects, with the exception of three FRIs out of the 24 objects in the Owen & Laing LRL sample are either FRII or compact steep-spectrum sources. All the $`z>0.8`$ 7C-iii sources at $`z>0.8`$ are well above the FRI/FRII boundary in radio luminosity, with $`L_{R(151)}\stackrel{>}{_{}}10^{26}\mathrm{WHz}^1\mathrm{sr}^1`$ compared to the FRI/FRII boundary at $`L_{R(151)}10^{25}\mathrm{WHz}^1\mathrm{sr}^1`$.
There is a clear trend for redshift and absolute magnitude to correlate, even if the incomplete HZ sample is excluded. The 3C, 6C and 7C samples are all complete samples, selected on the basis of low frequency radio flux only, and are nearly completely identified. Thus the only selection effect which needs to be considered for these samples is the tendency for redshift and luminosity to correlate within each flux limited sample. With the wide range in radio luminosities in the complete samples at $`z1`$, however, we can separate out the luminosity dependence. This is illustrated in Fig. 4, where we have plotted absolute magnitude against radio luminosity for 3C, 6C and 7C galaxies in the redshift range $`0.8<z<1.4`$. In this redshift range the mean magnitude in 3C is $`M_R=24.17\pm 0.13`$, whereas in 6C it is $`23.79\pm 0.14`$ and in 7C-iii $`23.78\pm 0.15`$. This suggests that only the most radio-luminous objects have slightly brighter (by $`0.4`$ mag) hosts \[see also Rawlings et al. (1998) where preliminary $`K`$-band photometry on the 7C-i and 7C-ii samples is presented\].
We can think of two possible explanations for this radio luminosity dependence of host magnitude at $`z1`$. The first is a contribution from AGN-related light, e.g. emission lines and reddened and/or scattered quasar light (Eales & Rawlings 1996). Rigler et al. (1992), however, show that the fraction of emission from $`z1`$ 3C radio galaxies which is aligned with the radio jet axis, and therefore probably closely related to the AGN activity, contributes only $`10`$% of the light in the observed $`K`$-band. Also, Simpson, Rawlings & Lacy (1999) argue that reddened quasar light is responsible for $`\stackrel{<}{_{}}10`$% of the observed near-infrared emission from 3C radio galaxies on the basis of 3$`\mu `$m imaging. Thus we believe that only $`1020`$ per cent of the host luminosity can be accounted for by the AGN, not the $`40`$ per cent required to explain the higher host luminosities of the 3C galaxies. The second possibility is that the radio luminosity may be correlated with the mass of the host. This could arise, for example, if correlation of black hole mass and the mass of the spheroidal component of galaxies which is claimed to be present at low redshift (e.g. Magorrian et al. 1998) is appearing at the highest radio luminosities (Roche et al. 1998). We have argued (Lacy, Ridgway & Trentham 2000; see also Willott et al. 1999) that most radio galaxies are sub-Eddington accretors, but at the highest luminosities corresponding to the highest black hole masses, even the radio galaxies may be accreting at near Eddington rates. Thus at these highest luminosities we might expect a host galaxy mass – AGN luminosity (and therefore host luminosity) correlation to appear. Also, the mass of the host may correlate with the density of the intergalactic medium confining the radio source, which would enhance the radio luminosity.
We have made a crude model of the luminosity dependence by assuming a contribution from AGN-related light to a base host optical luminosity $`L_0`$. We model this contribution as a combination of a component due to emission lines and another component due to other AGN-related light. (We separate the emission line contribution, as the equivalent width of the emission lines is a strong function of redshift, unlike the other AGN-related emission which we expect to be much less redshift dependent.) Both these components should be proportional to the AGN luminosity, which we assume to be proportional to $`L_R(151)`$ to the power 0.8 (Serjeant et al. 1998; Willott et al. 1999), i.e.
$$L_{\mathrm{host}}(z)=L_0(z)+(\alpha +\beta (z))L_{R(151)}^{0.8},$$
(1)
where $`\alpha L_{R(151)}^{0.8}`$ is the contribution of AGN-related light other than emission lines and $`\beta (z)L_{R(151)}^{0.8}`$ is the contribution due to emission lines. The base host luminosity at $`z1`$ was set to $`L_0(1)=10^{23}\mathrm{WHz}^1`$, or $`M_R23.7`$, close to the mean of the 7C and 6C data, and $`(\alpha +\beta (1))`$ fit by eye to be $`3.5\mathrm{W}^{0.2}\mathrm{Hz}^{0.2}`$ (Fig. 4). The emission line contribution to the total near-infrared luminosities of the 3C radio galaxies at these redshifts is $`10`$% (Rawlings et al. 1997; Rawlings Eales & Lacy 1990), or about 20% of the overall luminosity-dependent correction. We have therefore set $`\alpha =2.8\mathrm{W}^{0.2}\mathrm{Hz}^{0.2}`$ and $`\beta (1)=0.7\mathrm{W}^{0.2}\mathrm{Hz}^{0.2}`$.
Using this formula we can then correct the host magnitudes for AGN-related light and produce plot of $`L_0`$ versus redshift. \[Rather than attempt to model the redshift dependence of $`\beta `$ we have used emission line contamination estimates from Eales & Rawlings (1997) and van Breugel et al. (1998) to explicitly correct the magnitudes in the 6C and high redshift samples (where these were unknown the object was omitted from the plot and correlation analysis), and our own estimates to correct the 7C-iii magnitudes.\] The results are shown in Fig. 5, which shows that there is indeed a correlation of $`L_0`$ with redshift. This was confirmed using the Kendall Tau statistic generalized to include limits, as implemented in the iraf.stsdas program bhkmethod (Isobe & Feigelson 1990), which returned a $`\tau =0.52`$ for the 106 objects from the 3C, 6C, 7C and HZ samples after correction. This corresponds to a probability that there is no correlation between $`L_0`$ and redshift of $`<0.0001`$. However, this probability increases to 0.13 if the HZ sample is removed. We have also conducted the same statistical test before correction for radio luminosity-dependent effects in equation (1). This gave a $`\tau =0.95`$ and a probability of no correlation of $`<0.0001`$, removing the HZ sample reduced this to $`\tau =0.66`$, again with the probability of no correlation of $`<0.0001`$.
### 4.2 The dispersion in the host luminosities
The low scatter in the magnitudes of radio galaxies in the $`Kz`$ Hubble Diagram out to $`z\stackrel{>}{_{}}3`$ has long been used as an argument for a high redshift of radio galaxy formation (e.g. Lilly 1989). This low scatter is thought to come about because radio galaxies form at $`z>3`$ then evolve along similar passive evolution tracks to end at radio galaxy hosts today. With a few modifications this simple picture still seems to be basically valid. Fig. 3 and Table 4 show that the dispersion in absolute magnitudes to $`z4`$ is indeed very low. After correction for AGN-contamination, however, the dispersion in $`L_0`$ for the $`z>1.8`$ objects seems to be higher than for the uncorrected objects (Table 4). As discussed above, the correlation of absolute magnitude with redshift is also less tight after correction.
At first sight it seems paradoxical that correction for luminosity-dependent effects should increase the scatter in the absolute magnitudes. However, this can be be understood if the luminosity-dependent contributions $`(\alpha +\beta )L_{R(151)}^{0.8}`$ are approximately equal to the highest values of $`L_0`$. (The $`z>1.8`$ objects have a very narrow range in $`L_R`$ as the range in flux of these objects is low, all except VLA 123642+621331 having 151 MHz fluxes between 0.5 and 10 Jy.) The addition of the luminosity-dependent contributions can then boost the objects with low $`L_0`$ by a significant factor, whereas the addition to the objects with high $`L_0`$ results in only a relatively small fractional increase in the luminosity (see also de Vries 1999). The result of this is a reduction of the scatter in the total luminosities.
We have used the F-test for variances to examine the statistical significance of the increase in scatter of the $`L_0`$ values. Comparing the objects with $`0.8<z<1.8`$ to those with $`1.8<z<2.8`$, we find that the increase in the sample standard deviation from $`\sigma _{n1}=0.53`$ to $`\sigma _{n1}=0.93`$ is significant at the 0.1 % level. At $`z>2.8`$, our sample contains two upper limits on the magnitudes, so we give a lower limit to the scatter of $`\sigma _{n1}>0.97`$. This is again significantly higher than that in the range $`0.8<z<1.8`$ (at the 0.2 % level). Thus there is good evidence that the scatter in $`L_0`$ is increasing with redshift, suggesting that we are close to the formation epoch of at least a significant fraction of the radio galaxy population.
### 4.3 The scale sizes of the hosts
Roche et al. (1998) have shown that the scale sizes of the $`1<z<1.4`$ 6C galaxies are significantly smaller than their more radio luminous 3C counterparts as measured by Best et al. (1997), suggesting a strong dependence of scale size on AGN luminosity. In this section we re-examine this using our data on the 7C objects, and further data on 3C scale sizes published recently.
To compare with the 6C results, we have defined complete samples of 7C-iii and 3C radio galaxies from LRL in the redshift range $`0.8<z<1.4`$ (we took the lower limit of 0.8 so as to include more 3C and 7C objects; only two 6C objects lie in the redshift range $`0.8<z<1`$, and we do not feel their exclusion will significantly affect the results). For the 3C objects we have used size estimates from the 2D modelling of McLure & Dunlop (2000) or the Keck $`K`$-band images of Ridgway & Stockton (1997) where available, and otherwise from Best et al. (1997). No estimates were available for three of the sixteen objects in the 3C sample, and we have omitted 3C356 due to confusion over the true identification. One of the objects missing from the 3C sample were omitted from the sample of Best et al. (1997) (3C 263.1), and the other two, 3C 13 and 3C 368 have $`K`$-band structures which are clearly affected by AGN-related emission as they are closely aligned with the radio axes. Of the 7C-iii objects, only one, 7C 1816+6710, has no scale size estimate as its image was taken in poor seeing, although the scale size of 7C 1742+6346 is an upper limit only. All the 7C scale sizes were increased by a factor of 1.06 to correct for an expected mean ellipticity of 0.2.
In Fig. 6 we plot the scale sizes of the 3C, 6C and 7C radio galaxies in the samples defined above against host absolute magnitude. Although the scatter is large it does seem that the 7C-iii radio galaxies have scale sizes consistent with their magnitudes when compared to the 3C radio galaxies, although the 6C objects are plotting below the general trend. The mean scale size for the 3C sample is $`11.1\pm 1.4`$ kpc and that for the 6C sample $`4.7\pm 0.9`$kpc. For the 7C-iii sample, we find a median scale size of $`0.9\pm 0.4`$ arcsec, which, when corrected for a mean ellipticity of 0.2, corresponds to about $`8\pm 3`$ kpc. (Using the whole subset of 7C-iii sources in the range $`0.8<z<2.7`$ with measured scale sizes we obtained the same result.) The mean for the 7C-iii is thus between that of the 6C and 3C samples, and statistically consistent with both.
A further check is provided by the effective-radius – galaxy magnitude relationship discussed by Roche et al. (1998),
$$\mathrm{lg}(r_{hl}/\mathrm{kpc})=0.3(M_R+\mathrm{\Delta }R+20.01),$$
(2)
where $`\mathrm{\Delta }R`$ is the amount by which the hosts have brightened by passive evolution from $`z=0`$. To set $`\mathrm{\Delta }R`$, we used the mean scale size of the 3C galaxies in our $`0.8<z<1.4`$ sample of 11.1 kpc, and the mean magnitudes of the same objects using the photometry of Best et al. (1998). This gave $`\mathrm{\Delta }R=0.68`$. We have plotted the relationship of equation (2) on Fig. 6, for both $`\mathrm{\Delta }R=0`$ and $`\mathrm{\Delta }R=0.68`$. The 7C radio galaxies are consistent with this relation for $`\mathrm{\Delta }R=0.68`$ (within the scatter), but most of the 6C radio galaxies plot significantly below the line.
We therefore believe that a combination of poor seeing in both the Roche et al. (1998) and Best et al. (1998) studies, which lead to poor scale size estimates, combined with small number statistics in the Roche et al. study, can probably explain the result of Roche et al. without the requirement for scale-size to depend very strongly on radio luminosity, though clearly better images of both the 6C and 7C objects will be needed to be certain of this.
We have also investigated possible redshift dependences of scale size with the addition of data on objects in the HZ sample listed in Table 3. Fig. 7 shows the results of plotting the scale size in arcsec against redshift. The mean half-light radius for local hosts of classical double FRIIs from Owen & Laing (1989) and Govoni et al. (2000), $`13\pm 2.5`$ kpc, is also plotted as a line in this figure for three different cosmologies. There is some evidence for a weak trend for scale size to decrease with increasing redshift, but with a lot of scatter. In an attempt to reduce the scatter we have tried using a correction based on equation (1), by introducing a correction in the log of the scale size:
$$\mathrm{\Delta }\mathrm{lg}r_{\mathrm{hl}}=0.3(MM_0(z)),$$
where $`M`$ is the absolute $`R`$-band magnitude after correction for line contamination, and $`M_0(z)`$ is $`L_0(z)`$ converted to absolute magnitude. With this correction, using the 7C, 6C, 3C ($`0.8<z<1.4`$) and HZ samples there is a weak anticorrelation of scale size with redshift, significant at about the 5% level using the bhkmethod.
Inspection of Fig. 7, and the comparison to low redshift FRII hosts, shows that the scatter in scale sizes is large at all redshifts, and suggests that any evolution in scale size is weak. FRI hosts at low redshift, however, have significantly brighter hosts and larger scale sizes at a given radio luminosity (e.g. Owen & Ledlow 1994), suggesting that merging in the richer cluster environments associated with these objects may be significant if they are the descendents of more luminous sources at high redshift. An argument against this is provided by a recent study of brightest cluster galaxies (BCGs) by Burke, Collins & Mann (2000), who show that, when selection effects are properly taken into account, the BCGs of the most X-ray luminous clusters brighten with redshift (to a similar extent as the FRII hosts, in fact), and thus there is little evidence of significant merger activity for redshifts $`0<z<1`$. However, because the FRII hosts are significantly less luminous than BCGs, mergers that would enhance the luminosity of a BCG by only a few percent would enhance that of an FRII host by a much larger factor. Nevertheless, the lack of evidence for merger activity in even rich clusters suggests that FRII hosts, which are generally found in poor cluster environments (e.g. Wold et al. 2000), are probably not significantly affected by merging, at least at $`z<1`$.
### 4.4 Surface brightness evolution
To demonstrate that the redshift dependence of host magnitude is independent of the assumed cosmology we have plotted in Fig. 8 the effective surface brightness (defined here as the flux within $`r_{\mathrm{hl}}`$ divided by $`\pi r_{\mathrm{hl}}^2`$ converted into magnitudes) in the rest-frame $`R`$-band versus redshift for the sub-sample of 7C-iii galaxies with estimated scale sizes. We have compared this to the cosmological surface brightness dimming expected in a standard expanding Universe with a Friedmann-Robertson-Walker metric. The surface brightness decreases with redshift much less rapdly than the prediction for a non-evolving stellar population (solid line in Fig. 8). This demonstrates that the brightening of the hosts with redshift is a real evolutionary effect, and not one produced by an incorrect choice of $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$.
## 5 The evolution of radio galaxy hosts
The evolution of host magnitude and rest-frame surface brightness with redshift, the lack of strong evolution in scale sizes in the subset of objects with “undisturbed” morphologies, and the low scatter ($`<1`$ magnitude) in the absolute magnitudes up to at least $`z3`$ are all consistent with an early formation epoch ($`z\stackrel{>}{_{}}3`$) for most, and perhaps all, radio galaxy hosts. At first sight these results are in conflict with standard CDM-based models for galaxy formation, in which hierarchical structure formation occurs in a “bottom up” fashion with the most massive galaxies forming late. Radio galaxies at $`z2`$ are much brighter than the average quasar host predicted by the models of Kauffmann & Haehnelt (2000), for example. To compare our radio galaxy sample to these predictions we have assumed an equivalent quasar luminosity equal to the mean of the $`1<z<3`$ 7C quasars, $`M_B25`$ (Willott et al. 1998). For quasars of this luminosity, the Kauffmann & Haehnelt models predict a mean host magnitude of $`M_V21.5`$ compared to our radio galaxies with a luminosity-corrected median $`M_R23.5`$ (corresponding to $`M_V23`$). In contrast, the hosts of radio-quiet quasars at $`z2`$ seem to fit this model quite well (Ridgway et al. 2000).
Between $`z1.3`$ and $`z2.3`$, there is evidence for an increasing scatter in base host luminosity. This suggests we are close to the formation epoch at $`z2.3`$. There are two other pieces of evidence which point to strong evolution in the host galaxies and which appear at $`z\stackrel{>}{_{}}2.5`$. First, the observations of Pentericci (1999; and Pentericci et al. 2000) and van Breugel et al. (1998) show that the morphologies of radio galaxy hosts at $`z<2.5`$ are mostly relaxed ellipticals, whereas at higher redshifts ($`z>2.5`$ in the HST NICMOS $`H`$-band images of Pentericci et al. and $`z>3`$ in the ground-based $`K`$-band images of van Breugel et al.) hosts with clumpy structures aligned with the radio source axis become much more common<sup>1</sup><sup>1</sup>1van Breugel et al. argue that this is not simply due to the $`z>3`$ objects being observed at shorter rest-frame wavelengths, as lower redshift objects observed at the same rest-frame wavelengths as the $`z>3`$ objects still look like relaxed ellipticals.. Second, the detection rate of continuum submillimeter emission from hot dust seems to rise rapidly at $`z>2.5`$ (Archibald et al. 2000). The 6C and 7C samples are completely identified with luminous host galaxies out to at least $`z3`$, however, which would suggest the bulk of the stars had formed and/or merged onto the host by $`z3`$. One way around a high formation redshift for all objects would be to assume a systematic delay between a host forming the bulk of its stars and the switching on of powerful radio jets (for example, if the black hole needs time to accrete enough mass). Radio galaxies would then always be seen with luminous hosts. However, the strong evolution in the hosts seen by van Breugel et al. around $`z3`$ suggests that their nature is still changing up to at least that epoch.
Taken together, the evidence suggests that the epoch of radio galaxy host formation was essentially over by $`z3`$, in the sense that most objects had formed the bulk of their stars, although further star-formation and merger activity continued in many objects up to $`z2`$. As radio galaxies can be found out to $`z>5`$, however (van Breugel et al. 1999), and there is no sign of a substantial drop in the space density of radio galaxies out to $`z>4`$ relative to their peak space density at $`z2.5`$ (Jarvis et al. 1999), this formation epoch must last at least a Gyr. In contrast, field ellipticals with evolved stellar populations and luminosities $`L_{}`$ seem to become rarer at $`z>1`$ (Zepf 1998; Barger et al. 1999; Dickinson 1999), suggesting a later formation epoch, perhaps $`z12`$.
A recent variation on the hierarchical models, the so-called Anti-Hierarchical Baryonic Collapse Model (Granato et al. 2000) may be able to explain these observations. This predicts that the massive hosts of powerful AGN may have formed at $`z3`$. This model uses the same CDM assumptions as the standard scenarios, but assumes that the dense gas in the most massive haloes can collapse and form stars in a massive spheroidal component more rapidly than that in the smaller haloes associated with normal galaxies. Star formation is switched off in these objects by winds from the AGN and the starburst at high redshift. If radio galaxy hosts are indeed typical of the most massive elliptical galaxies and their progenitors this model should be applicable to them, and would explain their anomalous evolution relative to generally less luminous field ellipticals.
We thank Wim de Vries, Michael Gregg and Wil van Breugel helpful discussions, the referee for a useful report, and the staff at the IRTF and Lick Observatory for their assistance. We are also very grateful to Chris Willott for providing the photometric redshift estimates for those 7C-iii radio galaxies without spectroscopic redshifts prior to publication. The IRTF is operated by the University of Hawaii on behalf of NASA. AJB acknowledges support from the Cambridge Institute of Astronomy PPARC observational rolling grant, ref. no. PPA/G/O/1997/00793, and a NICMOS postdoctoral fellowship while at Berkeley (grant NAG 5-3043). This work was performed under the auspices of the U.S. Department of Energy by University of California Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48, and was partly based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract No. NAS5-26555. |
warning/0003/hep-th0003038.html | ar5iv | text | # LMU-TPW 00-8UAHEP 00-3hep-th/0003038 On the correspondence between gravity fields and CFT operators
## 1 Introduction
According to the AdS/CFT correspondence fields of type IIB supergravity on the $`AdS_5\times S^5`$ background are dual to gauge invariant operators in $`D=4`$, $`𝒩=4`$ supersymmetric Yang-Mills theory (SYM<sub>4</sub>) which belong to short representations of the conformal superalgebra $`SU(2,2|4)`$ and have protected scale dimensions. The short representations are generated by chiral primary operators (CPOs) transforming in the $`k`$-traceless symmetric representations of $`SO(6)`$. It is well-known that single-trace operators $`O_k^I=\mathrm{tr}(\varphi ^{(i_1}\mathrm{}\varphi ^{i_k)})`$ are chiral, and it was shown in that multi-trace operators of the form $`O_{k_1\mathrm{}k_n}^I=\mathrm{tr}(\varphi ^{(i_1}\mathrm{}\varphi ^{i_{k_1}})\mathrm{}\mathrm{tr}(\varphi ^{j_1}\mathrm{}\varphi ^{j_{k_n})})`$ are chiral too. There are also CPOs which are normal-ordered products of single- and multi-trace CPOs and their descendents. Thus, in general, CPOs are admixtures of single- and multi-trace operators with the same (protected) conformal dimension.
On the other hand the particle spectrum of type IIB supergravity on $`AdS_5\times S^5`$ contains only one set of fields which can couple to CPOs. These fields $`s^I`$ are mixtures of the five form field strength and the trace of the graviton on the sphere. Thus, one should understand which linear combinations of CPOs are dual to the gravity fields $`s^I`$. Although, it is customarily believed that they are dual to single-trace operators $`O_k^I`$, no complete reliable proof of this fact is known.
A way to solve the problem is to compute correlation functions in free field theory and in the supergravity approximation, and to compare them. Of course, one can compare only correlation functions subject to non-renormalization theorems. According to , to compute n-point functions in SYM<sub>4</sub> one has to know the type IIB supergravity action on $`AdS_5\times S^5`$ up to the $`n`$-th order. The quadratic action for physical fields was found in by using the “covariant” action of . The first step in finding interaction vertices was made in where quadratic and cubic actions for the scalars $`s^I`$ were found by expanding the covariant equations of motion for type IIB supergravity up to the second order. By using the actions, all 3-point functions of normalized CPOs dual to $`s^I`$ were computed, and, for generic values of conformal dimensions of the CPOs, appeared to coincide with 3-point functions of the single-trace CPOs $`O_k^I`$ calculated in the free field theory. It was conjectured in that the 3-point functions are not renormalized, and this was later proven in . One might conclude on the basis of this coincidence that the fields $`s^I`$ are dual to the single-trace CPOs. However, as was noted in , a 3-point function of CPOs computed in the supergravity approximation vanishes in the extremal case, for which the sum of conformal dimensions of two operators equals the conformal dimension of the third operator, e.g. $`k_1=k_2+k_3`$, because of the vanishing of the cubic couplings of the dual scalar fields.
There were proposed three different ways to resolve the puzzle. According to , to compute extremal 3-point functions, one should first analytically continue in the conformal dimensions $`k_1,k_2,k_3`$. Then, since the gravity coupling is proportional to $`k_2+k_3k_1`$, and the AdS integral behaves itself as $`1/(k_2+k_3k_1)`$, one obtains a finite extremal 3-point function. However, from the computational point of view the procedure of analytical continuation looks superfluous, because no actual singularity is involved. An extremal 3-point function vanishes due to the absence of the corresponding cubic coupling, and one does not have to evaluate any AdS integral.
In we explained the vanishing of the extremal cubic couplings by noting that the scalars $`s^I`$, and, in general, supergravity fields, may be dual to extended CPOs which are admixtures of single- and multi-trace CPOs. Nevertheless, the fact that the analytical continuation procedure seems to work in all known examples,<sup>1</sup><sup>1</sup>1In particular, this procedure works in the case of 3-point functions of operators dual to two scalars $`s^I`$ and a supergravity field, computed in , where the cubic couplings also vanish in extremal cases. allows one to assume that, in the large $`N`$ limit and for generic values of conformal dimensions, correlation functions of extended CPOs coincide with the ones of the single-trace CPOs. However, it is also clear that the analytical continuation procedure may work only in the large $`N`$ limit, because for finite $`N`$ only the single-trace CPOs $`O_k^I`$ with $`2kN`$ are independent, and a single-trace CPO with $`k>N`$ is equal to a linear combination of multi-trace CPOs. This also shows that the appearence of multi-trace CPOs is unavoidable for finite $`N`$.
Other arguments in favour of the proposal come from the study of quartic couplings of the scalars $`s^I`$ performed in . It is shown there that the quartic couplings vanish in the extremal case for which $`k_1=k_2+k_3+k_4`$. As was pointed out in the vanishing of extremal couplings is dictated by the AdS/CFT correspondence because in this case contact Feynman diagrams are ill-defined, and therefore, non-vanishing extremal quartic couplings would contradict to the AdS/CFT correspondence. By the same reason 2- and 4-derivative quartic couplings have to vanish in the subextremal case for which $`k_1=k_2+k_3+k_42`$, and 4-derivative quartic couplings should vanish in the sub-subextremal case when $`k_1=k_2+k_3+k_44`$. The vanishing of extremal couplings means that 4-point extremal correlators of CPOs dual to the scalars $`s^I`$ vanish, and, therefore, the scalars correspond not to single-trace CPOs but to extended CPOs.
Then, it is shown in that the quartic action is consistent with the Kaluza-Klein (KK) reduction down to five dimensions, and admits a truncation to the massless multiplet, which can be identified with the field content of the gauged $`𝒩=8`$, $`d=5`$ supergravity . Consistency means that there is no term linear in massive KK modes in the untruncated supergravity action, so that all massive KK fields can be put to zero without any contradiction with equations of motion. From the AdS/CFT correspondence point of view the consistent truncation implies that $`any`$ $`n`$-point correlation function of $`n1`$ operators dual to the fields from the massless multiplet and one operator dual to a massive KK field vanishes because, as one can easily see there is no exchange Feynman diagram in this case. This in particular implies that the scalars $`s^I`$ are dual to extended CPOs. Indeed, if we assume that the scalars $`s^I`$ correspond to the single-trace CPOs $`O_k^I`$, we derive from the consistency of the KK reduction that correlators of the form $`O_2^{I_1}O_2^{I_2}\mathrm{}O_2^{I_{n1}}O_k^{I_n}`$ vanish for $`k4`$, that is not the case for such correlators of single-trace CPOs.
Finally, the third way of solving the puzzle was proposed in , where it was noted that since the scalars $`s^I`$ used in differ from the original scalars appearing in the covariant equations of motion for type IIB supergravity, one could obtain a nonvanishing extremal 3-point function by using an action for the original scalars. This was demonstrated for fields from the descendent sector where the relevant part of the type IIB supergravity action is known. The action used in contains higher-derivative terms, and although the bulk extremal couplings vanish on shell, there appear boundary terms which provide nonvanishing contribution to extremal 3-point functions. However, as was also noted in , one can make a nonlinear off-shell transformation of the gravity fields and remove all higher-derivative terms and all nonvanishing (off-shell) extremal couplings. No boundary terms appear as a result of the field transformation, and the transformed action leads to vanishing extremal correlators.
Thus, these results seem to indicate that although the original gravity fields may be dual to single-trace operators, the transformed fields are already dual to mixtures of single- and multi-trace operators. From this point of view a redefinition of the gravity fields corresponds to a change of an operator basis in CFT.
To justify this point of view we study how 4-point correlation functions are changed under derivative-dependent gravity fields redefinitions of the form used in to reduce the non-Lagrangian equations of motion to a Lagrangian form. The field transformations discussed in are their particular case. We begin with the quartic action for scalars $`s^I`$ found in and show that these transformations indeed can change some correlators, in particular, the extremal 3- and 4-point functions and the subextremal 4-point functions for which $`k_1=k_2+k_3+k_42`$.
The subextremal correlators are of special interest because, as has been shown in , and checked in to first order in perturbation theory, they are not renormalized. Thanks to the non-renormalization theorem one can also employ subextremal 4-point correlators to test the AdS/CFT correspondence. In particular the non-renormalization implies that the subextremal quartic couplings vanish, and we show that this is indeed the case. This fact together with the absence of the exchange Feynman diagrams leads to the vanishing of the subextremal 4-point functions of extended CPOs dual to the scalars $`s^I`$.
We show that any field redefinition induces a change of correlation functions which is always given by a product of 2- and 3-point functions. By this reason, and due to the non-renormalization theorems for extremal and subextremal correlators, it seems possible to find such a field transformation that the extremal and subextremal 3- and 4-point functions coincide with the ones of single-trace CPOs.
As a by-product of our study, we also find that if conformal dimensions of at least two CPOs do not coincide then some structures in a 4-point function of these CPOs can be also changed by a field redefinition. Thus, although such a 4-point function in general is not protected by a non-renormalization theorem, this seems to be an indication that the coefficients of the changing structures of the 4-point function are not renormalized.
The plan of the paper is as follows. In section 2 we recall the definition of normalized single-trace CPOs and extended CPOs, and discuss the general properties of the supergravity action used to compute 3- and 4-point functions of the extended CPOs. In section 3 we study how the 3- and 4-point correlation functions change under a derivative-dependent field redefinition. In appendix we show that the quartic couplings of vanish in the subextremal case, and that 4-derivative couplings vanish in the sub-subextremal case.
## 2 Extended CPOs and quartic supergravity action
We follow defining the normalized single-trace CPOs as
$`O^I(\stackrel{}{x})={\displaystyle \frac{(2\pi )^k}{\sqrt{k\lambda ^k}}}C_{i_1\mathrm{}i_k}^I\mathrm{tr}(\varphi ^{i_1}(\stackrel{}{x})\mathrm{}\varphi ^{i_k}(\stackrel{}{x})),`$ (2.1)
where $`C_{i_1\mathrm{}i_k}^I`$ are totally symmetric traceless rank $`k`$ orthonormal tensors of $`SO(6)`$: $`C^IC^J=C_{i_1\mathrm{}i_k}^IC_{i_1\mathrm{}i_k}^J=\delta ^{IJ}`$, and $`\varphi ^i`$ are scalars of SYM<sub>4</sub>.
The two- and three-point functions of CPOs can be easily computed in free field theory and in the large $`N`$ limit, and are given by
$`O^I(\stackrel{}{x})O^J(\stackrel{}{y})={\displaystyle \frac{\delta ^{IJ}}{|\stackrel{}{x}\stackrel{}{y}|^{2k}}},`$ (2.2)
$`O^{I_1}(\stackrel{}{x})O^{I_2}(\stackrel{}{y})O^{I_3}(\stackrel{}{z})={\displaystyle \frac{1}{N}}{\displaystyle \frac{C^{I_1I_2I_3}}{|\stackrel{}{x}\stackrel{}{y}|^{2\alpha _3}|\stackrel{}{y}\stackrel{}{z}|^{2\alpha _1}|\stackrel{}{z}\stackrel{}{x}|^{2\alpha _2}}},`$ (2.3)
where $`\alpha _i=\frac{1}{2}(k_j+k_lk_i)`$, $`jli`$, $`C^{I_1I_2I_3}=\sqrt{k_1k_2k_3}C^{I_1}C^{I_2}C^{I_3}`$, and $`C^{I_1}C^{I_2}C^{I_3}`$ is the unique $`SO(6)`$ invariant obtained by contracting $`\alpha _1`$ indices between $`C^{I_2}`$ and $`C^{I_3}`$, $`\alpha _2`$ indices between $`C^{I_3}`$ and $`C^{I_1}`$, and $`\alpha _3`$ indices between $`C^{I_2}`$ and $`C^{I_1}`$. As was discussed in the Introduction, single-trace CPO cannot be dual to the scalar fields $`s^I`$ used in to compute their 3-point functions. However, as was shown in , one can define an extended CPO which corresponds to a scalar $`s^I`$ by adding to a single-trace CPO a proper combination of multi-trace CPOs:
$`\stackrel{~}{O}^{I_1}=O^{I_1}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{I_2+I_3=I_1}{}}C^{I_1I_2I_3}O^{I_2}O^{I_3}.`$ (2.4)
One can easily check that in the large $`N`$ limit these operators have the normalized two-point functions (2.2), the three-point functions (2.3) in the non-extremal case, and vanishing three-point functions in the extremal case. Note that these operators require further modification to be consistent with all $`n`$-point functions computed in the framework of the AdS/CFT correspondence. In general, an extended CPO is a linear combination of a CPO and chiral primary composite operators which are normal-ordered products of CPOs and their descendants.
The quartic action for the scalars $`s^I`$ dual to the extended CPOs was found in , and the part of the action depending only on the scalars can be written in the form
$`S`$ $`=`$ $`{\displaystyle _{AdS_5}}({\displaystyle \frac{1}{2}}(_as_I^as_I+m_I^2s_I^2)+\lambda _{IJK}s_Is_Js_K+\lambda _{IJKL}^{(0)}s_Is_Js_Ks_L`$ (2.5)
$`+`$ $`\lambda _{IJKL}^{(2)}_as_I^as_Js_Ks_L+\lambda _{IJKL}^{(4)}_as_I^as_J_bs_K^bs_L)`$ (2.6)
$`=`$ $`{\displaystyle _{AdS_5}}(s^I).`$ (2.7)
Since the action does not contain higher-derivative terms, the Hamiltonian reformulation of the quartic action is straightforward, and, therefore, as was shown in , there is no need to add boundary terms.
There are also cubic terms describing the interaction of the scalars $`s^I`$ with other scalars, with vector fields, and with massive symmetric tensor fields of the second rank, but we omit them for the sake of simplicity.
Considering the contribution of contact Feynman diagrams to 3- and 4-point functions, one can easily observe that the integrals over the $`AdS_5`$ space diverge in several cases: $`(i)`$ if cubic couplings do not vanish in the extremal case for which, e.g. $`k_1=k_2+k_3`$, $`(ii)`$ if quartic couplings do not vanish in the extremal case when $`k_1=k_2+k_3+k_4`$, $`(iii)`$ if 4-derivative and 2-derivative quartic couplings do not vanish in the subextremal case when $`k_1=k_2+k_3+k_42`$, and $`(iv)`$ if 4-derivative quartic couplings do not vanish in the sub-subextremal case for which $`k_1=k_2+k_3+k_44`$. Thus the AdS/CFT correspondence requires vanishing all these couplings. Moreover, although the AdS integral involved in the subextremal non-derivative quartic graph does not diverge, the non-derivative quartic couplings also have to vanish in the subextremal case, because as was proven in the subextremal 4-point functions are non-renormalized, and, therefore, have a free field form (a product of 2- and 3-point functions of a free CFT). On the other hand a nonvanishing quartic subextremal coupling would lead to a 4-point function which does not have a free field form, and, this would contradict to the AdS/CFT correspondence.
Since one can easily show that all exchange Feynman diagrams vanish in the extremal and subextremal cases, the vanishing of the quartic couplings means that extremal and subextremal 4-point functions of operators dual to the scalars $`s^I`$ in (2.7) also vanish. This is certanly not the case for the correlators of the single-trace CPOs $`O_k^I`$, and, therefore, we interpret the scalars $`s^I`$ as to be dual to the extended CPOs of the form (2.4). However, to obtain action (2.7) a number of nonlinear derivative-dependent field redefinitions was performed. Thus, a natural question arises whether it is possible to make such a field redefinition of the scalars $`s^I`$ that the redefined scalars $`𝐬^I`$ would correspond to the single-trace CPOs. In the next section we study the response of 3- and 4-point correlation functions to such changes and show that the desirable field redefinitions may exist.
## 3 Field redefinitions and 4-point functions
According to the proposal by , the generating functional of connected Green functions in SYM<sub>4</sub> at large $`N`$ and at strong ’t Hooft coupling coincides with the on-shell value of the type IIB supergravity action on $`AdS_5\times S^5`$ subject to the Dirichlet boundary conditions imposed on supergravity fields at the boundary of $`AdS_5\times S^5`$. To have a well-defined functional of the boundary fields we cut the AdS space<sup>2</sup><sup>2</sup>2 We use the AdS metric of the form: $`ds^2=\frac{1}{z^2}(dz^2+dx_i^2)`$. off at $`z=\epsilon `$ and consider the part of AdS with $`z\epsilon `$. We impose the Dirichlet boundary conditions on the scalars $`s^I`$: $`s^I(\epsilon ,\stackrel{}{x})s^I(\stackrel{}{x})`$, and denote the on-shell value of the action (2.7) as $`S(s)`$. To compute 3- and 4-point functions we only need equations of motion for the scalars $`s_I`$ decomposed up to the second order in fields:
$$(_a^2m_I^2)s_I+3\lambda _{IJK}s_Js_K=0.$$
(3.1)
The solution of the equation that satisfies the Dirichlet boundary conditions can be written in the form
$$s_I=s_I^{(0)}+s_I^{(1)}.$$
(3.2)
Here $`s_I^{(0)}`$ solves the linear part of (3.1) with the Dirichlet boundary conditions at $`z=\epsilon `$, and $`s_I^{(1)}`$ has the vanishing boundary conditions at $`z=\epsilon `$, and solves the equation
$$(_a^2m_I^2)s_I^{(1)}+3\lambda _{IJK}s_J^{(0)}s_K^{(0)}=0.$$
(3.3)
The on-shell value of action (2.7) is obtained by substituting (3.2) into it:
$$S(s)=_\epsilon ^{\mathrm{}}𝑑z𝑑\stackrel{}{x}z^5(s_I^{(0)}+s_I^{(1)}).$$
(3.4)
Let us now consider the following off-shell transformation of the fields $`s_I`$
$`s_I`$ $`=`$ $`𝐬_I+C_{IJK}^{(0)}𝐬_J𝐬_K+C_{IJK}^{(2)}_a𝐬_J^a𝐬_K`$ (3.5)
$`+`$ $`C_{IJKL}^{(0)}𝐬_J𝐬_K𝐬_L+C_{IJKL}^{(2)}_a𝐬_J^a𝐬_K𝐬_L+C_{IJKL}^{(4)}_a𝐬_J_b𝐬_K^a^b𝐬_L`$ (3.6)
$`=`$ $`𝐬_I+(\delta s)_I.`$ (3.7)
This transformation is of the same form as the most general $`s`$-dependent field redefinition that was used in to reduce the original non-Lagrangian equations of motion to a Lagrangian form. We also assume that the constants $`C_{IJK}^{(0)}`$ and $`C_{IJK}^{(2)}`$ do not vanish only if any of the conformal dimensions does not exceed its extremal value: $`\mathrm{\Delta }_I\mathrm{\Delta }_J+\mathrm{\Delta }_K`$.
The equations of motion for the redefined fields look as follows
$$(_a^2m_I^2)𝐬_I+3\lambda _{IJK}𝐬_J𝐬_K+(_a^2m_I^2)(\delta _2s)_I=0,$$
(3.8)
where
$$(\delta _2s)_I=C_{IJK}^{(0)}𝐬_J𝐬_K+C_{IJK}^{(2)}_a𝐬_J^a𝐬_K.$$
(3.9)
The simplest way to study the influence of the field redefinition (3.7) on the 3- and 4-point functions is to impose on the redefined fields $`𝐬_I`$ the same boundary conditions as on $`s_I`$, and to compare the on-shell values of the original and transformed actions. Thus we write the solution to (3.8) in the form
$$𝐬_I=s_I^{(0)}+𝐬_I^{(1)}=s_I^{(0)}+s_I^{(1)}(\delta _2s)_I^{(0)}+\sigma _I^{(0)}.$$
(3.10)
Here
$$(\delta _2s)_I^{(0)}=C_{IJK}^{(0)}s_J^{(0)}s_K^{(0)}+C_{IJK}^{(2)}_as_J^{(0)}^as_K^{(0)},$$
(3.11)
and $`\sigma _I^{(0)}`$ solves the linear part of the equation (3.1), and satisfies the following boundary condition
$$\sigma _I^{(0)}|_{z=\epsilon }=(\delta _2s)_I^{(0)}|_{z=\epsilon }.$$
(3.12)
Substituting (3.10) into (3.7), we find
$`𝐬_I+(\delta s)_I=s_I^{(0)}+s_I^{(1)}+\sigma _I^{(0)}+(\delta _3s)_I^{(0)}+2C_{IJK}^{(0)}s_J^{(0)}𝐬_K^{(1)}+2C_{IJK}^{(2)}_as_J^{(0)}^a𝐬_K^{(1)},`$
where
$`(\delta _3s)_I^{(0)}=C_{IJKL}^{(0)}s_J^{(0)}s_K^{(0)}s_L^{(0)}+C_{IJKL}^{(2)}_as_J^{(0)}^as_K^{(0)}s_L^{(0)}+C_{IJKL}^{(4)}_as_J^{(0)}_bs_K^{(0)}^a^bs_L^{(0)}.`$
Thus the on-shell value of the transformed Lagrangian is given by
$`\stackrel{~}{}(𝐬_I)`$ $`=`$ $`(𝐬_I+(\delta s)_I)=(s_I^{(0)}+s_I^{(1)})`$ (3.13)
$``$ $`_a\left(s_I^{(0)}+s_I^{(1)}\right)^a\left(\sigma _I^{(0)}+(\delta _3s)_I^{(0)}+2C_{IJK}^{(0)}s_J^{(0)}𝐬_K^{(1)}+2C_{IJK}^{(2)}_bs_J^{(0)}^b𝐬_K^{(1)}\right)`$ (3.14)
$``$ $`m_I^2\left(s_I^{(0)}+s_I^{(1)}\right)\left(\sigma _I^{(0)}+(\delta _3s)_I^{(0)}+2C_{IJK}^{(0)}s_J^{(0)}𝐬_K^{(1)}+2C_{IJK}^{(2)}_as_J^{(0)}^a𝐬_K^{(1)}\right)`$ (3.15)
$``$ $`{\displaystyle \frac{1}{2}}_a\sigma _I^{(0)}^a\sigma _I^{(0)}{\displaystyle \frac{1}{2}}m_I^2\sigma _I^{(0)}\sigma _I^{(0)}+3\lambda _{IJK}s_I^{(0)}s_J^{(0)}\sigma _K^{(0)}.`$ (3.16)
The first term on the r.h.s. of this equation is equal to the on-shell value of the original Lagrangian (2.7), and other terms represent relevant corrections to it. By using equations of motion we can rewrite (3.16) as follows
$`\stackrel{~}{}(𝐬_I)`$ $`=`$ $`(s_I^{(0)}+s_I^{(1)})`$ (3.17)
$``$ $`_a\left(^a\left(s_I^{(0)}+s_I^{(1)}\right)\left(\sigma _I^{(0)}+(\delta _3s)_I^{(0)}+2C_{IJK}^{(0)}s_J^{(0)}𝐬_K^{(1)}+2C_{IJK}^{(2)}_bs_J^{(0)}^b𝐬_K^{(1)}\right)\right)`$ (3.18)
$``$ $`{\displaystyle \frac{1}{2}}_a\left(^a\sigma _I^{(0)}\sigma _I^{(0)}\right).`$ (3.19)
Thus the on-shell values of the original action and the redefined one only differ by a boundary term. Omitting nonessential terms, which cannot change the 3- and 4-point functions, this boundary term can be written in the form
$`I={\displaystyle 𝑑\stackrel{}{x}\epsilon ^{d+1}}`$ $`(`$ $`\epsilon ^2C_{IJK}^{(2)}_0s_I^{(0)}_0s_J^{(0)}_0s_K^{(0)}+\epsilon ^2C_{IJK}^{(2)}_0s_I^{(1)}_0s_J^{(0)}_0s_K^{(0)}`$ (3.20)
$`+`$ $`_0s_I^{(0)}\left(C_{IJKL}^{(2)}_as_J^{(0)}^as_K^{(0)}s_L^{(0)}+C_{IJKL}^{(4)}_as_J^{(0)}^bs_K^{(0)}^a_bs_L^{(0)}\right)`$ (3.21)
$`+`$ $`2\epsilon ^2C_{IJK}^{(2)}_0s_I^{(0)}_0s_J^{(0)}_0(s_K^{(1)}(\delta _2s)_K^{(0)}+\sigma _K^{(0)})+{\displaystyle \frac{1}{2}}_0\sigma _I^{(0)}\sigma _I^{(0)}).`$ (3.22)
The first term on the r.h.s. of (3.22) represents the change in a 3-point function induced by the field redefinition (3.7). It was shown in that this term gives a nonvanishing contribution only to the extremal 3-point functions, and always leads to a contribution of the free-field form. In particular, one can choose a field redefinition of such a form that all 3-point functions will coincide with the 3-point functions of the single-trace CPOs.
We are going to study the influence of the field redefinition on the 4-point functions and begin with considering the simplest term<sup>3</sup><sup>3</sup>3For the sake of generality we consider the boundary term (3.22) in $`d`$ dimensions, and for arbitrary scalars $`s_I`$ dual to operators of conformal dimensions $`\mathrm{\Delta }_I`$ with the only restriction $`\mathrm{\Delta }_Id/2`$.
$$I_1=𝑑\stackrel{}{x}\epsilon ^{d+1}C_{IJKL}^{(2)}_0s_I^{(0)}_as_J^{(0)}^as_K^{(0)}s_L^{(0)}.$$
(3.23)
It is obvious that only derivatives in the radial direction can give a nonvanishing contribution to a 4-point function, thus we can replace $`I_1`$ by
$$I_1=𝑑\stackrel{}{x}\epsilon ^{d+3}C_{IJKL}^{(2)}_0s_I^{(0)}_0s_J^{(0)}_0s_K^{(0)}s_L^{(0)}.$$
(3.24)
It is well-known that the Fourier transform of the solution of the Dirichlet problem is given by (see, e.g. )
$`s_I^{(0)}(z,\stackrel{}{k})=K_I(z,\stackrel{}{k})s_I(\stackrel{}{k}),`$
where
$$K_I(z,\stackrel{}{k})=\left(\frac{z}{\epsilon }\right)^{d/2}\frac{𝒦_\nu (kz)}{𝒦_\nu (k\epsilon )},\nu =\mathrm{\Delta }_I\frac{d}{2}$$
(3.25)
and $`𝒦_\nu (kz)`$ is a Macdonald function. By using this formula, we find
$$z_0K_I(z,\stackrel{}{k})|_{z=\epsilon }=d\mathrm{\Delta }_I+a_I(k\epsilon )^{2(\mathrm{\Delta }_I\frac{d}{2})}\mathrm{log}(k\epsilon )+\mathrm{},$$
(3.26)
where $`\mathrm{}`$ refer to terms which do not contribute to 4-point functions in the limit $`\epsilon 0`$. Thus a relevant contribution of $`I_1`$ to the Fourier transform of a 4-point function is proportional to
$`\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4)\epsilon ^{2\mathrm{\Delta }_I+2\mathrm{\Delta }_J+2\mathrm{\Delta }_K4d}k_1^{\mathrm{\Delta }_I\frac{d}{2}}k_2^{\mathrm{\Delta }_J\frac{d}{2}}k_3^{\mathrm{\Delta }_K\frac{d}{2}}\mathrm{log}(k_1\epsilon )\mathrm{log}(k_2\epsilon )\mathrm{log}(k_3\epsilon ).`$
This expression is always the Fourier transform of a product of three 2-point functions. It gives a contribution to a 4-point function only if <sup>4</sup><sup>4</sup>4Note that the correct scaling behaviour of a $`4`$-point function is $`𝒪\left(\epsilon ^{\mathrm{\Delta }_I+\mathrm{\Delta }_J+\mathrm{\Delta }_K+\mathrm{\Delta }_L4d}\right)`$.
$$2\mathrm{\Delta }_I+2\mathrm{\Delta }_J+2\mathrm{\Delta }_K4d=\mathrm{\Delta }_I+\mathrm{\Delta }_J+\mathrm{\Delta }_K+\mathrm{\Delta }_L4d,$$
i.e., in the extremal case $`\mathrm{\Delta }_I+\mathrm{\Delta }_J+\mathrm{\Delta }_K=\mathrm{\Delta }_L`$. Note that in the non-extremal cases $`\mathrm{\Delta }_L<\mathrm{\Delta }_I+\mathrm{\Delta }_J+\mathrm{\Delta }_K`$, in particular in the subextremal one, the boundary term scales too fast to give a contribution.
The second integral to be considered is
$$I_2=𝑑\stackrel{}{x}\epsilon ^{d+3}C_{IJL}^{(2)}_0s_I^{(0)}_0s_J^{(0)}_0\sigma _L^{(0)}.$$
(3.27)
To analyse the integral we use
$`\sigma _L^{(0)}(z,\stackrel{}{k})`$ $`=`$ $`K_L(z,\stackrel{}{k}){\displaystyle }d\stackrel{}{k}_3d\stackrel{}{k}_4\delta (\stackrel{}{k}_3+\stackrel{}{k}_4\stackrel{}{k})(C_{LMN}^{(2)}\epsilon ^2_0s_M^{(0)}(\epsilon ,\stackrel{}{k}_3)_0s_N^{(0)}(\epsilon ,\stackrel{}{k}_4)`$ (3.28)
$`+`$ $`C_{LMN}^{(2)}\epsilon ^2_is_M(\stackrel{}{k}_3)_is_N(\stackrel{}{k}_4)+C_{LMN}^{(0)}s_M(\stackrel{}{k}_3)s_N(\stackrel{}{k}_4)).`$ (3.29)
One can easily see that only the first term can give a nonvanishing contribution to a 4-point function. Combining (3.27) with (3.29) we get that this contribution to a 4-point function is proportional to
$`C_{IJL}^{(2)}C_{LMN}^{(2)}\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4)\epsilon ^{d+5}_0K_I(\epsilon ,\stackrel{}{k}_1)_0K_J(\epsilon ,\stackrel{}{k}_2)_0K_L(\epsilon ,\stackrel{}{k}_3+\stackrel{}{k}_4)`$ (3.30)
$`\times _0K_M(\epsilon ,\stackrel{}{k}_3)_0K_N(\epsilon ,\stackrel{}{k}_4).`$ (3.31)
By using (3.26), we see that there are several cases when we can get a nonlocal contribution: $`(i)`$ the five-logs case, $`(ii)`$ the four-logs case, and $`(iii)`$ the three-logs case. It is not difficult to show that there is no contribution in the five- and four-logs cases. Three-logs case has three subcases. The first one is
$$\mathrm{log}(k_3+k_4)\mathrm{log}k_3\mathrm{log}k_4\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4)$$
In this case we get $`\delta (x_2x_1)`$ after integrating over momenta.
The second case is
$$\mathrm{log}k_2\mathrm{log}k_3\mathrm{log}k_4\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4)$$
It is obvious that in this case we get a product of three 2-point functions, and a nonvanishing contribution will be only in the extremal case $`\mathrm{\Delta }_J+\mathrm{\Delta }_M+\mathrm{\Delta }_N=\mathrm{\Delta }_I`$.
The third case is
$$\mathrm{log}(k_3+k_4)\mathrm{log}k_2\mathrm{log}k_4\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4)$$
One can easily see that in this case we also obtain a product of three 2-point functions, and a nonvanishing contribution will be only if
$$\mathrm{\Delta }_I+\mathrm{\Delta }_J+2\mathrm{\Delta }_L\mathrm{\Delta }_M+\mathrm{\Delta }_N=0.$$
Taking into account that
$$\mathrm{\Delta }_I+\mathrm{\Delta }_J+\mathrm{\Delta }_L0,\mathrm{\Delta }_M+\mathrm{\Delta }_N+\mathrm{\Delta }_L0$$
we get that there is a solution to this equation if
$$\mathrm{\Delta }_L=\mathrm{\Delta }_I\mathrm{\Delta }_J=\mathrm{\Delta }_M\mathrm{\Delta }_N.$$
This is a new case, and it is tempting to assume that the coefficient of the corresponding structure in the 4-point function is not renormalized. To understand why this may be so it is instructive to write down the changing term in the 4-point function:
$`O^I(\stackrel{}{x}_1)O^J(\stackrel{}{x}_2)O^M(\stackrel{}{x}_3)O^N(\stackrel{}{x}_4){\displaystyle \frac{1}{x_{12}^{2\mathrm{\Delta }_J}x_{13}^{2\mathrm{\Delta }_L}x_{34}^{2\mathrm{\Delta }_N}}}+\mathrm{}.`$
We see that this term is obtained by plunging the operators $`O^J`$ and $`O^N`$ into the operators $`O^I`$ and $`O^M`$ respectively. The perturbative non-renormalization of the Feynman diagrams of such a type was checked in to first order in perturbation theory where it was noted that this effectively is equivalent to the proof of the non-renormalization of 2-point functions of CPOs given in .
The next integral to be considered is
$$I_3=𝑑\stackrel{}{x}\epsilon ^{d+3}C_{IJL}^{(2)}_0s_I^{(0)}_0s_J^{(0)}_0s_L^{(1)}.$$
(3.32)
Taking into account that $`s_L^{(1)}`$ solves the equation (3.3), we obtain the formula
$$s_L^{(1)}(x_0,\stackrel{}{x})=3\lambda _{LMN}d^{d+1}y\sqrt{g}G_L^\epsilon (x,y)s_M^{(0)}(y)s_N^{(0)}(y),$$
where the Green function can be found in , and satisfies
$`{\displaystyle \frac{}{x_0}}G_L^\epsilon (x,y)|_{x_0=\epsilon }=\epsilon ^{d1}{\displaystyle \frac{d\stackrel{}{k}}{(2\pi )^d}e^{i\stackrel{}{k}(\stackrel{}{x}\stackrel{}{y})}K_L(y_0,\stackrel{}{k})}=\epsilon ^{\mathrm{\Delta }_L1}K_{\mathrm{\Delta }_L}(y_0,\stackrel{}{x}\stackrel{}{y}),`$
where $`K_{\mathrm{\Delta }_L}(x_0,\stackrel{}{x}\stackrel{}{y})`$ is the bulk-to-boundary propagator defined in .
It is convenient to analyse (3.32) in the x-space, where the solution $`s_I^{(0)}`$ can be written as
$`s_I^{(0)}(x_0,\stackrel{}{x})={\displaystyle \frac{1}{\epsilon ^{d\mathrm{\Delta }_I}}}\left({\displaystyle 𝑑\stackrel{}{y}K_{\mathrm{\Delta }_I}(x_0,\stackrel{}{x}\stackrel{}{y})s(\stackrel{}{y})}+o(\epsilon )\right).`$ (3.33)
Thus, for $`s_L^{(1)}`$ we have
$`_0s_L^{(1)}(\epsilon ,\stackrel{}{x})=3\lambda _{LMN}\epsilon ^{\mathrm{\Delta }_L+\mathrm{\Delta }_M+\mathrm{\Delta }_N2d1}\times `$ (3.34)
$`{\displaystyle 𝑑\stackrel{}{y}_3𝑑\stackrel{}{y}_4s_M(\stackrel{}{y}_3)s_N(\stackrel{}{y}_4)\left(d^{d+1}y\sqrt{g}K_{\mathrm{\Delta }_L}(y_0,\stackrel{}{x}\stackrel{}{y})K_{\mathrm{\Delta }_M}(y_0,\stackrel{}{y}\stackrel{}{y}_3)K_{\mathrm{\Delta }_N}(y_0,\stackrel{}{y}\stackrel{}{y}_4)+o(\epsilon )\right)}.`$
Since the cubic couplings $`\lambda _{LMN}`$ vanish in the extremal case, the integral
$$_\epsilon ^{\mathrm{}}𝑑y_0𝑑\stackrel{}{y}\sqrt{g}K_{\mathrm{\Delta }_L}(y_0,\stackrel{}{x}\stackrel{}{y})K_{\mathrm{\Delta }_M}(y_0,\stackrel{}{y}\stackrel{}{y}_3)K_{\mathrm{\Delta }_N}(y_0,\stackrel{}{y}\stackrel{}{y}_4),$$
which appears in evaluation of a 3-point function, is finite in the limit $`\epsilon 0`$ (it diverges only in the extremal case) and, therefore, can be approximated as
$$\mathrm{\Lambda }_{LMN}(\stackrel{}{x};\stackrel{}{y}_3,\stackrel{}{y}_4)+o(\epsilon ),$$
where $`\mathrm{\Lambda }_{LMN}(\stackrel{}{x},\stackrel{}{y}_3,\stackrel{}{y}_4)`$ is defined as
$$\mathrm{\Lambda }_{LMN}(\stackrel{}{x};\stackrel{}{y}_3,\stackrel{}{y}_4)=_0^{\mathrm{}}𝑑y_0𝑑\stackrel{}{y}\sqrt{g}K_{\mathrm{\Delta }_L}(y_0,\stackrel{}{x}\stackrel{}{y})K_{\mathrm{\Delta }_M}(y_0,\stackrel{}{y}\stackrel{}{y}_3)K_{\mathrm{\Delta }_N}(y_0,\stackrel{}{y}\stackrel{}{y}_4).$$
Thus, $`_0s_L^{(1)}`$ behaves itself as
$`_0s_L^{(1)}(\epsilon ,\stackrel{}{x})=3\lambda _{LMN}\epsilon ^{\mathrm{\Delta }_L+\mathrm{\Delta }_M+\mathrm{\Delta }_N2d1}\left(A_{LMN}(\stackrel{}{x})+o(\epsilon )\right)`$ (3.35)
with $`A_{LMN}(\stackrel{}{x})=𝑑\stackrel{}{y}_3𝑑\stackrel{}{y}_4s(\stackrel{}{y}_3)s(\stackrel{}{y}_4)\mathrm{\Lambda }_{LMN}(\stackrel{}{x};\stackrel{}{y}_3,\stackrel{}{y}_4)`$ and we, therefore, find
$$I_3=3C_{IJL}^{(2)}\lambda _{LMN}\epsilon ^{\mathrm{\Delta }_L+\mathrm{\Delta }_M+\mathrm{\Delta }_N3d}𝑑\stackrel{}{x}\epsilon _0s_I^{(0)}\epsilon _0s_J^{(0)}(A_{LMN}(\stackrel{}{x})+o(\epsilon )).$$
The last formula allows one to determine the behavior of the corresponding correlation function in the momentum space. Namely, the leading in $`\epsilon `$ contribution to the 4-point correlation function is proportional to
$`C_{IJL}^{(2)}\lambda _{LMN}\epsilon ^{\mathrm{\Delta }_L+\mathrm{\Delta }_M+\mathrm{\Delta }_N3d}\epsilon _0K_I(\epsilon ,\stackrel{}{k}_1)\epsilon _0K_J(\epsilon ,\stackrel{}{k}_2)\mathrm{\Lambda }_{LMN}(\stackrel{}{k}_1+\stackrel{}{k}_2;\stackrel{}{k}_3,\stackrel{}{k}_4),`$ (3.36)
where function $`\mathrm{\Lambda }_{LMN}(\stackrel{}{k};\stackrel{}{k}_3,\stackrel{}{k}_4)`$ stands now for the Fourier transform of $`\mathrm{\Lambda }_{LMN}(\stackrel{}{x};\stackrel{}{y}_3,\stackrel{}{y}_4)`$ in all its arguments.
To find the relevant contribution to the 4-point function we use (3.26) so that the leading contribution (3.36) is given by the sum of two different terms: the first one contains only one log, while the second one contains the product of two logs. We first consider the one-log case and show that it provides in particular a contribution to a subextremal 4-point correlation function. For definiteness we pick up here the following term
$`C_{IJL}^{(2)}\lambda _{LMN}\epsilon ^{\mathrm{\Delta }_L+\mathrm{\Delta }_M+\mathrm{\Delta }_N+2\mathrm{\Delta }_J4d}k_2^{2\mathrm{\Delta }_Jd}\mathrm{log}k_2\mathrm{\Lambda }_{LMN}(\stackrel{}{k}_1+\stackrel{}{k}_2;\stackrel{}{k}_3,\stackrel{}{k}_4),`$ (3.37)
which gives a nonvanishing contribution if
$$2\mathrm{\Delta }_J+\mathrm{\Delta }_L+\mathrm{\Delta }_M+\mathrm{\Delta }_N4d=\mathrm{\Delta }_I+\mathrm{\Delta }_J+\mathrm{\Delta }_M+\mathrm{\Delta }_N4d$$
i.e. $`\mathrm{\Delta }_I=\mathrm{\Delta }_J+\mathrm{\Delta }_L`$. Clearly, this equality is not too restrictive and it allows in particular the solution for the subextremal case<sup>5</sup><sup>5</sup>5The extremal case is of no interest here since the coupling $`\lambda _{LMN}`$ vanishes., i.e., when $`\mathrm{\Delta }_I=\mathrm{\Delta }_J+\mathrm{\Delta }_M+\mathrm{\Delta }_N2`$ and $`\mathrm{\Delta }_L=\mathrm{\Delta }_M+\mathrm{\Delta }_N2`$. Let us now show that for these values of conformal dimensions (3.37) indeed represents the relevant momentum space structure of a subextremal 4-point correlation function.
Due to the non-renormalization theorem a subextremal 4-point correlation function of single-trace CPOs is given by the sum of products of two-point functions that is further restricted by the conformal invariance to the form
$`O^I(\stackrel{}{x}_1)O^J(\stackrel{}{x}_2)O^M(\stackrel{}{x}_3)O^N(\stackrel{}{x}_4)={\displaystyle \frac{A_{IJMN}}{x_{12}^{2(\mathrm{\Delta }_J+\alpha 1)}x_{13}^{2(\mathrm{\Delta }_M+\beta 1)}x_{14}^{2(\mathrm{\Delta }_N+\gamma 1)}x_{23}^{2\gamma }x_{24}^{2\beta }x_{34}^{2\alpha }}},`$ (3.38)
where $`\mathrm{\Delta }_I=\mathrm{\Delta }_J+\mathrm{\Delta }_M+\mathrm{\Delta }_N2`$ and $`\alpha ,\beta ,\gamma `$ are integers obeying the condition $`\alpha +\beta +\gamma =1`$, so that only one of them is non-zero and equals to 1. Thus we have three different subextremal structures which have in general three different coefficients $`A`$.
Consider the structure with $`\beta =\gamma =0`$ and perform the Fourier transform. The corresponding structure in the momentum space looks as
$`O^I(\stackrel{}{k}_1)O^J(\stackrel{}{k}_2)O^M(\stackrel{}{k}_3)O^N(\stackrel{}{k}_4){\displaystyle 𝑑\stackrel{}{x}_1𝑑\stackrel{}{x}_2𝑑\stackrel{}{x}_3𝑑\stackrel{}{x}_4\frac{e^{i\stackrel{}{k}_1\stackrel{}{x}_1+i\stackrel{}{k}_2\stackrel{}{x}_2+i\stackrel{}{k}_3\stackrel{}{x}_3+i\stackrel{}{k}_4\stackrel{}{x}_4}}{x_{12}^{2\mathrm{\Delta }_J}x_{13}^{2(\mathrm{\Delta }_M1)}x_{14}^{2(\mathrm{\Delta }_N1)}x_{34}^2}}`$
$`\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4)k_2^{2\mathrm{\Delta }_Jd}\mathrm{log}k_2{\displaystyle 𝑑v𝑑w\frac{e^{i\stackrel{}{k}_3vi\stackrel{}{k}_4w}}{v^{2(\mathrm{\Delta }_M1)}w^{2(\mathrm{\Delta }_N1)}(vw)^2}},`$
where new integration variables $`v=x_{13}`$ and $`w=x_{14}`$ were introduced. Coming back to (3.37) it remains to note that in the x-space $`\mathrm{\Lambda }_{LMN}(\stackrel{}{x};\stackrel{}{y}_3,\stackrel{}{y}_4)`$ is fixed by conformal invariance to be
$`\mathrm{\Lambda }_{LMN}(\stackrel{}{x};\stackrel{}{y}_3,\stackrel{}{y}_4)={\displaystyle \frac{C}{(\stackrel{}{x}\stackrel{}{y}_3)^{2(\mathrm{\Delta }_M1)}(\stackrel{}{x}\stackrel{}{y}_4)^{2(\mathrm{\Delta }_N1)}(\stackrel{}{y}_3\stackrel{}{y}_4)^2}},`$
where we have used subextremality condition $`\mathrm{\Delta }_L=\mathrm{\Delta }_M+\mathrm{\Delta }_N2`$ and $`C`$ is the numerical (non-zero) constant. Transforming this expression to the momentum space we therefore find
$`\mathrm{\Lambda }_{LMN}(\stackrel{}{k}_1+\stackrel{}{k}_2;\stackrel{}{k}_3,\stackrel{}{k}_4)=\delta (\stackrel{}{k}_1+\stackrel{}{k}_2+\stackrel{}{k}_3+\stackrel{}{k}_4){\displaystyle 𝑑v𝑑w\frac{e^{i\stackrel{}{k}_3vi\stackrel{}{k}_4w}}{v^{2(\mathrm{\Delta }_M1)}w^{2(\mathrm{\Delta }_N1)}(vw)^2}}.`$
Thus, we have shown that the field redefinition induces a non-trivial contribution to the subextremal 4-point functions.
The general case $`\mathrm{\Delta }_I=\mathrm{\Delta }_J+\mathrm{\Delta }_L`$ is considered in the same way, and we get that the changing term in a 4-point function has the form
$`O^I(\stackrel{}{x}_1)O^J(\stackrel{}{x}_2)O^M(\stackrel{}{x}_3)O^N(\stackrel{}{x}_4){\displaystyle \frac{1}{x_{12}^{2\mathrm{\Delta }_J}x_{13}^{\mathrm{\Delta }_L+\mathrm{\Delta }_M\mathrm{\Delta }_N}x_{14}^{\mathrm{\Delta }_L+\mathrm{\Delta }_N\mathrm{\Delta }_M}x_{34}^{\mathrm{\Delta }_M+\mathrm{\Delta }_N\mathrm{\Delta }_L}}}+\mathrm{}.`$
We see that this term is obtained by plunging the operator $`O^J`$ into the operator $`O^I`$. The perturbative non-renormalization of the Feynman diagrams of such a type seems to be equivalent to the proof of the non-renormalization of 2- and 3-point functions of CPOs given in .
Consideration of the term involving two logs shows that it scales too fast and by this reason does not lead to any contribution to 4-point functions.
The last integral to be considered is
$$I_4=𝑑\stackrel{}{x}\epsilon ^{d+1}C_{IJKL}^{(4)}_0s_I^{(0)}_as_J^{(0)}^bs_K^{(0)}^a_bs_L^{(0)}$$
(3.39)
The only case when the integral gives a contribution to a 4-point function is $`a=b=0`$. However, in this case we can use the equations of motion for scalars $`s^I`$ to express $`^0_0s_L^{(0)}`$ as $`(m_L^2^i_i)s_L^{(0)}`$. Thus, this integral is equivalent to the integral $`I_1`$ (3.23), and can give a nonvanishing contribution to an extremal 4-point function.
This completes our consideration of the boundary terms (3.22).
## 4 Conclusion
In this paper we studied nonlinear derivative-dependent transformations of gravity fields, and showed that they change 3- and 4-point functions in a boundary CFT. We interpreted such a change of correlation functions as a manifestation of an operator basis transformation in CFT. Thus, a derivative-dependent field redefinition invokes a transformation of operators in CFT, and, as the consequence, a transformation of correlation functions. However, this transformation has a very restrictive form as by a derivative-dependent field redefinition it is possible to change only the coefficients of non-renormalized structures of correlation functions. In particular, one probably can find such a field redefinition of the gravity fields that the redefined scalars $`s^I`$ would be dual to the single-trace CPOs. Still the analysis performed in the paper does not allow one to conclude this definitely. The point is that we do not have enough parameters in the field transformations because we only considered field redefinitions of the scalars $`s^I`$. In general, one should take into account the scalar dependent redefinitions of vector and tensor fields as well. This would give us enough parameters to transform the correlation functions of the extended CPOs into the ones of the single-trace CPOs.
## 5 Appendix
In the quartic action for scalars $`s^I`$ was found in the form
$`S(s)={\displaystyle _{AdS_5}}\left(_4^{(4)}+_4^{(2)}+_4^{(0)}\right),`$
where the quartic terms contain the 4-derivative couplings
$`_4^{(4)}=\left(S_{I_1I_2I_3I_4}^{(4)}+A_{I_1I_2I_3I_4}^{(4)}\right)s^{I_1}_as^{I_2}_b^2(s^{I_3}^as^{I_4}),`$
the 2-derivative couplings
$`_4^{(2)}=\left(S_{I_1I_2I_3I_4}^{(2)}+A_{I_1I_2I_3I_4}^{(2)}\right)s^{I_1}_as^{I_2}s^{I_3}^as^{I_4}`$
and the couplings without derivatives
$`_4^{(0)}=S_{I_1I_2I_3I_4}^{(0)}s^{I_1}s^{I_2}s^{I_3}s^{I_4}.`$
The corresponding vertices have the following symmetry properties
$`S_{I_1I_2I_3I_4}^{(4)}=S_{I_2I_1I_3I_4}^{(4)}=S_{I_3I_4I_1I_2}^{(4)},A_{I_1I_2I_3I_4}^{(4)}=A_{I_2I_1I_3I_4}^{(4)}=A_{I_3I_4I_1I_2}^{(4)},`$
$`S_{I_1I_2I_3I_4}^{(2)}=S_{I_2I_1I_3I_4}^{(2)}=S_{I_3I_4I_1I_2}^{(2)},A_{I_1I_2I_3I_4}^{(2)}=A_{I_2I_1I_3I_4}^{(2)}=A_{I_3I_4I_1I_2}^{(2)}`$
and their explicit values are given in . What is important for our discussion here is that all the couplings are represented as sums of the $`SO(6)`$ tensors of three different types:
$`F(I_5)a_{I_1I_2I_5}a_{I_3I_4I_5},F(I_5)t_{I_1I_2I_5}t_{I_3I_4I_5},F(I_5)p_{I_1I_2I_5}p_{I_3I_4I_5},`$
where $`F(I_5)`$ is a function of $`I_5`$ and the sum over $`I_5`$ is assumed. There also appear tensors obtained from these ones by different permutation of indices. Recall that $`a_{I_1I_2I_3}`$, $`t_{I_1I_2I_3}`$ and $`p_{I_1I_2I_3}`$ represent the following integrals involving the scalar $`Y^I`$, the vector $`Y_\alpha ^I`$ and the tensor $`Y_{(\alpha \beta )}^I`$ spherical harmonics respectively
$`a_{I_1I_2I_3}={\displaystyle _{S^5}}Y^{I_1}Y^{I_2}Y^{I_3},t_{I_1I_2I_3}={\displaystyle _{S^5}}^\alpha Y^{I_1}Y^{I_2}Y_\alpha ^{I_3},p_{I_1I_2I_3}={\displaystyle _{S^5}}^\alpha Y^{I_1}^\beta Y^{I_2}Y_{(\alpha \beta )}^{I_3}.`$
To prove the vanishing of the couplings in the subextremal case as well as the vanishing of the 4-derivative couplings in the sub-subextremal case we find convenient to pass to the 4-derivative vertices of the form (2.7). This is achieved by using the following relations valid on-shell:
$`A_{1234}^{(4)}{\displaystyle s_1_as_2_b^2(s_3^as_4)}`$ $`=`$ $`2A_{1234}^{(4)}{\displaystyle _as_1_bs_2^as_3^bs_4}`$
$``$ $`4A_{1234}^{(4)}{\displaystyle s_1_as_2s_3^as_4}`$
$``$ $`{\displaystyle \frac{1}{4}}A_{1234}^{(4)}(m_1^2m_2^2)(m_3^2m_4^2){\displaystyle s_1s_2s_3s_4}.`$
$`S_{1234}^{(4)}{\displaystyle s_1_as_2_b^2(s_3^as_4)}`$ $`=`$ $`S_{1234}^{(4)}{\displaystyle _as_1^as_2_bs_3^bs_4}`$ (5.1)
$`+`$ $`S_{1234}^{(4)}\left(m_1^2+m_2^2+m_3^2+m_4^24\right){\displaystyle s_1_as_2s_3^as_4}`$
$`+`$ $`{\displaystyle \frac{1}{4}}S_{1234}^{(4)}(m_1^2+m_2^2)(m_3^2+m_4^2){\displaystyle s_1s_2s_3s_4},`$
where here and below we write concisely the summation index $`I_1`$ simply as 1 and similar for the others, $`m`$ denotes the $`AdS`$ mass of a scalar field.
First we consider the subextremal case and assume for definiteness that $`k_1=k_2+k_3+k_42`$. It is easy to show, by using the description of spherical harmonics as restrictions of functions, vectors and tensors on the $`𝐑^6`$ in which the sphere $`S^5`$ is embedded , that the tensor<sup>6</sup><sup>6</sup>6We do not assume here summation over $`I_5`$. $`t_{125}t_{345}`$ does not vanish in the subextremal case only for $`k_5=k_3+k_41`$ while for $`p_{125}p_{345}`$ it is the case only if $`k_5=k_3+k_42`$. As for the tensor $`a_{125}a_{345}`$, it differs from zero in two cases: when $`k_5=k_3+k_4`$ or when $`k_5=k_3+k_42`$. Analogously, the only non-trivial values of $`k_5`$ for $`a_{135}a_{245}`$ are $`k_5=k_2+k_4`$ and $`k_5=k_2+k_42`$, and for $`a_{145}a_{235}`$ they are $`k_5=k_2+k_3`$ and $`k_5=k_2+k_32`$. Thus in all vertices we can replace $`k_5`$ by a corresponding function of $`k_2,k_3,k_4`$, and, then the only dependence on $`k_5`$ is in tensors $`t_{125}t_{345}`$, $`p_{125}p_{345}`$, $`a_{125}a_{345}`$, $`a_{135}a_{245}`$ and $`a_{145}a_{235}`$. However, not all of these tensors are independent. Indeed, $`a_{125}a_{345}`$, $`a_{135}a_{245}`$ and $`a_{145}a_{235}`$ are subjected to the following three identities :
$`a_{125}a_{345}=a_{135}a_{245}=a_{145}a_{235},`$ (5.2)
$`f_5(a_{125}a_{345}+a_{135}a_{245}+a_{235}a_{145})=(f_1+f_2+f_3+f_4)a_{125}a_{345},`$
where $`f_i=f(k_i)=k_i(k_i+4)`$. For the sake of simplicity it is useful to introduce the notation
$`\begin{array}{cc}l_1=a_{125}a_{345}|_{k_5=k_3+k_4},\hfill & l_2=a_{125}a_{345}|_{k_5=k_3+k_42},\hfill \\ m_1=a_{145}a_{235}|_{k_5=k_2+k_3},\hfill & m_2=a_{145}a_{235}|_{k_5=k_2+k_32},\hfill \\ n_1=a_{135}a_{245}|_{k_5=k_2+k_4},\hfill & n_2=a_{135}a_{245}|_{k_5=k_2+k_42},\hfill \end{array}`$ (5.6)
where, e.g., $`l_1`$ denotes tensor $`a_{125}a_{345}`$ for the value of $`k_5`$ equal to $`k_3+k_4`$. Hence we have six tensors corresponding to different values of $`k_5`$ and to different order of indices, which are confined by three relations (5.2). Therefore, restricting eqs.(5.2) to the subextremal case, i.e., putting $`k_1=k_2+k_3+k_42`$ one can solve them for any three tensors. If we choose here $`l_1`$, $`m_1`$ and $`n_1`$ as independent variables, then $`l_2`$, $`m_2`$ and $`n_2`$ are expressed as
$`l_2`$ $`=`$ $`{\displaystyle \frac{(m_1+n_1l_1)(k_2+1)+m_1k_3+n_1k_4}{k_1+k_2+k_3+2}},`$
$`m_2`$ $`=`$ $`{\displaystyle \frac{(n_1+l_1m_1)(k_4+1)+l_1k_3+n_1k_2}{k_1+k_2+k_3+2}},`$ (5.7)
$`n_2`$ $`=`$ $`{\displaystyle \frac{(m_1n_1+l_1)(k_3+1)+l_1k_4+m_1k_2}{k_1+k_2+k_3+2}}.`$
For $`t_{125}t_{345}`$ and $`p_{125}p_{345}`$ we will need the following three identities found in :
$`t_{125}t_{345}={\displaystyle \frac{(f_1f_2)(f_3f_4)}{4f_5}}a_{125}a_{345}+{\displaystyle \frac{1}{4}}f_5(a_{145}a_{235}a_{245}a_{135}),`$ (5.8)
$`(1f_5)t_{125}t_{345}={\displaystyle \frac{1}{4}}(f_5^2f_5(f_1+f_2+f_3+f_44))(a_{145}a_{235}a_{135}a_{245})`$ (5.9)
$`{\displaystyle \frac{4f_5}{4f_5}}(f_1f_2)(f_3f_4)a_{125}a_{345},`$
$`p_{125}p_{345}={\displaystyle \frac{(f_1f_2)(f_3f_4)}{2(f_55)}}t_{125}t_{345}{\displaystyle \frac{5}{4f_5(f_55)}}d_{125}d_{345}`$
$`{\displaystyle \frac{1}{20}}(f_1+f_2f_5)(f_3+f_4f_5)a_{125}a_{345}+{\displaystyle \frac{1}{8}}(f_1+f_3f_5)(f_2+f_4f_5)a_{135}a_{245}`$
$`+{\displaystyle \frac{1}{8}}(f_1+f_4f_5)(f_2+f_3f_5)a_{145}a_{235},`$ (5.10)
where
$$d_{123}=_{S^5}^{(\alpha }^{\beta )}Y^{I_3}_\alpha Y^{I_1}_\beta Y^{I_2}=(\frac{1}{10}f_2f_3+\frac{1}{10}f_1f_3+\frac{1}{2}f_1f_2\frac{1}{4}f_1^2\frac{1}{4}f_2^2+\frac{3}{20}f_3^2)a_{125}.$$
Since in the subextremal case $`t_{125}t_{345}`$ is non-zero only for one value of $`k_5`$ we may use formula (5.8) and eqs.(5.7) to express $`t_{125}t_{345}`$ in terms of $`l_1`$, $`m_1`$ and $`n_1`$. Similarly, combining eq.(5.10) with (5.8) and with eqs.(5.7) one obtains an analogous representation for $`p_{125}p_{345}`$. In this way we have expressed all the quartic vertices via independent tensors $`l_1`$, $`m_1`$ and $`n_1`$.
Now we single out the field $`s^{I_1}`$ and write the relevant part of the quartic 4-derivative vertices as functions of $`l_1`$, $`m_1`$ and $`n_1`$ in the form
$`Ł^{(4)}=4{\displaystyle \underset{I_2,I_3,I_4}{}}\left(S_{I_1I_2I_3I_4}^{(4)}A_{I_1I_3I_2I_4}^{(4)}+A_{I_1I_4I_3I_2}^{(4)}\right)_as^{I_1}^as^{I_2}_bs^{I_3}^bs^{I_4},`$ (5.11)
where we sum over the representations satisfying the subextremality condition. Now, we substitute the values of $`k_5`$ discussed above, and $`k_1=k_2+k_3+k_42`$ in the quartic couplings, and obtain zero.
To analyse 2-derivative terms we represent the 2-derivative Lagrangian as follows
$`Ł^{(2)}`$ $`=`$ $`4{\displaystyle \underset{I_2,I_3,I_4}{}}(({\displaystyle \frac{1}{2}}\stackrel{~}{S}_{I_1I_3I_2I_4}^{(2)}+\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)})s^{I_1}^as^{I_2}s^{I_3}_as^{I_4}`$
$`+`$ $`{\displaystyle \frac{1}{4}}(\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)}(m_4^2m_3^2)\stackrel{~}{S}_{I_1I_2I_3I_4}^{(2)}(m_4^2+m_3^2))s^{I_1}s^{I_2}s^{I_3}s^{I_4}),`$
where using (5.1) we define
$`\stackrel{~}{S}_{I_1I_2I_3I_4}^{(2)}`$ $`=`$ $`S_{I_1I_2I_3I_4}^{(2)}+S_{I_1I_2I_3I_4}^{(4)}\left(m_1^2+m_2^2+m_3^2+m_4^24\right),`$
$`\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)}`$ $`=`$ $`A_{I_1I_2I_3I_4}^{(2)}4A_{I_1I_2I_3I_4}^{(4)}.`$
This time substituting $`k_5`$ and $`k_1`$ and symmetrizing the expression obtained in $`I_2`$ and $`I_4`$, we get a non-zero function which is, however, completely symmetric in $`I_2`$, $`I_3`$ and $`I_4`$. Thus we can remove the 2-derivative term by using the shift
$$s^{I_1}s^{I_1}\frac{2}{3\kappa _1}\left(\frac{1}{2}\stackrel{~}{S}_{I_1I_3I_2I_4}^{(2)}+\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)}\right)s^{I_2}s^{I_3}s^{I_4},$$
where $`\kappa _1=\frac{32k_1(k_11)(k_1+2)}{k_1+1}`$. This shift also produces an additional contribution to the non-derivative terms which is equal to
$$\frac{2}{3}\left(\frac{1}{2}\stackrel{~}{S}_{I_1I_3I_2I_4}^{(2)}+\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)}\right)(m_2^2+m_3^2+m_4^2m_1^2)s^{I_1}s^{I_2}s^{I_3}s^{I_4}.$$
After accounting this contribution the non-derivative terms acquire the form
$`Ł^{(0)}`$ $`=`$ $`4{\displaystyle \underset{I_2,I_3,I_4}{}}(S_{I_1I_2I_3I_4}^{(0)}{\displaystyle \frac{1}{6}}({\displaystyle \frac{1}{2}}\stackrel{~}{S}_{I_1I_3I_2I_4}^{(2)}+\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)})(m_2^2+m_3^2+m_4^2m_1^2)`$
$`+`$ $`{\displaystyle \frac{1}{4}}\left(\stackrel{~}{A}_{I_1I_2I_3I_4}^{(2)}(m_4^2m_3^2)\stackrel{~}{S}_{I_1I_2I_3I_4}^{(2)}(m_4^2+m_3^2)\right)`$
$`+`$ $`{\displaystyle \frac{1}{4}}S_{1234}^{(4)}(m_1^2+m_2^2)(m_3^2+m_4^2){\displaystyle \frac{1}{4}}A_{1234}^{(4)}(m_1^2m_2^2)(m_3^2m_4^2))s^{I_1}s^{I_2}s^{I_3}s^{I_4}.`$
Substituting $`k_5`$ and $`k_1`$ and symmetrizing the coefficient obtained in $`I_2`$, $`I_3`$ and $`I_4`$ we end up with zero. Thus, we have shown that after the additional field redefinition all subextremal quartic couplings vanish.
The treatment of the 4-derivative quartic couplings in the sub-subextremal case is quite analogous to the previous one. For definiteness we assume that $`k_1=k_2+k_3+k_44`$. Then $`a_{125}a_{345}`$ is non-zero in three cases: $`k_5=k_3+k_4`$, $`k_5=k_3+k_42`$ and $`k_5=k_3+k_44`$. Similarly $`a_{135}a_{245}`$ is non-zero only for $`k_5`$ equal to $`k_2+k_4`$, $`k_2+k_42`$ or $`k_2+k_44`$, while $`a_{145}a_{235}`$ admits for $`k_5`$ one of the following values $`k_2+k_3`$, $`k_2+k_32`$ or $`k_2+k_34`$. Denote
$`\begin{array}{ccc}l_1=a_{125}a_{345}|_{k_5=k_3+k_4},\hfill & l_2=a_{125}a_{345}|_{k_5=k_3+k_42},\hfill & l_3=a_{125}a_{345}|_{k_5=k_3+k_44},\hfill \\ m_1=a_{145}a_{235}|_{k_5=k_2+k_3},\hfill & m_2=a_{145}a_{235}|_{k_5=k_2+k_32},\hfill & m_3=a_{145}a_{235}|_{k_5=k_2+k_34},\hfill \\ n_1=a_{135}a_{245}|_{k_5=k_2+k_4},\hfill & n_2=a_{135}a_{245}|_{k_5=k_2+k_42},\hfill & n_3=a_{135}a_{245}|_{k_5=k_2+k_44}.\hfill \end{array}`$ (5.16)
Then identities (5.2) allow one to express $`l_3,m_3,n_3`$ via six independent tensors $`l_1,m_1,n_1`$ and $`l_2,m_2,n_2`$, e.g.,
$`l_3`$ $`=`$ $`{\displaystyle \frac{1}{k_2+k_3+k_4+2}}(m_2+n_2+l_2+(2m_1m_2+l_2)k_3`$
$`+`$ $`(2m_1m_22n_1n_2+2l_1+2l_2)k_2+(2n_1n_2+l_2)k_4).`$
The formulas for $`m_3`$ and $`n_3`$ are obtained from this one by permutations of indices.
Except the tensor structures we have just considered the quartic couplings of 4-derivative vertices contain a tensor $`(f_51)^2t_{125}t_{345}`$ (see Appendix A of ). In the sub-subextremal case $`t_{125}t_{345}`$ differs from zero only for $`k_5=k_3+k_41`$ or $`k_5=k_3+k_43`$. It is then suitable to represent
$`(f_51)^2t_{125}t_{345}=\left((f_5\alpha )(f_5\beta )+a(f_51)+b\right)t_{125}t_{345},`$ (5.18)
where
$`a`$ $`=`$ $`\alpha +\beta 2,`$
$`b`$ $`=`$ $`(\alpha 1)(\beta 1)`$
and pick up for $`\alpha `$ and $`\beta `$ the following values $`\alpha =f(k_3+k_41)`$ and $`\beta =f(k_3+k_43)`$. Clearly, in the sub-subextremal case the first term in the r.h.s. of (5.18) is absent and we may use identities (5.8) and (5.9) to rewrite $`(f_51)^2t_{125}t_{345}`$ via $`l_1,\mathrm{},n_2`$. Hence, as in the subextremal case, we reduced all quartic couplings of 4-derivative vertices to the independent tensor structures.
Now upon substituting in eq.(5.11) the 4-derivative quartic couplings evaluated for the proper values of $`k_5`$ and putting $`k_1=k_2+k_3+k_44`$ we obtain zero.
ACKNOWLEDGMENT
We would like to thank A. Tseytlin for prompt reading of the manuscript and useful remarks, and D. Freedman for useful correspondence, G.A. is grateful to S. Theisen and S. Kuzenko, and S.F. is grateful to A. Tseytlin and S. Mathur for valuable discussions. The work of G.A. was supported by the Alexander von Humboldt Foundation and in part by the RFBI grant N99-01-00166, and the work of S.F. was supported by the U.S. Department of Energy under grant No. DE-FG02-96ER40967 and in part by RFBI grant N99-01-00190. |
warning/0003/cond-mat0003115.html | ar5iv | text | # Tunnel splittings for one dimensional potential wells revisited
## I Introduction
The purpose of this article is to reexamine some formal aspects of the problem of calculating the quantum mechanical tunnel splittings in a smooth, symmetric, one-dimensional double well potential, such as that in Fig. 1. Readers may justifiably wonder if anything new remains to be said on so mature a subject, and we assure them that, by and large, there isn’t. The physical phenomenon is certainly very well understood, as are the basic mathematical ideas behind the calculations. We find, nevertheless, that especially in the case of the ground state splitting, there is some confusion regarding the correct WKB answer for this splitting, which is often taken in a form that is less accurate than it needs to be. In particular, we note that the formula for the splitting in the masterly text book by Landau and Lifshitz leads to a prefactor which is incorrect for the ground state . A second problem is that the form in which this result is usually presented is poorly suited to calculation of the ground state splitting. In applying it to model problems such as the quartic double well potential, e.g., the unwary student is unnecessarlily led into asymptotic expansions of elliptic integrals, which must then be taken from standard tables of formulas. Because of this confusion, there is also confusion regarding the equivalence of the WKB and the instanton methods for calculating tunnel splittings.
Like many other problems in physics of a similar nature, the above discrepancies and their resolution are at the same time “well-known” and not known at all! They are well-known to some experts. Thus, the correct answers for the splitting are implicit throughout early WKB studies of the anharmonic oscillator , and can be ferreted out with some work. They are also clearly known to many field theorists , and to authors of more pedagogical articles . At the same time, the confusion persists, and continues to crop up from time to time. Thus it seems worthwhile to address this issue here.
It should also be stated at the outset that our article is of no interest if one only wants the splitting to “exponential accuracy”, i.e., if one only wants the Gamow factor. In most physical applications, this is all that can sensibly be done, and the prefactor is best estimated as an “attempt frequency.” From a mathematical point of view, however, the tunnel splitting is found as an asymptotic approximation in the limit $`\mathrm{}0`$. An “exponentially accurate” answer for the splitting $`\mathrm{\Delta }`$ is one for which one only has an asymptotically correct result for $`\mathrm{ln}\mathrm{\Delta }`$. An asymptotically correct answer for $`\mathrm{\Delta }`$ itself requires worrying about the prefactor. More importantly, any formalism which did not give this prefactor, or which was incapable of giving it correctly even in principle, would be regarded as logically unsound. It is from this perspective that the prefactor is important. The confusion that is referred to above can be said to pertain solely to this prefactor, and readers who are not concerned with such arcana should stop reading here. On the other hand, the WKB analysis presented here is easily incorporated into a graduate level discussion of tunnel splittings, with little additional effort beyond that of the standard treatment, and there is no reason not to do so. The problem may also serve as a “real-life” physics example of some nontrivial asymptotic analysis. (To avoid misunderstanding, let us also state at the outset what we mean by WKB theory. We use the term in the broad sense used by Bender and Orszag , or by Berry and Mount , i.e., as a body of mathematical techniques that yield global understanding and systematic asymptotic approximations to solutions of many linear differential equations,
including Schrödinger’s equation, and shares ideas with other methods such as asymptotic matching and patching, boundary layer theory, etc. Some readers may, however, define WKB to encomapss only the approximate exponential form (11) (and its oscillatory counterpart) and connection formulas at linear turning points. A subset of these readers may recognize that our treatment in Sec. II is tantamount to the use of quadratic turning point connection formulas (for which see Berry and Mount ), and object that this goes beyond WKB. Such a semantic restriction of the scope of the term ‘WKB’, should in our view, be avoided. Since the method did not originate with Wentzel, Kramers, and Brillouin, but had antecedents in the work of Jeffreys, Rayleigh, Carlini, Green, Liouville, and perhaps others, there is no compelling historical reason for this restriction (as opposed to say, the usages ‘Einstein model’ for lattice specific heat, or the ‘Kronig-Penney model’ for electrons in a periodic potential), and it is surely more useful to present the subject to students as one which allows for systematic improvement, and is not a closed subject of study even today.)
It is with this motivation that we revisit the problem of calculating tunnel splittings in symmetric double well potentials. The article has three distinct aims. The first is to carefully calculate the prefactor multiplying the exponential of the turning-point-to-turning-point action integral in the standard WKB expression for the splitting $`\mathrm{\Delta }`$. The correct formula is Eq. (1) below. Our second aim is to present another formula \[See Eqs. (46)\] for the ground state splitting , that reduces everything to the evaluation of two integrals involving the potential $`V(x)`$, and does not require looking up any asymptotic expansions. The last aim is to show that the correct WKB and instanton methods do yield the same ground state splitting. We will do this by starting with Coleman’s instanton method expression for the splitting and showing that it reduces to our simple formula.
We expect that our WKB based disucssion will be accessible to students who have seen some graduate level quantum mechanics, even though our starting formula for the splitting is rarely ever mentioned in the common text books. One book where it does appear is again that of Landau and Lifshitz , who give an exceptionally lucid and self-contained discussion. We will comment further on this formula when we come to it in Sec. II, and give a separate derivation of it in Appendix A. The instanton formalism, and the path integral approach to quantum mechanics upon which it is in turn based, are less likely to be familiar, but clear and accessible discussions have been given in this journal by Holstein , and by Holstein and Swift . The former discusses exactly the same problem as us, namely, the tunnel splitting, but for the special case of a quartic double well. We enthusiastically recommend Coleman’s somewhat longer but authoratitive article on instantons to readers who wish to learn more about this technique. A related paper by Carlitz and Nicole may also be read profitably.
A plan of our paper and summary of our results is as follows. We consider a potential which is smooth, reflection symmetric about $`x`$=0 \[$`V(x)=V(x)`$\], and which has quadratic minima at $`x=\pm a`$. The correct WKB answer for the ground state splitting is
$$\mathrm{\Delta }=\frac{\mathrm{}\omega }{\sqrt{e\pi }}\mathrm{exp}\left[_a^{}^a^{}\frac{|p|}{\mathrm{}}𝑑x\right].$$
(1)
Here, $`\omega `$ is the frequency of small amplitude oscillations in the wells about $`x=\pm a`$, $`\pm a^{}`$ are the classical turning points given by the equation
$$V(a^{})=E_0=\frac{1}{2}\mathrm{}\omega +o(\mathrm{}),$$
(2)
and $`p`$ is the momentum
$$p(x)=[2m(V(x)E)]^{1/2}.$$
(3)
Note that $`p(x)`$ is imaginary in the classically forbidden region $`a^{}<x<a^{}`$, and that we have taken $`V(\pm a)=0`$.
The proof of Eq. (1) is given in Sec. II. This entails matching the WKB wavefunction in the classically forbidden region to the exact harmonic oscillator wavefunction near the classical turning point, rather than use the connection formulas. This matching calculation may actually be found in a paper by Furry , but we include it because it is very short, and so as to have the complete argument in one place in a consistent notation.
Prefactor corrections similar to those in Eq. (1) accompany the excited states too, and are discussed in Appendix B.
The expression (1) is not easy to use because the integrand in the exponential is close to a singularity near the limits, which means that the next to leading dependence of $`\mathrm{\Delta }`$ on $`\mathrm{}`$ is not manifest. In fact, the true preexponential factor varies as $`\mathrm{}^{1/2}`$, and this is not obvious from Eq. (1). In Sec. III, we shall, therefore, carefully extract the singular contributions to the action integral from the end points, and show that the splitting may be written for a general potential as
$$\mathrm{\Delta }=2\mathrm{}\omega \left(\frac{m\omega a^2}{\pi \mathrm{}}\right)^{1/2}e^Ae^{S_0/\mathrm{}},$$
(4)
where $`S_0`$ (note the limits of integration) is the action integral
$$S_0=_a^a\left(2mV(x)\right)^{1/2}𝑑x,$$
(5)
and
$$A=_0^a\left[\frac{m\omega }{\sqrt{2mV(x)}}\frac{1}{ax}\right]𝑑x.$$
(6)
This expression does not have the complexities mentioned above. The $`\mathrm{}^{1/2}`$ dependence of the prefactor is apparent, and the answer involves only integral functionals of $`V(x)`$. These formulas will be applied to two model problems in Sec. V.
In Sec. IV, we turn to the instanton approach . Here, the splitting is expressed as
$$\mathrm{\Delta }=2\mathrm{}K\left(\frac{S_0}{2\pi \mathrm{}}\right)^{1/2}\mathrm{exp}(S_0/\mathrm{}),$$
(7)
where $`S_0`$ is as in Eq. (5), and
$$K=\left[\frac{det(_\tau ^2+\omega ^2)}{det^{}[_\tau ^2+m^1V^{\prime \prime }(x_{\mathrm{cl}}(\tau ))]}\right]^{1/2}$$
(8)
is the ratio of fluctuation determinants. Here, $`x_{\mathrm{cl}}(\tau )`$ is the instanton, which obeys the (Euclidean) classical equation of motion
$$m\frac{d^2x_{\mathrm{cl}}(\tau )}{d\tau ^2}V^{}(x_{\mathrm{cl}})=0,$$
(9)
with the boundary conditions $`x_{\mathrm{cl}}(\pm \mathrm{})=\pm a`$, and the additional condition $`x_{\mathrm{cl}}(0)=0`$ to fix the time translation degree of freedom. Further, $`V^{}=dV/dx`$, $`V^{\prime \prime }=d^2V/dx^2`$, and the prime on the det in Eq. (8) means that the zero eigenvalue of the operator argument is to be excluded from a product of all eigenvalues.
A short and pedagogical discussion of the instanton method may be found in Holstein . Holstein does not explain how to calculate the quantity $`K`$ (which is proportional to his quantity $`K_1`$), so a few remarks that help elucidate its nature may not be out of place. Note first, that each determinant appearing in Eq. (8) is of a one-dimensional Schrodinger-like operator, in which $`\tau `$ plays the role of position, and either $`\omega ^2`$ or $`m^1V^{\prime \prime }(x_{\mathrm{cl}}(\tau ))`$ plays the role of the potential energy. Next we note, that the determinant of such an operator may be defined as an infinite product of all its eigenvalues. The spectrum of operators at hand may be rendered completely discrete by adding hard walls at $`\tau =\pm T`$, and letting $`T\mathrm{}`$ at the end. This scheme has the advantage that it allows us to put the eigenvalues of both operators in one-to-one correspondence. This in turn enables us to argue that one eigenvalue is missing in the denominator, and since the eigenvalues have the dimensions of $`_\tau ^2`$ or $`\omega ^2`$, we see that $`K`$ has the dimensions of frequency. For a smooth potential $`V(x)`$, the curvature $`m^1V^{\prime \prime }(x_{\mathrm{cl}}(\tau ))`$ is never very different from $`\omega ^2`$ in the interval $`axa`$, so there is only one frequency scale in the problem. It follows that $`K`$ is of order $`\omega `$.
To the author’s knowledge, the equivalence of Eqs. (1) and (7) has only been shown for particular examples, such as the quartic double well potential. A general demonstration is lacking, although it is implicit in the fact that both methods start from the same point and are correctly executed. A direct demonstration is therefore of some value. Secondly, Eq. (7) is complicated just as Eq. (1) was. Indeed, the ratio of determinants seems more intimidating and harder to calculate than the integral in Eq. (1). Langer’s original calculation for $`K`$ is quite involved, and while it has been greatly simplified by Coleman , it still requires solving the auxilliary problem of finding the instanton $`x_{\mathrm{cl}}(\tau )`$.
To relate the instanton and WKB answers, ws therefore show in Sec. IV, that Coleman’s result for $`K`$ may be expressed in terms of the integral $`A`$ in Eq. (6). Thus, although we do not provide a complete derivation of the instanton result for $`\mathrm{\Delta }`$, once Coleman’s formulas are accepted, our analysis shows the equivalence of the two approaches.
## II WKB formula for ground state splitting
Our starting point is the formula
$$\mathrm{\Delta }=\frac{2\mathrm{}^2}{m}\psi _0(0)\psi _0^{}(0),$$
(10)
where $`\psi _0(x)`$ is an approximate solution to Schrödinger’s equation with energy $`E_0`$, which is localized in the right hand well, and decays away from that well, in the entire central, classically forbidden region. Further, it is normalized to yield unit total probability in the right hand well. Note that it is not necessary to examine or specify the behaviour of $`\psi _0(x)`$ near $`x=a^{}`$. Lastly, $`\psi _0^{}(x)=d\psi _0(x)/dx`$.
We digress briefly at this point to comment on the formula (10), which does not seem to be widely known. It is sometimes named after Conyers Herring, who derived it (in a slightly more general form, in fact) in the course of evaluating the g-u splitting of the two lowest electronic states of the H$`{}_{}{}^{+}{}_{2}{}^{}`$ molecular ion . The actual derivation is simple, and very similar to the usual half-page of analysis used to argue that the the Hamiltonian operator in the position representation is Hermitean and has real eigenvalues, and so we refer readers to Landau and Lifshitz for the details. Instead, we present an alternative derivation in Appendix A.
Resuming our argument, we note that the WKB approximation for $`\psi _0(x)`$ in the region $`(a^{}x)(\mathrm{}/m\omega )^{1/2}`$ is
$$\psi _0(x)=\frac{C_0}{\sqrt{|v(x)|}}\mathrm{exp}\frac{1}{\mathrm{}}_0^x|p(x^{})|𝑑x^{},$$
(11)
where $`v(x)=p(x)/m`$, and $`C_0`$ is a constant to be found by matching on to the solution in the well. To the accuracy of this solution, $`\psi _0^{}(x)(m|v(x)|/\mathrm{})\psi _0(x)`$, so that
$$\mathrm{\Delta }=2\mathrm{}C_0^2.$$
(12)
To find $`C_0`$, we first note that near $`x=a`$, $`\psi _0(x)`$ is very accurately given by the ground state harmonic oscillator wave function
$$\psi _0(x)=\left(\frac{m\omega }{\pi \mathrm{}}\right)^{1/4}\mathrm{exp}\left[\frac{m\omega }{2\mathrm{}}(xa)^2\right].$$
(13)
In fact, this form holds well into the forbidden region, and so can be directly matched to Eq. (11) without invoking connection formulas. In the overlap region, we may write
$$|p(x)|=m\omega \left[(ax)^2u_0^2\right]^{1/2},$$
(14)
where
$$u_0=aa^{}(\mathrm{}/m\omega )^{1/2}.$$
(15)
The term $`u_0^2`$ in Eq. (14) may be neglected in evaluating the $`|v(x)|^{1/2}`$ prefactor in Eq. (11), so that we may write
$$\psi _0(x)\frac{C_0}{\sqrt{\omega (ax)}}\mathrm{exp}\left[\frac{1}{\mathrm{}}_0^a^{}|p(x^{})|𝑑x^{}+\mathrm{\Phi }(x)\right],$$
(16)
where
$`\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \frac{m\omega }{\mathrm{}}}{\displaystyle _x^a^{}}\left[(ax^{})^2u_0^2\right]^{1/2}𝑑x^{}`$ (17)
$``$ $`{\displaystyle \frac{m\omega (ax)^2}{2\mathrm{}}}+\mathrm{ln}\left({\displaystyle \frac{2(ax)}{u_0}}\right)^{1/2}+{\displaystyle \frac{1}{4}}+O\left({\displaystyle \frac{u_0}{ax}}\right)^2.`$ (18)
Comparing Eqs. (16) and (18) with Eq. (13), we obtain
$$C_0=\left(\frac{\omega ^2}{4\pi e}\right)^{1/4}\mathrm{exp}\left[\frac{1}{\mathrm{}}_0^a^{}|p(x)|𝑑x\right].$$
(19)
Substituting Eq. (19) in Eq. (12), and making use of the reflection symmetry of $`V(x)`$, we obtain
$$\mathrm{\Delta }=\frac{\mathrm{}\omega }{\sqrt{e\pi }}\mathrm{exp}\left[_a^{}^a^{}\frac{|p|}{\mathrm{}}𝑑x\right],$$
(20)
which is Eq. (1).
Landau and Lifshitz obtain $`C_0`$ by using connection formulas near a linear turning point, and the standard WKB result for the normalization of a bound state. For the ground state, this procedure is not accurate enough, and we must proceed as above. (Some readers may recognize in our procedure, the use of quadratic turning point formulas .) Once this is realized, similar prefactor corrections are expected for the low lying excited states. These corrections are obtained in Appendix B.
## III Simpler Expression for Ground State Splitting
In this section we will show that the WKB expression (1) leads to the simpler result (46). The objective, clearly, is to let the action integral run from $`a`$ to $`a`$ instead of from $`a^{}`$ to $`a^{}`$. To that end, let us denote $`y=ax`$, $`V(x)=U(y)`$, and define, for a general $`u`$,
$$I(u)=_u^a\left[2m\left(U(y)U(u)\right)\right]^{1/2}𝑑y.$$
(21)
The quantity $`I(u_0)`$ is clearly half the action integral in Eq. (1), i.e.,
$$I(u_0)=\frac{1}{2\mathrm{}}_a^{}^a^{}|p(x)|𝑑x.$$
(22)
Our goal is to expand $`I(u)`$ for small $`u`$. The analysis in Sec. II reveals \[see Eq. (18)\] that this expansion contains a term varying as $`u^2\mathrm{ln}u`$, which means that the expansion cannot be found by simply differentiating $`I(u)`$. More precisely, the term of order $`u^2\mathrm{ln}u`$ is easily found, but the term of order $`u^2`$, which we also need, is rather harder to get.
We therefore resort to a subtraction and write $`I(u)=I_1(u)+I_1^{}(u)`$, where
$`I_1(u)`$ $`=`$ $`{\displaystyle _u^a}\sqrt{2mU(y)}𝑑y,`$ (23)
$`I_1^{}(u)`$ $`=`$ $`{\displaystyle _u^a}\left(\sqrt{2m[U(y)U(u)]}\sqrt{2mU(y)}\right)𝑑y.`$ (24)
$`I_1(u)`$ can be directly expanded:
$`I_1(u)`$ $``$ $`{\displaystyle _0^a}\sqrt{2mU(y)}𝑑y{\displaystyle _0^u}\left(m\omega y+O(y^2)\right)𝑑y`$ (25)
$`=`$ $`I(0){\displaystyle \frac{1}{2}}m\omega u^2+O(u^3),`$ (26)
while for $`I_1^{}(u)`$ we have, as $`u0`$,
$$I_1^{}(u)(m\omega u)^2_u^a\frac{dy}{\sqrt{2m[U(y)U(u)]}+\sqrt{2mU(y)}}.$$
(27)
Note that in writing this expression, we have taken $`U(u)m\omega ^2u^2/2`$ in the numerator, as $`u`$ is small. It still can not be expanded directly, however, so we perform another subtraction, and write $`I_1^{}=I_2+I_3`$, where
$$I_2(u)=\frac{1}{2}(m\omega u)^2_u^a\frac{dy}{\sqrt{2m[U(y)U(u)]}}.$$
(28)
The difference $`I_3`$ contains a proportionality factor $`U(u)`$, which we again approximate by $`m\omega ^2u^2/2`$. This leads to
$$I_3(u)\frac{1}{2}(m\omega u)^4_u^a\frac{dy}{\sqrt{2m[U(y)U(u)]}\left[\sqrt{2m[U(y)U(u)]}+\sqrt{2mU(y)}\right]^2}.$$
(29)
Let us consider $`I_2(u)`$ first. We can clearly write the integral as
$$_u^a\frac{dy}{m\omega \sqrt{y^2u^2}}+_u^a\left[\frac{1}{\sqrt{2m[U(y)U(u)]}}\frac{1}{m\omega \sqrt{y^2u^2}}\right]𝑑y.$$
(30)
The first integral can be done exactly, while in the second we can set $`u=0`$ inside the integrand and in the limits to leading order. We thus find
$$I_2(u)\frac{1}{2}m\omega u^2\mathrm{ln}\frac{2a}{u}\frac{1}{2}(m\omega u)^2_0^a\left[\frac{1}{\sqrt{2mU(y)}}\frac{1}{m\omega y}\right]𝑑y+O(u^4\mathrm{ln}u).$$
(31)
The leading behaviour of $`I_3(u)`$, on the other hand, is controlled by the lower limit in the integral (29), where we can again approximate $`U(y)`$ by $`m\omega ^2y^2/2`$. Hence,
$`I_3(u)`$ $``$ $`{\displaystyle \frac{1}{2}}m\omega u^4{\displaystyle _u^a}{\displaystyle \frac{dy}{\sqrt{y^2u^2}\left[\sqrt{y^2u^2}+y\right]^2}}`$ (32)
$``$ $`{\displaystyle \frac{1}{4}}m\omega u^2+O(u^4).`$ (33)
Collecting together Eqs. (26), (31), and (33), and putting $`u=u_0=(\mathrm{}/m\omega )^{1/2}`$, we get
$$I(u_0)I(0)\frac{\mathrm{}}{2}\mathrm{ln}\frac{2a}{u_0}\frac{\mathrm{}}{4}\frac{\mathrm{}}{2}_0^a\left[\frac{m\omega }{\sqrt{2mU(y)}}\frac{1}{y}\right]𝑑y+\mathrm{}.$$
(34)
The last integral in the above equation is nothing but the quantity $`A`$ defined in Eq. (6). Using $`u_0=(\mathrm{}/m\omega )^{1/2}`$ once again, and putting together Eqs. (34), (22), and (1), we get
$$\mathrm{\Delta }=\frac{2\mathrm{}\omega }{\sqrt{e\pi }}\left(\frac{m\omega a^2}{\mathrm{}}\right)^{1/2}\sqrt{e}e^Ae^{2I(0)/\mathrm{}}.$$
(35)
Since $`I(0)=S_0/2`$, this is nothing but Eq. (4).
## IV Equivalence of WKB and Instanton Results
Our next step is to show that the instanton result (7) also leads to Eqs. (46), and thus prove the equivalence of the WKB and instanton results for the ground state splitting. To do this, we use Coleman’s result for the ratio $`K`$ . According to him,
$$K=\sqrt{2\omega }\beta ,$$
(36)
where $`\beta `$ is related to the asymptotic, $`\tau \pm \mathrm{}`$ behaviour of the instanton velocity via
$$x_1(\tau )\left(\frac{m}{S_0}\right)^{1/2}\frac{dx_{\mathrm{cl}}}{d\tau }\beta e^{\omega \tau }\text{as }\tau \pm \mathrm{}\text{.}$$
(37)
It is easy to integrate the equation of motion (9) for the instanton, and obtain $`x_1(\tau )`$. Using the fact that $`x_{\mathrm{cl}}(0)=0`$, we have
$$\tau =m_0^{x_{\mathrm{cl}}}\frac{dx}{\sqrt{2mV(x)}}.$$
(38)
This diverges as $`x_{\mathrm{cl}}a`$, as it should. We extract the divergence by subtracting and adding the leading singular part of the integrand. This yields,
$`\tau `$ $`=`$ $`m{\displaystyle _0^{x_{\mathrm{cl}}}}\left[{\displaystyle \frac{1}{\sqrt{2mV(x)}}}{\displaystyle \frac{1}{m\omega (ax)}}\right]𝑑x+{\displaystyle \frac{1}{\omega }}\mathrm{ln}\left({\displaystyle \frac{a}{ax_{\mathrm{cl}}}}\right)`$ (39)
$``$ $`{\displaystyle \frac{1}{\omega }}\mathrm{ln}\left({\displaystyle \frac{a}{ax_{\mathrm{cl}}}}\right)+{\displaystyle \frac{A}{\omega }}\text{as }x_{\mathrm{cl}}a\text{.}`$ (40)
$`A`$ is, of course, the quantity defined in Eq. (6).
Thus, as $`\tau \mathrm{}`$,
$`x_{\mathrm{cl}}(\tau )`$ $``$ $`aae^Ae^{\omega \tau },`$ (41)
$`{\displaystyle \frac{dx_{\mathrm{cl}}}{d\tau }}`$ $``$ $`a\omega e^Ae^{\omega \tau }.`$ (42)
Comparing with Eq. (37), we can read off $`\beta `$ immediately:
$$\beta =a\omega \left(\frac{m}{S_0}\right)^{1/2}e^A.$$
(43)
Hence,
$$K=a\omega \left(\frac{2m\omega }{S_0}\right)^{1/2}e^A,$$
(44)
and
$$\mathrm{\Delta }=2\mathrm{}\omega \left(\frac{m\omega a^2}{\pi \mathrm{}}\right)^{1/2}e^Ae^{S_0/\mathrm{}},$$
(45)
which is what we set out to show.
## V Illustrative Examples
We conclude with two elementary examples to which we apply Eqs. (46).
The first example is that of the quartic double well,
$$V(x)=V_0(x^2a^2)^2/a^4.$$
(46)
Here, $`V_0`$ is the barrier height. The frequency $`\omega `$ is given by
$$\omega =(8V_0/ma^2)^{1/2}.$$
(47)
It is simple to perform the integrals, and show that
$`{\displaystyle \frac{S_0}{\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{2mV_0a^2}{\mathrm{}^2}}\right)^{1/2}={\displaystyle \frac{16}{3}}{\displaystyle \frac{V_0}{\mathrm{}\omega }},`$ (48)
$`A`$ $`=`$ $`\mathrm{ln}2.`$ (49)
Further, $`(m\omega a^2/\pi \mathrm{})^{1/2}=(3S_0/2\pi \mathrm{})^{1/2}`$, so that
$$\mathrm{\Delta }=4\sqrt{3}\mathrm{}\omega \left(\frac{S_0}{2\pi \mathrm{}}\right)^{1/2}e^{S_0/\mathrm{}},$$
(50)
which is a well known form for the answer.
Our second example involves the spin Hamiltonian
$$=\gamma S_z^2\alpha S_x,$$
(51)
where $`S_\alpha `$ ($`\alpha =x,y,z`$) are components of the dimensionless spin operator $`𝐒`$ obeying the commutation rules $`[S_\alpha ,S_\beta ]=iϵ_{\alpha \beta \gamma }S_\gamma `$, and we also take $`\alpha >0`$, $`\gamma >0`$.
In the limit where the magnitude $`S`$ of the spin is very large , Eq. (51) can be viewed as tending to a classical Hamiltonian, wherein the dynamics are defined by giving the Poisson brackets $`[S_\alpha ,S_\beta ]_{\mathrm{PB}}=ϵ_{\alpha \beta \gamma }S_\gamma `$. The classical energy has two degenerate minima when the spin lies in the xz plane at angles $`\theta _0`$ or $`\pi \theta _0`$ to the z axis, where
$$\mathrm{sin}\theta _0=\alpha /2\gamma S.$$
(52)
For the quantum mechanical problem, with large but finite $`S`$, we expect the ground states around these classical orientations to be admixed by tunneling, thus giving rise to a small splitting of the energy levels.
Compared to the tunneling of massive particles with a position coordinate, the tunneling of spins is somewhat less familiar, but it is a perfectly bona fide instance of the general tunneling phenomenon. Spin tunnel splittings have been calculated by a variety of means for some time now , but obtaining the prefactors in the tunnel splittings correctly to order $`S^0`$ as $`S\mathrm{}`$, has proven to be a fairly difficult task, and the instanton calculations are especially subtle. The discrete WKB method uses only elementary methods of analysis, but the calculations to date are still lengthy. There is an even simpler method, however, which is due to Scharf, Wreszinski, and van Hemmen . These authors expand a general state $`|\psi `$ in the $`S_z`$ eigenbasis $`\{|m\}`$ ($`S_z|m=m|m`$) as $`|\psi =_mD_m|m`$, and construct a generating function $`_mD_mx^m`$. After a few changes of variables, the Schrödinger equation for $`|\psi `$ is turned into the following Schrödinger equation \[See their Eq. (2.17)\] for a wavefunction $`y(z)`$ related to the generating function:
$$\gamma ^2\frac{d^2y}{dz^2}+V(z)y(z)=\gamma Ey(z),$$
(53)
with the potential
$`V(z)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\alpha ^2(\mathrm{cosh}z\mathrm{cosh}z_0)^2,`$ (54)
$`\mathrm{cosh}z_0`$ $`=`$ $`(2S+1)\gamma /\alpha .`$ (55)
The quantity $`E`$ in Eq. (53) is the energy eigenvalue of the Hamiltonian (51).
The potential (54) is even about $`z=0`$, and has minima at $`\pm z_0`$. Our formalism is directly applicable if we identify $`\mathrm{}^2/2m`$ with $`\gamma ^2`$, and $`m\omega ^2`$ with $`V^{\prime \prime }(\pm z_0)=\frac{1}{2}\alpha ^2\mathrm{sinh}^2z_0`$. The action integral (5) is easily seen to be
$`{\displaystyle \frac{S_0}{\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{\gamma }}{\displaystyle _0^{z_0}}(\mathrm{cosh}z_0\mathrm{cosh}z)𝑑z`$ (56)
$`=`$ $`{\displaystyle \frac{\alpha }{\gamma }}(z_0\mathrm{cosh}z_0\mathrm{sinh}z_0),`$ (57)
while the correction (6) is given by
$$A=_0^{z_0}\left[\frac{\mathrm{sinh}z_0}{\mathrm{cosh}z_0\mathrm{cosh}z}\frac{1}{z_0z}\right]𝑑z.$$
(58)
Since the integrand is nonsingular at $`z=z_0`$ (or if one so wishes to view matters, has a removable singularity at that point), the integral can be found by changing the upper limit to $`z_0\delta `$, integrating the two terms separately, and letting $`\delta 0`$ at the end. The integrals involved are elementary, and the final result is
$$A=\mathrm{ln}\left(\frac{2\mathrm{sinh}z_0}{z_0}\right).$$
(59)
It remains to substitute Eqs. (57) and (59) into our general formula (4). Recalling the factor of $`\gamma `$ on the right hand side of Eq. (53), and the equivalences $`\mathrm{}^2/2m\gamma ^2`$, $`m\omega ^2V^{\prime \prime }(z_0)`$, we obtain
$$\mathrm{\Delta }=\frac{4\alpha ^{3/2}\mathrm{sinh}^{5/2}z_0}{\sqrt{2\pi \gamma }}\mathrm{exp}\left[\frac{\alpha }{\gamma }(z_0\mathrm{cosh}z_0\mathrm{sinh}z_0)\right].$$
(60)
It is easier to interpret this result if it is cast in terms of the angle $`\theta _0`$ in Eq. (52). After a certain amount of tedious algebra, one obtains, correct to order $`S^0`$ as $`S\mathrm{}`$,
$$\mathrm{\Delta }=\frac{8\gamma S^{3/2}\mathrm{cos}^{5/2}\theta _0}{\pi ^{1/2}\mathrm{sin}\theta _0}\left(\frac{1\mathrm{cos}\theta _0}{1+\mathrm{cos}\theta _0}\right)^{S+1/2}e^{2S\mathrm{cos}\theta _0}.$$
(61)
This is precisely the form quoted in Ref. . Its relation to previous calculations is discussed there.
From the viewpoint of this article, the interesting comparison is with Scharf, Wreszinski, and van Hemmen . These authors only consider the case where $`\alpha `$ and $`\gamma `$ are fixed as $`S\mathrm{}`$, so that $`\theta _0=O(1/S)`$. In that limit, they follow the Landau-Lifshitz prescription, and obtain a ground state splitting $`\mathrm{\Delta }^{}=(2S/\pi )e^J`$, with
$$J=2S\mathrm{ln}\left(\frac{8}{\alpha }\gamma ^2S^2\right)+2S+(2S+1)\mathrm{ln}\left[\gamma \left(2S+\frac{1}{2}\right)\right]\frac{1}{2}\mathrm{ln}\left[\gamma ^2\left(S+\frac{1}{4}\right)\right].$$
(62)
With a little work one can show that $`\mathrm{\Delta }^{}/\mathrm{\Delta }`$ is exactly $`(e/\pi )^{1/2}`$. The discrepancy can in fact be seen in the comparison between numerical and analytical results in Table 1 of Ref. . While $`(e/\pi )^{1/2}=0.930`$, the ratio of their analytical answer for the splitting to the numerical one (see the last two columns) decreases from 0.987 to 0.951 as $`S`$ increases from 5 to 11, which is quite close. A Richardson transformation of this ratio does suggest that it tending to 0.93, but one cannot be certain of this conclusion, as there are two erratic values at 0.90 and 0.91.
###### Acknowledgements.
This work is supported by the National Science Foundation via grant number DMR-9616749.
## A Herring’s formula
Let us denote the two states localized in the left and right wells by $`|R`$ and $`|L`$ respectively. These states are degenerate in the absence of tunneling with an energy $`E_0`$. The tunnel splitting is $`\mathrm{\Delta }`$. The Hamiltonian matrix in this two-state subspace is given by
$$=\left(\begin{array}{cc}E_0& \mathrm{\Delta }/2\\ \mathrm{\Delta }/2& E_0\end{array}\right)$$
(A1)
Assuming that the system starts in the state $`|R`$ at $`t=0`$, it is straightforward to show that the probability $`P_R(t)`$ for finding it in the same state at a later time $`t`$ is given by
$$P_R(t)=\mathrm{cos}^2(\mathrm{\Delta }t/2\mathrm{}).$$
(A2)
In particular,
$$\frac{dP_R}{dt}=\frac{\mathrm{\Delta }}{2\mathrm{}}\mathrm{sin}\left(\frac{\mathrm{\Delta }t}{\mathrm{}}\right).$$
(A3)
To relate this abstract space description to that in position space, we make use of the continuity equation for probability. For a general wave function $`\psi (x,t)`$ that obeys Schrödinger’s equation, the probability density $`P(x,t)=|\psi (x,t)|^2`$ obeys
$$\frac{P(x,t)}{t}=\frac{\mathrm{}}{m}\frac{}{x}\mathrm{Im}[\psi ^{}(x,t)\psi ^{}(x,t)],$$
(A4)
where $`\psi ^{}\psi /x`$. The probability $`P_R(t)`$ for being in the right well is then given by
$$P_R(t)=_0^{\mathrm{}}P(x,t)𝑑x.$$
(A5)
Differentiating with respect to $`t`$, making use of the continuity equation, and the fact that $`\psi (x,t)0`$ as $`x\mathrm{}`$ for any well behaved solution, we obtain
$$\frac{dP_R(t)}{dt}=\frac{\mathrm{}}{m}\mathrm{Im}[\psi ^{}(0,t)\psi ^{}(0,t)].$$
(A6)
Let us finally consider the states $`|R`$ and $`|L`$ in position space. It is evident that we should take $`x|R`$ to be the function described as $`\psi _0(x)`$ in Sec. II. By symmetry, $`x|L=\psi _0(x)`$, and the energy eigenstates are $`(\psi _0(x)\pm \psi _0(x))/\sqrt{2}`$. It then follows that with the same initial conditions as above, the wave function at an arbitrary time $`t`$ is given by
$$\psi (x,t)=\psi _0(x)\mathrm{cos}(\mathrm{\Delta }t/2\mathrm{})+i\psi _0(x)\mathrm{sin}(\mathrm{\Delta }t/2\mathrm{}).$$
(A7)
Hence,
$$\mathrm{Im}[\psi ^{}(0,t)\psi ^{}(0,t)]=\psi _0(0)\psi _0^{}(0)\mathrm{sin}(\mathrm{\Delta }t/\mathrm{}).$$
(A8)
Substituting this result in Eq. (A6) and comparing with Eq. (A3), we obtain
$$\mathrm{\Delta }=(2\mathrm{}^2/m)\psi _0(0)\psi _0^{}(0),$$
(A9)
which is Eq. (10), Herring’s formula.
It is apparent from our derivation that Herring’s formula holds whenever the energy eigenfunctions are well approximated by the combinations $`(\psi _0(x)\pm \psi _0(x))/\sqrt{2}`$. (More precise statements of the conditions for its validity are given in Ref. .) Use of this formula greatly simplifies the labour required to solve all the standard double-well problems: the double square well with infinite side walls , the double delta function , and the double harmonic oscillator with a kink . For smooth potentials, where WKB is indicated, it is far superior to the approach where one uses connection formulas at all four turning points . Further, the formula brings out the physical point that a tunnel splitting is intimately related to a tunneling amplitude, since it relates the amplitude to make a transition between wells per unit time, $`i\mathrm{\Delta }/\mathrm{}`$, to the probability current.
The reader will also have noticed that the above argument cannot be carried out for an asymmetric potential, and indeed the very concept of tunnel splitting is then very delicate. A tunneling amplitude can of course still be defined, but since this amplitude is exponentially small in general (on account of the Gamow factor), mixing between left and right well states will be negligible unless the bottoms of the two wells are tuned to the same exponential sensitivity. The mathematical formulation of these points leads to extremely unpleasant transcendental equations, and the situation is not significantly improved in the case of symmetric potentials if one approaches the problem solely in terms of enforcing continuity of the wave function and its derivative. It is probably because of this fact that most introductory or intermediate quantum mechanics texts do not consider double-well tunnel splittings when they discuss the WKB method. Use of Herring’s formula (with or without WKB) would make the problem much more tractable, and its widespread adoption is thus greatly to be desired.
## B WKB formula for higher state splittings
It is to be expected that the formula given by Landau and Lifshitz is increasingly accurate for higher pairs of states, assuming of course, that the WKB approximation is still applicable. We will find that this is indeed so. The splitting for the $`n`$th pair of levels, $`\mathrm{\Delta }_n`$, is given by
$$\mathrm{\Delta }_n=g_n\frac{\mathrm{}\omega }{\pi }\mathrm{exp}\left[_{a_n^{}}^{a_n^{}}\frac{|p|}{\mathrm{}}𝑑x\right],$$
(B1)
where $`\pm a_n^{}`$ are the classical turning points for the $`n`$th energy level pair $`(n+\frac{1}{2})\mathrm{}\omega `$, and
$$g_n=\frac{\sqrt{2\pi }}{n!}\left(n+\frac{1}{2}\right)^{n+\frac{1}{2}}e^{(n+\frac{1}{2})}.$$
(B2)
Note that if $`g_n`$ were unity, we would have the formula of Ref. . The corrections are indeed small: $`g_0=(\pi /e)^{1/2}1.075`$, $`g_11.028`$, $`g_21.017`$, and so on. Stirling’s formula for $`n!`$ shows that $`g_n1`$ as $`n\mathrm{}`$.
The derivation of Eq. (B1) and the conditions under which it holds are straightforward though long. The starting point is still Herring’s formula, Eq. (10). The procedure of Sec. II can be applied word for word, with $`a^{}`$, $`u_0`$, and $`C_0`$ replaced by $`a_n^{}`$, $`u_n`$, and $`C_n`$, respectively. The essential part is to find $`\psi (x)`$ in the classically forbidden region and match it onto the $`n`$th harmonic oscillator state in the vicinity of the well $`xa`$. We leave it as an exercise — or see Ref. — to show that the function $`\mathrm{\Phi }(x)`$, defined in Eq. (16), is given in the overlap region by
$$\mathrm{\Phi }(x)=\frac{1}{2}\frac{m\omega }{\mathrm{}}(ax)^2+\frac{2n+1}{2}\mathrm{ln}\left(\frac{2(ax)}{u_n}\right)+\frac{2n+1}{4}+\mathrm{},$$
(B3)
so that the leading approximation to $`\psi _0(x)`$ is given by $`C_n^{}(ax)^n\mathrm{exp}(x^2/2u_0)`$, where
$$C_n^{}=\frac{C_n}{\omega ^{1/2}}\left(\frac{2}{u_n}\right)^{n+\frac{1}{2}}\mathrm{exp}\left(\frac{1}{\mathrm{}}_0^{a_n^{}}|p|𝑑x+\frac{2n+1}{4}\right)$$
(B4)
This is of the precise form required to be match on to the $`n`$th excited state harmonic oscillator wave function:
$$\psi _0(x)\left(\frac{m\omega }{\pi \mathrm{}}\right)^{1/4}\frac{2^{n/2}}{\sqrt{n!}}e^{\xi ^2/2}\left(\xi ^n+O(\xi )^{n2}\right),$$
(B5)
with $`\xi (ax)/u_0`$. The match yields the constant $`C_n`$, and the splitting $`\mathrm{\Delta }_n`$, which equals $`2\mathrm{}|C_n|^2`$, is easily shown to be given by Eq. (B1) and (B2).
The calculation sketched above assumes that the turning point and the wavefunction in its vicinity are well approximated by taking the potential well to be parabolic. This is clearly less accurate as $`n`$ gets large. Defining $`y`$ and $`U(y)`$ as in Sec. IV, let us keep the cubic and quartic terms in $`y`$ in $`U(y)`$:
$$U(y)=\frac{1}{2}m\omega ^2y^2+\alpha y^3+\beta y^4.$$
(B6)
The wavefunction for the $`n`$th state can be found by perturbation theory assuming that $`\alpha `$ and $`\beta `$ are small. The key requirement is that the correction terms be small compared to the unperturbed wavefunction near the turning point $`y=u_n`$. The dominant corrections are those that entail the Hermite polynomials of order $`n+1`$ through $`n+4`$. Since the coefficient of the largest power of $`x`$ in these polynomials is known, the conditions for the corrections to be small are not difficult to find. They are
$`\alpha `$ $``$ $`{\displaystyle \frac{6}{13n+11}}(2n+1)^{1/2}\left({\displaystyle \frac{m^3\omega ^5}{\mathrm{}}}\right)^{1/2},`$ (B7)
$`\beta `$ $``$ $`{\displaystyle \frac{4}{(6n+7)(2n+1)}}{\displaystyle \frac{m^2\omega ^3}{\mathrm{}}}.`$ (B8)
It is also possible to show that when these conditions hold, the corrections to $`u_n`$ and the energies of the harmonic oscillator states are negligible, so that the conditions are self-consistently derived.
It need hardly be said that the conditions (B7) and (B8) are only necessary, not sufficient, for Eq. (B1) to apply. In addition, one must have the requirement for the WKB approximation itself to hold, which in the present context can be stated as $`(n+\frac{1}{2})\mathrm{}\omega V(0)V(a)`$. |
warning/0003/cond-mat0003324.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Random Matrix Models have been introduced in order to give an approximate statistical description of quantum systems involving disorder, chaos, complexity or whatever prevents from solving the equations of motion exactly. Those models are described by a matrix (Hamiltonian, transfer matrix, or scattering matrix) of large size $`N`$, which is too complicated to be diagonalized exactly, and for which only statistical observations of the spectrum are available (see for a review on RMT).
Most of the quantities of interest and observables are related to the short range (in energy scale) behavior of the spectrum (indeed small energies correspond to long time evolution, i.e. to equilibrium thermodynamical properties). This is why the short range correlation functions are the most studied.
At first, the simplest models assumed a gaussian weight for the random matrix and gave good agreement with observations, provided that the ensemble of matrices(hermitian, orthogonal, quaternionic…) has the required symmetries (time reversibility,…) .
It has been observed that the correlation functions of the spectrum possess universal properties at sort range, which do not depend on the probability weight, gaussian or not. This universality has been proved for a wide range of models by several approaches , but a very general proof and the exact hypothesis which lead to it are still under investigations.
Moreover, the long range correlation functions appear to share also some universal properties, which depend on the probability weight only through a few parameters . The most striking example is the 2-point correlation function that we shall discuss below.
In the following we will restrict our attention to the so called Hermitian-One-Matrix-Model<sup>5</sup><sup>5</sup>5 The other ensembles are of course worth considering, but this one is the simplest. (hermiticity corresponds to a system with broken time reversibility, for instance in the presence of a magnetic field).
We consider a hermitian matrix $`M`$ of size $`N\times N`$ with a probability law of the form:
$$𝒫(M)=\mathrm{e}^{N\mathrm{tr}V(M)}$$
where $`V`$ is a polynomial potential bounded from below.
We wish to study the statistical properties of the eigenvalues $`(\lambda _1,\mathrm{},\lambda _N)`$ of $`M`$ in the large $`N`$ limit, in particular the density of eigenvalues $`\rho (\lambda )`$, and the correlation function $`R(\lambda ,\mu )`$, which measures the probability that two of the eigenvalues take the values $`\lambda `$ and $`\mu `$.
Roughly speaking, the eigenvalues tend to occupy a finite interval centered around the bottom of the potential well, and in the large $`N`$ limit, the density of eigenvalues $`\rho (\lambda )`$ is a continuous function with a compact support.
The simplest case, where the support is connected, known as the “1-cut case”, has been extensively studied . It is then found that the density $`\rho (\lambda )`$ is not universal, it depends on the details of the potential $`V`$, while the connected correlation function $`R_c(\lambda ,\mu )`$ is universal in the short range regime ($`|\lambda \mu |O(1/N)`$), but also in the long range regime ($`|\lambda \mu |O(1)`$) once the short range oscillations (of period $`O(1/N)`$) have been smoothed out.
What happens when the potential $`V`$ possesses several wells, of approximatively the same depths? Then the density $`\rho (\lambda )`$ has a disconnected support, $`[a_1,b_1][a_2,b_2]\mathrm{}[a_s,b_s]`$, each interval $`[a_i,b_i]`$ being centered around one well of the potential $`V`$. This case is known as the multicut (or multiband) case (here $`s`$ cuts).
In the multicut case the density is still not universal, whereas the 2-point correlation function is universal in the short range regime and seems to have some universal properties in the long range regime after smoothing: in an explicit form of the 2-point connected correlation function was given, and is claimed to be universal: indeed, according to the authors of “it depends only on the number of connected components of the support and on the position of the endpoints, but not on the potential”. However, more recently several authors have studied the two-cut case $`s=2`$: they concentrated on the case of an even potential $`V`$ (the two cuts are thus symmetric $`[a,b]`$ and $`[b,a]`$). Using an ansatz for the asymptotic expression of orthogonal polynomials in the large $`N`$ limit, and rederiving the two-point function from this ansatz, they observed that the connected correlation function is still universal in the short distance regime (which was expected), but more surprisingly, that the smoothed connected correlation function in the long range regime depends on the parity of $`N`$ ($`N`$ being the size of the matrix). This seems to contradict the former result of !
In this paper, we will solve this paradox.
We will show that the semi-classical method of gives the 2-point connected correlation function only up to an additional non-universal term, which is already present in the free energy, but subdominant at large $`N`$ in this case. We correct the semi-classical argument of , and give a simple (and physically appealing) derivation of the origin of the additional term, that we compute explicitly. This allows us to recover the results of for the symmetric case, and to generalize them to non symmetric potentials, without using orthogonal polynomials.
Using the same semi-classical argument, we are able to derive large $`N`$ asymptotics for the orthogonal polynomials, recovering the results of as well as the general $`s`$ cuts asymptotics which appeared recently in the mathematical literature , and to extend these results to the case of complex potentials.
The effect leading to the new term in the semi-classical calculation is simple enough to be explained briefly in this introductory section.
In , the free energy $`F`$ of the matrix model is derived by a saddle point approximation. In particular, one has to extremize the action with respect to variations of the number $`n_i`$ ($`i=1\mathrm{}s`$) of eigenvalues in each connected part of the support, or in other words with respect to the occupation ratio $`x_i=n_i/N`$:
$$F=F(x_c)\mathrm{where}\frac{F}{x}|_{x=x_c}=0$$
(1.1)
However, one has here missed the crucial fact that $`n_i=Nx_i`$ are not real numbers but integers. When $`Nx_c`$ is not an integer, the extremum of $`F(x)`$ is never reached, and the saddle point approximation has to be slightly modified. Roughly speaking the discrete sum cannot be approximated by an integral:
$$\underset{n}{}\mathrm{e}^{g(nNx_c)^2}\text{d}x\mathrm{e}^{N^2g(xx_c)^2}$$
(1.2)
The discrete sum actually depends on how far from an integer $`Nx_c`$ is. For instance, in the symmetric case, we have $`x_c=\frac{1}{2}`$, and the result depends on the parity of $`N`$.
We will show that (as expected from Eq. (1.2)) in general the result involves elliptic theta functions depending on $`Nx_c`$, thus leading to a quasi-periodic dependence on $`N`$. This effect is of order $`N^2`$ for the free energy, but is of order $`1`$ for the computation of the orthogonal polynomials and for the correlation functions. It implies in particular that there is no regular large $`N`$ topological expansion (involving only power series in $`N^2`$) for the 2-cut matrix model.
We will find out that the short range correlation function is universal, while the long range smoothed correlation depends on $`N`$ quasi-periodically.
The paper is divided as follows: in section 2 we introduce the method and notations for the 2-cut model, and we compute the free energy. In section 3 we derive the 2-point correlation function, and we recover the expression of in the symmetric case. In section 4, we give an asymptotic expression for the orthogonal polynomials, which we use to rederive the universal short range properties of the spectrum, as well as the smoothed long range 2-point correlation function.
The generalizations to a complex potential or to an arbitrary number of cuts are presented in Appendix B and C. Appendix A is a summary of some relationships between elliptical functions in case the reader is not familiar with them.
## 2 The free energy
### 2.1 Basics
We start from the standard Hermitian matrix model defined by the partition function
$$Z[V;N]=\mathrm{d}_N[M]\mathrm{e}^{N\mathrm{tr}V(M)}$$
(2.1)
where $`N`$ is the dimension of the matrix $`M`$, $`V`$ is an analytic – in general polynomial – and for the moment real – function, and $`\text{d}_N[M]`$ is the standard $`U(N)`$ invariant measure over Hermitian matrices
$$\mathrm{d}_N[M]=\underset{i=1}{\overset{N}{}}\text{d}M_{ii}\underset{1i<jN}{}2\text{d}\mathrm{Re}(M_{ij})\text{d}\mathrm{Im}(M_{ij})$$
(2.2)
Integrating out the ”angular part” of $`M`$, $`Z`$ can be rewritten as an integral over the $`N`$ eigenvalues $`\lambda _1,\mathrm{},\lambda _N`$ of $`M`$
$$Z[V;N]=𝐂_N\stackrel{~}{Z}[V;N]$$
(2.3)
$$\stackrel{~}{Z}[V;N]=\underset{k=1}{\overset{N}{}}\text{d}\lambda _k\mathrm{e}^{N_kV(\lambda _k)}\underset{k<l}{}(\lambda _k\lambda _l)^2=\underset{k=1}{\overset{N}{}}\text{d}\lambda _k\mathrm{e}^{N^2S(\lambda _1,\mathrm{},\lambda _N)}$$
(2.4)
with the measure factor
$$C_N=\mathrm{Vol}\left[\frac{\mathrm{U}(N)}{\mathrm{U}(1)^N\times 𝔖_N}\right]=\frac{1}{N!}\underset{K=1}{\overset{N}{}}\frac{(2\pi )^{K1}}{\mathrm{\Gamma }(K)}$$
(2.5)
and with the action $`S(\lambda _k)`$:
$$S(\lambda _1,\mathrm{},\lambda _N)=\frac{1}{N}\underset{k=1}{\overset{N}{}}V(\lambda _k)\frac{1}{2N^2}\underset{1klN}{}\mathrm{ln}(\lambda _k\lambda _l)^2$$
(2.6)
In the simplest ”one cut” case, corresponding in particular to a concave potential, it is known that the free energy $`F`$ defined as <sup>6</sup><sup>6</sup>6 $`Z`$ is normalized here so that $`F`$ is zero for the Gaussian model $`V=\frac{1}{2}M^2`$
$$Z[V;N]=\left(\frac{2\pi }{N}\right)^{\frac{N^2}{2}}\mathrm{e}^{F[V;N]}$$
(2.7)
has a topological large $`N`$ expansion
$$F=N^2F_0+F_1+N^2F_2+\mathrm{}$$
(2.8)
obtained for instance by re-organizing the perturbative expansion according to the topology of the Feynman diagrams. The large $`N`$ limit (planar limit) can be described by a ”master field” configuration where the eigenvalues are described by a continuous density $`\rho (\lambda )`$ with a connected compact support $`𝒞=[a,b]`$, with the constraints
$$_𝒞\text{d}\lambda \rho (\lambda )=1\mathrm{and}\rho (\lambda )0\mathrm{if}\lambda 𝒞$$
(2.9)
and the action 2.6 becomes
$$S[\rho ]=_𝒞\text{d}\lambda V(\lambda )\rho (\lambda )_{𝒞\times 𝒞}\text{d}\lambda \text{d}\mu \rho (\lambda )\rho (\mu )\mathrm{ln}|\lambda \mu |$$
(2.10)
The leading term of the free energy $`F_0`$ can be obtained by the saddle point method: the effective action $`S[\rho ]`$ is extremized for a continuous distribution $`\rho _c`$ and we have simply
$$F_0=S[\rho _c]$$
(2.11)
(up to an additive – potential independent – constant). To compute $`\rho _c`$ we include the constraint 2.9 in the effective action by a Lagrange multiplier $`\mathrm{\Gamma }`$
$$\overline{S}[\rho ]=S[\rho ]+\mathrm{\Gamma }(1_𝒞\rho )$$
(2.12)
The saddle point equation for $`\rho `$ reads:
$$\frac{\overline{S}}{\rho (\lambda )}=V(\lambda )2_𝒞\text{d}\mu \rho (\mu )\mathrm{ln}(|\lambda \mu |)\mathrm{\Gamma }=0\lambda 𝒞$$
(2.13)
which simply means that the real part of the effective potential
$$V_{\mathrm{eff}}(\lambda )=V(\lambda )2_𝒞\text{d}\mu \rho (\mu )\mathrm{ln}(\lambda \mu )$$
(2.14)
is constant on the e.v. support $`𝒞`$, and equal to $`\mathrm{\Gamma }`$. The derivative of 2.13 w.r.t. $`\lambda `$ gives the well-known equation
$$\mathrm{Re}\left(\omega _0(\lambda )\right)=_𝒞\text{d}\mu \rho (\mu )\frac{1}{\lambda \mu }=V^{}(\lambda )/2\lambda 𝒞$$
(2.15)
where $`\omega _0`$ is the large $`N`$ resolvent
$$\omega _0(\lambda )=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{Tr}\left[\frac{1}{\lambda M}\right]=_𝒞\text{d}\mu \frac{\rho (\mu )}{\lambda \mu }$$
(2.16)
Finally let us recall that in the one-cut case $`𝒞=[a,b]`$, if the potential $`V`$ is a polynomial of degree $`P`$, $`\omega `$ is of the form
$$\omega _0(\lambda )=\frac{V^{}(\lambda )}{2}\frac{M(\lambda )\sqrt{\sigma (\lambda )}}{2}\mathrm{with}\sigma (\lambda )=(\lambda a)(\lambda b)$$
(2.17)
where $`M(\lambda )`$ is a polynomial with degree $`P2`$. $`a`$, $`b`$ and $`M`$ are entirely determined by the constraint that
$$\omega _0(\lambda )=\lambda ^1+𝒪(\lambda ^2)\mathrm{for}\lambda \mathrm{}$$
(2.18)
The e.v. density is given by the discontinuity of $`\omega `$
$$\rho (\lambda )=\frac{\mathrm{i}}{2\pi }\left[\omega (\lambda +\mathrm{i0}_+)\omega (\lambda \mathrm{i0}_+)\right]=\frac{M(\lambda )\sqrt{|\sigma (\lambda )|}}{2\pi }$$
(2.19)
### 2.2 The 2-cut case
#### 2.2.1 Mean field:
If the potential $`V`$ is real but has more than one minimum, the large $`N`$ limit may be described by an e.v. distribution on several disconnected intervals. For simplicity we shall first consider the case where there are two intervals
$$𝒞=𝒞_1𝒞_2,𝒞_1=[a,b],𝒞_2=[c,d],a<b<c<d$$
(2.20)
In this case, as we shall see, there is no topological large $`N`$ expansion, even for the free energy $`F`$. As shown in , to describe the large $`N`$ limit, we have to consider as an additional variable the ”average” proportion of eigenvalues $`x_1=n_1/N`$ and $`x_2=n_2/N`$ in each interval $`𝒞_1`$ and $`𝒞_2`$, and introduce the associated Lagrange multipliers $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ for the constraints
$$x_\alpha =_{𝒞_\alpha }\rho (\lambda )\text{d}\lambda ,\alpha =\mathrm{\hspace{0.17em}1},2$$
(2.21)
The effective action 2.12 now reads, with
$$x=x_1$$
(2.22)
$$\overline{S}[\rho ;x]=S[\rho ]+\underset{\alpha =1}{\overset{2}{}}\mathrm{\Gamma }_\alpha (x_\alpha _{𝒞_\alpha }\rho (\lambda )\text{d}\lambda ),x_1+x_2=1$$
(2.23)
with $`S[\rho ]`$ given by 2.10 as before. The saddle point equation w.r.t. $`\rho (\lambda )`$ gives as before the equation 2.13, which implies that the effective potential defined by 2.14 is constant on each interval
$$V_{\mathrm{eff}}(\lambda )=\mathrm{\Gamma }_\alpha \mathrm{when}\lambda 𝒞_\alpha $$
(2.24)
but the corresponding e.v. density $`\rho _c(\lambda )`$ and the effective action $`\overline{S}_c`$ still depend explicitly of the e.v. proportion $`x`$, since we have
$$\overline{S}_c[x]=\frac{1}{2}\left(_𝒞\rho _c(\lambda )V(\lambda )\text{d}\lambda +\underset{\alpha }{}\mathrm{\Gamma }_\alpha x_\alpha \right)$$
(2.25)
The saddle point equation w.r.t. $`x`$ implies the equality of the effective potentials for each interval
$$\frac{\overline{S}}{x}=\mathrm{\Gamma }_1\mathrm{\Gamma }_2=0$$
(2.26)
This fixes the value of $`x`$, and it is known that with this last equation the e.v. density $`\rho _c`$ is uniquely determined in the 2-cut case. The large $`N`$ free energy is then given simply by
$$F_0=\overline{S}[\rho _c;x_c]=\overline{S}_c[x_c]$$
(2.27)
For an explicit polynomial potential $`V`$ of degree $`P`$, and for fixed $`x`$, the 2-cut mean field solution for the resolvent is
$$\omega _0(\lambda ,x)=\frac{V^{}(\lambda )}{2}\frac{M(\lambda )\sqrt{\sigma (\lambda )}}{2}\mathrm{with}\sigma (\lambda )=(\lambda a)(\lambda b)(\lambda c)(\lambda d)$$
(2.28)
and $`M(\lambda )`$ a polynomial with degree $`P3`$. The e.v. density $`\rho (\lambda ,x)`$ is still given by the discontinuity of $`\omega _0`$. The coefficients of $`M`$ and the 4 end-points $`a,b,c,d`$ are entirely determined by the constraint that $`\omega _0(\lambda )\lambda ^1`$ when $`\lambda \mathrm{}`$ and by the fact that $`x`$ must be given by
$$x=_a^b\rho (\lambda ,x)\text{d}\lambda =\frac{1}{2\pi }_a^b|M(\lambda ,x)|\sqrt{|\sigma (\lambda )|}\text{d}\lambda $$
(2.29)
Finally the equation 2.26 which fixes $`x=x_c`$ reads
$$0=V_{\mathrm{eff}}(b)V_{\mathrm{eff}}(c)=_b^c\text{d}\lambda \left(2\omega _0(\lambda ,x)V^{}(\lambda )\right)=_b^c\text{d}\lambda M(\lambda ,x)\sqrt{|\sigma (\lambda )|}$$
(2.30)
#### 2.2.2 Discreteness of number of e.v.’s:
This is sufficient if one is interested in the leading term in the large $`N`$ limit (planar approximation). However, in order to understand the structure of the subdominant terms of the large $`N`$ expansion, it turns out that we cannot neglect the fact that the number of e.v. $`n_\alpha =Nx_\alpha `$ in each interval $`𝒞_\alpha `$ must be an integer.
Let us consider the simple case where the potential $`V`$ has two separate minima $`z_1`$ and $`z_2`$. What has to be done is first to fix the number of eigenvalues $`n_1=n`$ (resp. $`n_2=Nn`$) in the vicinity of $`z_1`$ (resp. $`z_2`$) in the partition function 2.5 by writing
$$\stackrel{~}{Z}[V;N]=\underset{n=0}{\overset{N}{}}\frac{N!}{n!(Nn)!}\stackrel{~}{Z}[V;n,Nn]$$
(2.31)
where
$$\stackrel{~}{Z}[V;n,Nn]=_{\mathrm{}}^E\underset{in}{}\text{d}\lambda _i_E^+\mathrm{}\underset{j>n}{}\text{d}\lambda _j\mathrm{e}^{N_kV(\lambda _k)}\underset{k<l}{}(\lambda _k\lambda _l)^2$$
(2.32)
with $`E`$ a ”frontier” $`b<E<c`$ between the two semi-classical cuts $`[a,b]`$ and $`[c,d]`$. We now claim that each term of this discrete sum has a well defined large $`N`$ topological expansion. Indeed, we can rewrite 2.32 as a matrix integral over two separate matrices: a $`n_1\times n_1`$ matrix $`M_1`$ with the $`n_1=n`$ e.v. $`<E`$ and a $`n_2\times n_2`$ matrix $`M_2`$ with the $`n_2=Nn`$ e.v. $`>E`$, as
$$\stackrel{~}{Z}[V;n]=\frac{1}{C_nC_{Nn}}\mathrm{d}_{n_1}[M_1]\mathrm{d}_{n_2}[M_2]\mathrm{e}^{N\mathrm{Tr}\left(V(M_1)\right)N\mathrm{Tr}\left(V(M_2)\right)+2\mathrm{Tr}\left(\mathrm{ln}\left(M_1\mathrm{Id}+\mathrm{Id}M_2\right)\right)}$$
(2.33)
This last matrix integral has a topological large $`N`$ expansion of the form 2.8 in the ‘t Hooft limit $`N\mathrm{}`$, $`x=n/N`$ fixed, obtained by doing a classical perturbative expansion around the smallest minimum $`z_1`$ of $`V`$ for $`M_1`$ and around the largest minimum $`z_2`$ of $`V`$ for $`M_2`$, and by re-organizing the perturbative expansion according to the topology of the Feynman diagrams. Taking into account carefully the measure factors $`C_n`$ and $`C_{Nn}`$, and using their large $`N`$ asymptotics
$$C_N=\frac{1}{N!}\left(\frac{2\pi }{N}\right)^{\frac{N^2}{2}}\mathrm{e}^{\frac{3}{4}N^2}(2\pi )^NN^{\frac{1}{12}}\mathrm{cst}(1+𝒪(N^1))\text{ when }N\mathrm{}$$
(2.34)
(easily derived from Stirling formula),we obtain that
$$Z[V;N]=\left(\frac{2\pi }{N}\right)^{\frac{N^2}{2}}N^{\frac{1}{12}}\underset{n=0}{\overset{N}{}}\mathrm{e}^{F[V;N,x]}$$
(2.35)
where each $`F[V;N,x]`$ has a regular large $`N`$ asymptotic expansion of the form
$$F[V;N,x]=\underset{h=0}{\overset{\mathrm{}}{}}N^{22h}F_h[V,x]\mathrm{where}x=n/N$$
(2.36)
with each $`F_h[V,x]`$ a regular function of $`x=n/N`$. In particular, the leading large $`N`$ term is given (up to an additive – $`V`$ and $`x`$ independent – constant) by the classical effective action 2.27
$$F_0=\overline{S}[\rho _c;x_c]=\overline{S}_c[x_c]$$
(2.37)
Finally, let us stress that although this decomposition depends on the arbitrary parameter $`E`$, since $`E`$ is in the interval $`]b,c[`$ where the density of eigenvalues is exponentially small with $`N`$, the integral 2.32 depends on $`E`$ only through exponentially small terms of order $`\mathrm{e}^{\mathrm{cst}N}`$, which are “non-perturbative” in the topological expansion Eq. (2.36).
#### 2.2.3 Beyond mean-field:
We can now easily calculate the subleading terms of order $`𝒪(N^2)`$ for the full partition function. In the large $`N`$ limit we can approximate the sum 2.36 by
$$Z[V;N]\underset{n=0}{\overset{N}{}}\mathrm{e}^{N^2F_0[V;x]F_1[V;x]+\mathrm{}}$$
(2.38)
If $`x_c`$ denotes the saddle point of $`F_0[x]`$ given by 2.26, the sum is dominated by the $`n`$’s such that
$$|nNx_c|=𝒪(1)$$
(2.39)
Thus we can still use a quadratic approximation for $`F_0[x]`$
$$Z[V;N]\mathrm{e}^{\left(N^2F_0[V;x_c]+F_1[V;x_c]+\mathrm{}\right)}\underset{n}{}\mathrm{e}^{(nNx_c)^2F_0^{\prime \prime }[V;x_c]/2}$$
(2.40)
where $`F_0^{\prime \prime }=^2F_0/x^2`$ and where the $`\mathrm{}`$ represent terms of order $`𝒪(N^2)`$. The last sum over $`n`$ gives simply an elliptic Jacobi theta function $`\theta _3`$
$$\underset{n}{}\mathrm{e}^{(nNx_c)^2F_0^{\prime \prime }[V;x_c]/2}=\left(2\pi F_0^{\prime \prime }[V;x_c]\right)^{1/2}\theta _3(Nx_c|\tau )$$
(2.41)
with modular parameter $`\tau `$ given by
$$\tau =\frac{2\mathrm{i}\pi }{F_0^{\prime \prime }[V;x_c]}$$
(2.42)
and where the theta function is defined as
$$\theta _3(z|\tau )=\theta _3(z)=\underset{n}{}q^{n^2}\mathrm{e}^{2\mathrm{i}\pi nz}\mathrm{with}q=\mathrm{e}^{\mathrm{i}\pi \tau }$$
(2.43)
It obeys the periodicity relations
$$\theta _3(z+1)=\theta _3(z),\theta _3(z+\tau )=\mathrm{e}^{\mathrm{i}\pi (2z+\tau )}\theta _3(z)$$
(2.44)
(For details on elliptic functions see e.g. Refs ). Eventually we have for the free energy
$`\begin{array}{ccc}\hfill F[V;N]=& N^2& F_0[V,x_c]\hfill \\ & & \mathrm{ln}\left(\theta _3(Nx_c)\right)+F_1(V;x_c)+\frac{1}{2}\mathrm{ln}\left(2\pi F_0^{\prime \prime }[V;x_c]\right)\hfill \\ & +& 𝒪(N^2)\hfill \end{array}`$ (2.48)
where $`F_1`$ is the torus contribution in the topological expansion of 2.33. The next terms of this expansion can be calculated along the same line.
Let us stress that this is not a topological expansion, since the second term $`\mathrm{ln}\left(\theta _3(Nx_c)\right)`$, seemingly $`𝒪(1)`$ and contributing at the torus order, is not regular in $`N`$. Indeed, it is periodic in $`x_c`$ with period $`1/N`$. When computing some observables or quantities of the matrix model, one must take derivatives of $`F`$ w.r.t. some parameters of the potential $`V`$. Since the saddle point $`x_c`$ depends implicitly on $`V`$, every derivative will give a factor $`N`$, and this term may become of the same order than the first term $`N^2F_0[x_c]`$ given by the planar limit. Note that the last two terms depend on $`x`$ and not on $`Nx`$, and they will remain subdominant once we take derivatives of $`F`$.
### 2.3 The modular parameter
Finally, we can express simply the modular parameter $`\tau `$ defined by Eq. (2.42) in term of the end-points $`a`$, $`b`$, $`c`$, $`d`$ of the support of e.v. For this purpose, we introduce the function $`\sigma `$
$$\sigma (\lambda )=(\lambda a)(\lambda b)(\lambda c)(\lambda d)$$
(2.49)
and the function $`u`$
$$u(\lambda )=\frac{1}{2K}_d^\lambda \frac{\text{d}z}{\sqrt{\sigma (z)}}$$
(2.50)
where $`K`$ is
$$K=_b^c\frac{\text{d}z}{\sqrt{|\sigma (z)|}}=\frac{2}{\sqrt{(ca)(db)}}K[m]\mathrm{with}m=\frac{(da)(cb)}{(db)(ca)}$$
(2.51)
$`K[m]`$ is the complete elliptic integral of the first kind. Similarly we define
$$K^{}=_a^b\frac{\text{d}z}{\sqrt{|\sigma (z)|}}=\frac{2}{\sqrt{(ca)(db)}}K[m^{}]\mathrm{with}m^{}=\mathrm{\hspace{0.17em}1}m$$
(2.52)
We shall show that the modular parameter $`\tau `$ of Eq. (2.42) coincides with the standard modular parameter of the torus associated to the mapping $`u`$, i.e. of the elliptic curve $`y^2=\sigma (z)`$. Indeed, $`\tau `$ is simply given by
$$\tau =\mathrm{i}\frac{K^{}}{K}=\mathrm{i}\frac{K[1m]}{K[m]}$$
(2.53)
So we have
$$u(d)=0,u(a)=\frac{1}{2},u(b)=\frac{1+\tau }{2},u(c)=\frac{\tau }{2},u(\mathrm{})=u_{\mathrm{}}$$
(2.54)
and $`u`$ maps the upper half $`\lambda `$-plane onto the half-periods rectangle $`(1/2,\tau /2)`$ and the double-sheeted complex $`\lambda `$-plane onto the period rectangle $`(1,\tau )`$.
To show Eq. (2.53), we use the fact that in the two-cut case, if we fix $`x`$ (the e.v. ratio in the first cut) the semiclassical e.v. density (extrema of the effective action $`\overline{S}`$ ) is now a function $`\rho (\lambda ,x)`$ of $`\lambda `$ and $`x`$, and the end-points $`a`$, $`b`$, $`c`$, $`d`$ depend on $`x`$. Therefore the large $`N`$ resolvent $`\omega _0`$ is of the form
$$\omega _0(\lambda ,x)=\frac{V^{}(\lambda )}{2}\frac{M(\lambda ,x)\sqrt{\sigma (\lambda )}}{2}$$
(2.55)
with $`M(\lambda ,x)`$ a polynomial with degree $`P3`$ in $`\lambda `$ ($`P`$ being the degree of $`V`$), entirely fixed by the constraints 2.18 and 2.29. Therefore the partial derivative of $`\omega _0(\lambda ,x)`$ w.r.t. $`x`$ is necessarily of the form
$$\frac{\omega _0(\lambda ,x)}{x}=\frac{C}{\sqrt{\sigma (\lambda )}}$$
(2.56)
with $`C=C(\lambda ,x)`$ a priori a polynomial in $`\lambda `$. Since 2.18 still holds independently of $`x`$ we must have
$$\frac{\omega _0(\lambda ,x)}{x}=𝒪(\lambda ^2)\mathrm{for}\lambda \mathrm{}$$
(2.57)
which implies that $`C(\lambda ,x)`$ is of degree $`0`$ in $`\lambda `$, i.e. is a constant (depending only on $`x`$)
$$C=C(x)$$
(2.58)
This constant can be easily determined by using that
$$x=_a^b\rho (\lambda ,x)\text{d}\lambda =_𝒞^{}\frac{\text{d}\lambda }{2\mathrm{i}\pi }\omega _0(\lambda ,x)$$
(2.59)
with $`𝒞^{}`$ a clockwise contour encircling the interval $`[a,b]`$. Therefore we have
$$\frac{x}{x}=\mathrm{\hspace{0.17em}1}=_𝒞^{}\frac{\text{d}\lambda }{2\mathrm{i}\pi }\frac{C}{\sqrt{\sigma (\lambda )}}=\frac{CK^{}}{\pi }C=\frac{\pi }{K^{}}$$
(2.60)
with $`K^{}`$ the half-period defined in Eq. (2.52). Now we use Eq. (2.26), Eq. (2.27) and the definition of the effective potential $`V_{\mathrm{eff}}`$ of Eq. (2.14) to write the derivative of the free energy w.r.t. $`x`$ as
$$\frac{F_0}{x}=V_{\mathrm{eff}}(b)V_{\mathrm{eff}}(c)=_b^c\text{d}\lambda \left(2\omega _0(\lambda )V^{}(\lambda )\right)$$
(2.61)
Now we take the derivative w.r.t. $`x`$ of this equation and obtain
$$F_0^{\prime \prime }=\frac{^2F_0}{x^2}=\mathrm{\hspace{0.17em}2}_b^c\text{d}\lambda \frac{\omega _0(\lambda )}{x}=\mathrm{\hspace{0.17em}2}_b^c\text{d}\lambda \frac{C}{\sqrt{\sigma (\lambda )}}=\mathrm{\hspace{0.17em}2}CK=\frac{2\pi K}{K^{}}$$
(2.62)
Using Eq. (2.42) we thus obtain the result 2.53.
## 3 2-points correlation function
### 3.1 The basic formula
As a first application we compute the large $`N`$ smoothed connected two-point correlation function (first obtained by ), defined as
$$\omega ^c(\lambda ,\mu )=\mathrm{Tr}\left[\frac{1}{\lambda M}\right]\mathrm{Tr}\left[\frac{1}{\mu M}\right]\mathrm{Tr}\left[\frac{1}{\lambda M}\right]\mathrm{Tr}\left[\frac{1}{\mu M}\right]$$
(3.1)
Adding source terms to the potential of the form
$$V_{ϵ_\lambda }=V(z)ϵ_\lambda \frac{1}{\lambda z},V_{ϵ_\lambda ,ϵ_\mu }(z)=V(z)ϵ_\lambda \frac{1}{\lambda z}ϵ_\mu \frac{1}{\mu z}$$
(3.2)
we have
$$\omega ^c(\lambda ,\mu )=\frac{1}{N^2}\frac{}{ϵ_\lambda }\frac{}{ϵ_\mu }F[V_{ϵ_\lambda ,ϵ_\mu },N]|_{ϵ_\lambda =ϵ_\mu =0}$$
(3.3)
and for the resolvent (one-point function)
$$\omega (\lambda )=\frac{1}{N}\mathrm{Tr}\left[\frac{1}{\lambda M}\right]=\frac{1}{N^2}\frac{}{ϵ_\lambda }F[V_{ϵ_\lambda },N]|_{ϵ_\lambda =0}$$
(3.4)
If the $`ϵ_\lambda `$’s are small and the $`\lambda `$’s not too close to the cuts, the mean-field solution is still a two-cut e.v. distribution, with $`x_c=x_c(ϵ)`$ an explicit function of the $`\lambda `$’s. So from Eq. (2.48) for the free energy we have for the two-point function
$$\omega ^c(\lambda ,\mu )=\left[\frac{}{ϵ_\lambda }\frac{}{ϵ_\mu }F_0[V_{ϵ_\lambda ,ϵ_\mu }]+\frac{x_c}{ϵ_\lambda }\frac{x_c}{ϵ_\mu }\left[\mathrm{ln}\left(\theta _3(Nx_c)\right)\right]^{\prime \prime }\right]_{ϵ_\lambda =ϵ_\mu =0}+𝒪(N^1)$$
(3.5)
The first term in the r.h.s. of 3.5 is the mean-field contribution already calculated in , the second term involving a second derivative of an elliptic function, characterizes the multi-cut solution.
### 3.2 The mean-field contribution
For completeness let us first rederive the mean field contribution of . Taking the derivative with respect to $`ϵ_\lambda `$ we obtain the mean-field resolvent for the potential $`V_{ϵ_\mu }`$
$$\frac{}{ϵ_\lambda }F_0[V_{ϵ_\lambda ,ϵ_\mu }]|_{ϵ_\lambda =0}=\omega _0(\lambda ;V_{ϵ_\mu })$$
(3.6)
which must be of the form
$$\omega _0(z;V_{ϵ_\mu })=\frac{1}{2}\left[V^{}(z)\frac{ϵ_\mu }{(z\mu )^2}+\frac{M(z)\sqrt{\sigma (z)}}{(z\mu )^2}\right]$$
(3.7)
with $`M(z)`$ a polynomial of degree $`P1`$ (here both the coefficients of $`M`$ and of $`\sigma `$ depend on $`\mu `$ and $`ϵ_\mu )`$. In addition to the $`P2`$ constraints (coming from 2.18 and 2.30) $`\omega _0(z;V_{ϵ_\mu })`$ must be regular at $`z=\mu `$. This determines entirely $`M`$. Taking the derivative w.r.t. $`ϵ_\mu `$ and using the symmetry $`\lambda \mu `$ we get
$$\omega _0^c(\lambda ,\mu )=\frac{}{ϵ_\lambda }\frac{}{ϵ_\mu }F_0=\frac{1}{2}\frac{1}{(\lambda \mu )^2}\left[1+\frac{Q(\lambda ,\mu )}{\sqrt{\sigma (\lambda )}\sqrt{\sigma (\mu )}}\right]$$
(3.8)
with $`Q(\lambda ,\mu )`$ a symmetric polynomial in $`\lambda `$ and in $`\mu `$. The constraints on $`\omega _0^c`$ are: (i) $`\omega _0^c=𝒪(\lambda ^2)`$ as $`\lambda \mathrm{}`$ which implies that $`Q`$ is of degree at most 2; (ii) $`\omega _0^c`$ is regular at $`\lambda =\mu `$ which implies that $`Q(\lambda ,\mu )=\sigma ((\lambda +\mu )/2)+𝒪((\lambda \mu )^2)`$; (iii) finally the equality of the effective potential on the two cuts implies that
$$_b^c\omega _0^c(\lambda ,\mu )\text{d}\mu =\mathrm{\hspace{0.17em}0}$$
(3.9)
Conditions (i) and (ii) fix uniquely $`Q`$
$$Q(\lambda ,\mu )=\frac{1}{2}\left[\begin{array}{c}(\lambda a)(\mu b)(\mu c)(\lambda d)\\ +(\mu a)(\lambda b)(\lambda c)(\mu d)\end{array}\right]+S(\lambda \mu )^2$$
(3.10)
up to a constant $`S`$ fixed by condition (iii), which is found to be
$$S=\frac{1}{2}(ca)(db)\frac{E[m]}{K[m]}\mathrm{with}m=\frac{(cb)(da)}{(ca)(db)}$$
(3.11)
and where $`K[m]`$ and $`E[m]`$ are the standard elliptic integrals of the first and second kind.
### 3.3 The non-regular contribution
In order to compute the second contribution to 3.5, we simply need $`\frac{x_c}{ϵ_\lambda }`$. Since $`x_c`$ is fixed by the constraint $`\frac{F_0}{x}=0`$ we can write
$$\frac{x_c}{ϵ_\lambda }=\frac{^2F_0}{xϵ_\lambda }/\frac{^2F_0}{x^2}$$
(3.12)
Using the results of subsection 2.3 we have
$$\frac{^2F_0}{x^2}=\frac{2\pi K}{K^{}}\mathrm{and}\frac{^2F_0}{xϵ_\lambda }=\frac{}{x}\omega _0(\lambda ,x)=\frac{\pi }{K^{}}\frac{1}{\sqrt{\sigma (\lambda )}}$$
(3.13)
So we have eventually
$$\frac{x_c}{ϵ_\lambda }=\frac{1}{2K\sqrt{\sigma (\lambda )}}$$
(3.14)
and the second non-regular term is
$$\frac{x_c}{ϵ_\lambda }\frac{x_c}{ϵ_\mu }\left[\mathrm{ln}\left(\theta _3(Nx_c|\tau )\right)\right]^{\prime \prime }=\frac{1}{2K\sqrt{\sigma (\lambda )}}\frac{1}{2K\sqrt{\sigma (\mu )}}\left[\mathrm{ln}\left(\theta _3(Nx_c|\tau )\right)\right]^{\prime \prime }$$
(3.15)
with $`K`$ defined by Eq. (2.51). Using standard relations on elliptic functions, this can be rewritten as
$$\frac{(ca)(db)}{4\sqrt{\sigma (\lambda )}\sqrt{\sigma (\mu )}}\left[\frac{E[m]}{K[m]}+\mathrm{dn}^2(Nx_c+{\scriptscriptstyle \frac{1}{2}})\right]$$
(3.16)
with
$$\mathrm{dn}(u)=\mathrm{dn}(2K[m]u|m)$$
(3.17)
where $`\mathrm{dn}(u|m)`$ is the Jacobi elliptic function $`\mathrm{dn}`$. Its periods are $`2K[m]`$ and $`4\mathrm{i}K[m^{}]`$, and $`\mathrm{dn}^2(z)`$ has periods $`1`$ and $`\tau `$.
### 3.4 The final result
Combining 3.8, 3.10, 3.11 and 3.16 we obtain the final result for the 2-point correlation function
$`\omega ^c(\lambda ,\mu )`$ $`=`$ $`{\displaystyle \frac{1}{4(\lambda \mu )^2}}[(1\sqrt{{\displaystyle \frac{(\lambda a)(\lambda b)(\mu c)(\mu d)}{(\mu a)(\mu b)(\lambda c)(\lambda d)}}})+(\lambda \mu )]`$ (3.18)
$`{\displaystyle \frac{(ca)(db)}{4\sqrt{\sigma (\lambda )}\sqrt{\sigma (\mu )}}}\mathrm{sn}^2(Nx_c+{\scriptscriptstyle \frac{1}{2}})`$
We have used the relation $`\mathrm{dn}^2(u)=1m\mathrm{sn}^2(u)`$, where similarly to 3.17 we note
$$\mathrm{sn}(u)=\mathrm{sn}(2K[m]u|m)$$
(3.19)
Surprisingly, the ratio $`E[m]/K[m]`$ characteristic of the mean-field solution of has disappeared.
The smoothed 2-point connected density correlator $`\rho ^c(\lambda ,\mu )`$, defined as
$$\rho ^c(\lambda ,\mu )=\mathrm{Tr}\left[\delta (\lambda M)\right]\mathrm{Tr}\left[\delta (\mu M)\right]\mathrm{Tr}\left[\delta (\lambda M)\right]\mathrm{Tr}\left[\delta (\mu M)\right]$$
(3.20)
can be obtained easily from the discontinuity of $`\omega ^c(\lambda ,\mu )`$. One obtains in the large $`N`$ limit, if $`\lambda `$ and $`\mu `$ are on the support of e.v.
$`\rho ^c(\lambda ,\mu )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}[{\displaystyle \frac{1}{(\lambda \mu )^2}}(\sqrt{\left|{\displaystyle \frac{(\lambda a)(\lambda b)(\mu c)(\mu d)}{(\mu a)(\mu b)(\lambda c)(\lambda d)}}\right|}+\lambda \mu )`$ (3.21)
$`+\epsilon _\lambda \epsilon _\mu {\displaystyle \frac{(ca)(db)}{\sqrt{|\sigma (\lambda )|}\sqrt{|\sigma (\mu )|}}}\mathrm{sn}^2(Nx_c+{\scriptscriptstyle \frac{1}{2}})]`$
$$\epsilon _\lambda =\mathrm{\hspace{0.17em}1}\mathrm{if}\lambda [\mathrm{a},\mathrm{b}],1\mathrm{if}\lambda [\mathrm{c},\mathrm{d}],$$
(3.22)
and zero otherwise.
The new non-regular term $`\mathrm{sn}^2(Nx_c+{\scriptscriptstyle \frac{1}{2}})`$ is an even periodic function of $`Nx_c`$ with period $`1`$ which varies between $`0`$ and $`1`$. Therefore, as $`N`$ varies, depending on the rationality or the irrationality of $`x_c`$, the two-point function will be varying with $`N`$ in a periodic or quasiperiodic way.
### 3.5 The symmetric case
It is now very easy to recover the results of for a symmetric potential. Indeed, if the potential $`V`$ is symmetric, the two cuts are also symmetric
$$a=d,b=c$$
(3.23)
and we have automatically
$$x_c=\frac{1}{2}$$
(3.24)
so that
$$\mathrm{sn}^2(Nx_c+{\scriptscriptstyle \frac{1}{2}})=\{\begin{array}{ccc}\mathrm{sn}^2({\scriptscriptstyle \frac{1}{2}})=1\text{ if }N\text{ is even }\hfill & & \\ \mathrm{sn}^2(0)=0\text{ if }N\text{ is odd}\hfill & & \end{array}$$
(3.25)
Eq. (3.18) and Eq. (3.22) become
$$\omega ^c(\lambda ,\mu )=\frac{1}{2(\lambda \mu )^2}\left[1\frac{(a^2\lambda \mu )(b^2\lambda \mu )}{\sqrt{\sigma (\lambda )}\sqrt{\sigma (\mu )}}\right]\frac{(1)^N}{2}\frac{ab}{\sqrt{\sigma (\lambda )}\sqrt{\sigma (\mu )}}$$
(3.26)
$$\rho ^c(\lambda ,\mu )=\frac{1}{2\pi ^2}\frac{\epsilon _\lambda \epsilon _\mu }{\sqrt{|\sigma (\lambda )|}\sqrt{|\sigma (\mu )|}}\left(\frac{(a^2\lambda \mu )(b^2\lambda \mu )}{(\lambda \mu )^2}(1)^Nab\right)$$
(3.27)
with $`\sigma (\lambda )=(\lambda ^2a^2)(\lambda ^2b^2)`$.
### 3.6 The two-point function as an elliptic function
It is interesting to consider the two-point correlator in terms of the elliptic coordinates defined by Eq. (2.50)
$$u=u(\lambda ),v=u(\mu )$$
(3.28)
Let us thus consider
$$\overline{\omega }^c(u,v)=\frac{\lambda }{u}\frac{\mu }{v}\omega ^c(\lambda ,\mu )=\mathrm{\hspace{0.17em}2}K\sqrt{\sigma (\lambda )}\mathrm{\hspace{0.17em}2}K\sqrt{\sigma (\mu )}\omega ^c(\lambda ,\mu )$$
(3.29)
It is easy to see (from the properties of $`\omega ^c`$) that $`\overline{\omega }^c(u,v)`$ satisfies:
1. $`\overline{\omega }^c(u,v)`$ is a doubly periodic function of $`u`$ (and of $`v`$) with periods $`1`$ and $`\tau `$;
2. $`\overline{\omega }^c(u,v)`$ is regular at $`u`$ and $`v=u(a),u(b),u(c),u(d)`$ and $`u_{\mathrm{}}`$;
3. $`\overline{\omega }^c(u,v)`$ is regular when $`u=v`$, but has a double pole at $`u=v`$ (corresponding to the double pole of $`\omega ^c(\lambda ,\mu )`$ when $`\lambda =\mu `$ but with $`\lambda `$ in the first sheet and $`\mu `$ in the second sheet), with residue $`1`$.
This implies that $`\overline{\omega }^c(u,v)`$ is a Weierstrass elliptic function
$$\overline{\omega }^c(u,v)=\mathrm{}(u+v|\tau )+\mathrm{constant}$$
(3.30)
where the constant depends on $`Nx_c`$ ($`\mathrm{}`$ has periods $`1`$ and $`\tau `$). Using classical identities between the Weirstrass $`\mathrm{}`$ function and the Jacobi elliptic functions, it can be easily calculated. We find the remarkably simple result
$$\overline{\omega }^c(u,v)=\mathrm{}(u+v|\tau )\mathrm{}(Nx_c+{\scriptscriptstyle \frac{\tau }{2}}|\tau )$$
(3.31)
or equivalently
$$\overline{\omega }^c(u,v)=\left[\mathrm{ln}\left(\theta _1(u+v|\tau )\right)\right]^{\prime \prime }+\left[\mathrm{ln}\left(\theta _3(Nx_c|\tau )\right)\right]^{\prime \prime }$$
(3.32)
## 4 The orthogonal polynomials
Let us briefly recall some basic facts about the well-known method of orthogonal polynomials , which is a powerful tool for studying the spectral properties of random matrices . Asymptotic expressions for the orthogonal polynomials have been obtained recently in the mathematical literature, by solving a Rieman-Hilbert problem. Here we will derive them from the free energy directly.
Consider the partition function (2.4):
$$\stackrel{~}{Z}=\text{d}\lambda _1\mathrm{}\text{d}\lambda _N\mathrm{e}^{N_iV(\lambda _i)}\underset{i<j}{}(\lambda _i\lambda _j)^2$$
(4.1)
The last term is a Vandermonde determinant :
$$\underset{i<j}{}(\lambda _i\lambda _j)=\underset{i,j}{det}\left((\lambda _i)^{j1}\right)=\underset{i,j}{det}\left(𝒫_{j1}(\lambda _i)\right)$$
(4.2)
where the last equality is obtained by linearly mixing columns of the determinant, and holds for arbitrary monic polynomials $`𝒫_n(\lambda )`$ with leading coefficient $`𝒫_n(\lambda )=\lambda ^n+\mathrm{}`$.
The method of orthogonal polynomials consists in choosing a family of polynomials suitable for the computation of (4.1), namely, the family of polynomials orthogonal with respect to the weight $`\mathrm{exp}NV(\lambda )`$:
$$\text{d}\lambda 𝒫_n(\lambda )𝒫_m(\lambda )\mathrm{e}^{NV(\lambda )}=h_n\delta _{nm}$$
(4.3)
With this particular choice of polynomials, the integral (4.1) is merely:
$$\stackrel{~}{Z}=N!\underset{n=0}{\overset{N1}{}}h_n$$
(4.4)
and the joint probability density of all the eigenvalues takes the form of a Slater determinant:
$$R_N(\lambda _1,\mathrm{},\lambda _N)=\frac{1}{N!}\left(\underset{0n<N}{det}_{1iN}\left[\psi _{n1}(\lambda _i)\right]\right)^2$$
(4.5)
where the wave functions $`\psi _n(\lambda )=\frac{1}{\sqrt{h_n}}𝒫_n(\lambda )\mathrm{e}^{\frac{N}{2}V(\lambda )}`$ are orthonormal.
### 4.1 The Kernel $`K(\lambda ,\mu )`$
The square of a determinant can be rewritten as the determinant of a product:
$$\left(\underset{n,i}{det}\left(\psi _{n1}(\lambda _i)\right)\right)^2=\underset{1i,jN}{det}\left[\underset{n=0}{\overset{N1}{}}\psi _n(\lambda _i)\psi _n(\lambda _j)\right]$$
we are thus led to introduce the kernel $`K(\lambda ,\mu )`$ :
$$K(\lambda ,\mu )=\frac{1}{N}\underset{n=0}{\overset{N1}{}}\psi _n(\lambda )\psi _n(\mu )$$
(4.6)
In terms of which the joint density of eigenvalues is now a determinant
$$R_N(\lambda _1,\mathrm{},\lambda _N)=\frac{N^N}{N!}det\left[K(\lambda _i,\lambda _j)\right]$$
(4.7)
The orthonormality properties of the polynomials imply the projection relations
$$\text{d}\lambda K(\lambda ,\lambda )=1\mathrm{and}\text{d}\lambda K(\mu ,\lambda )K(\lambda ,\nu )=\frac{1}{N}K(\mu ,\nu )$$
(4.8)
which make any partial integration of (4.7) easy to perform (theorem of Dyson ).
In particular, the integration over $`N1`$ eigenvalues gives the density of eigenvalues
$$\rho (\lambda _1)=\text{d}\lambda _2\mathrm{}\text{d}\lambda _NR_N(\lambda _1,\mathrm{},\lambda _N)=K(\lambda _1,\lambda _1)$$
and the integration over $`N2`$ eigenvalues gives the correlation function:
$$\begin{array}{cc}\hfill R_2(\lambda _1,\lambda _2)& =\text{d}\lambda _3\mathrm{}\text{d}\lambda _NR_N(\lambda _1,\mathrm{},\lambda _N)\hfill \\ & =\frac{N}{N1}\left(K(\lambda _1,\lambda _1)K(\lambda _2,\lambda _2)K(\lambda _1,\lambda _2)K(\lambda _2,\lambda _1)\right)\hfill \end{array}$$
In short:
$$\rho (\lambda )=K(\lambda ,\lambda ),\rho (\lambda ,\mu )=(K(\lambda ,\lambda )K(\mu ,\mu )K(\lambda ,\mu )^2)$$
(4.9)
In addition, the Darboux-Christoffel theorem , asserts that
$$K(\lambda ,\mu )=\frac{1}{Nh_{N1}}\frac{𝒫_N(\lambda )𝒫_{N1}(\mu )𝒫_N(\mu )𝒫_{N1}(\lambda )}{\lambda \mu }\mathrm{e}^{\frac{N}{2}(V(\lambda )+V(\mu ))}$$
(4.10)
which means that we need to evaluate $`𝒫_n`$ only for $`n=N`$ and $`n=N1`$.
Thus, we shall now aim at finding asymptotic expressions for the orthogonal polynomials $`𝒫_n(\lambda )`$, and the kernel $`K(\lambda ,\mu )`$ in the large $`N`$ limit, and $`n`$ close to $`N`$. This has been done in the 1-cut case and in the symmetric 2-cut case . Here we will generalize it to the non-symmetric case, with the method used in .
### 4.2 WKB approximation for the orthogonal polynomials $`𝒫_n(\lambda )`$
The orthogonal polynomials have the following integral representation (see Appendix 1 of or ):
$$𝒫_n(\lambda )=\frac{\text{d}M_{n\times n}det(\lambda M)\mathrm{e}^{N\mathrm{tr}V(M)}}{\text{d}M_{n\times n}\mathrm{e}^{N\mathrm{tr}V(M)}}$$
(4.11)
where the integral is restricted to hermitian matrices of size $`n\times n`$.
Thus the orthogonal polynomial is given by the ratio of two matrix integrals of the same type as the partition function 2.1:
$$𝒫_n(\lambda )=\frac{Z[V+\delta V_1+\delta V_2;n]}{Z[V+\delta V_1;n]}=\frac{\mathrm{e}^{F[V+\delta V_1+\delta V_2;n]}}{\mathrm{e}^{F[V+\delta V_1;n]}}$$
(4.12)
where
$$\delta V_1(z)=\frac{Nn}{n}V(z)\mathrm{and}\delta V_2=\frac{1}{n}\mathrm{ln}(z\lambda )$$
(4.13)
We have seen in the previous section (eq.2.48) that
$$F[V;n]=n^2F_0[V;x_c]\mathrm{ln}\theta _3(nx_c[V])+\mathrm{}$$
(4.14)
We will use the fact that under a variation $`\delta V`$ of the potential, the variation of $`F_0`$ is :
$$\delta F_0=\frac{1}{2i\pi }\omega (z)\delta V(z)\text{d}z$$
(4.15)
where the anti-clockwise contour encloses the support of the density of eigenvalues, and $`\omega (z)`$ is the resolvent (eq 3.4 and 2.28):
$$\omega (z)=\frac{1}{2}\left(V^{}(z)M(z)\sqrt{\sigma (z)}\right)$$
It is convenient to introduce two sources $`t_1`$ and $`t_2`$ for the variations $`\delta V_1`$ and $`\delta V_2`$ of the potential, and consider a generalized potential $`𝒱(z)`$:
$$𝒱(z)=V(z)+t_1\delta V_1(z)+t_2\delta V_2(z)=V(z)+t_1V(z)+t_2\mathrm{ln}|\lambda z|$$
Since $`t_1=\frac{Nn}{n}`$ and $`t_2=\frac{1}{n}`$ are both small of order $`O(N^1)`$, we will expand $`F`$ in Taylor’s series:
$$F[𝒱(z);n]=F[V(z);n]+t_1_1F+t_2_2F+\frac{t_1^2}{2}_{11}F+t_1t_2_{12}F+\frac{t_2^2}{2}_{22}F+\mathrm{}$$
all the derivatives being taken at the point $`t_1=t_2=0`$.
This will give
$$𝒫_n(\lambda )\mathrm{e}^{n_2F_0}\mathrm{e}^{(Nn)_{12}F_0}\mathrm{e}^{\frac{1}{2}_{22}F_0}\frac{\theta _3(nx+(Nn)_1x_2x)}{\theta _3(nx+(Nn)_1x)}(1+O(N^1))$$
(4.16)
Now, let us compute the derivatives of $`F_0`$ and $`x=x_c`$ with respect to $`t_1`$ and $`t_2`$. The method proceeds similarly to section 3.1.
#### 4.2.1 Derivatives of $`F_0`$ with respect to $`t_1`$ and $`t_2`$
using (4.15) with (4.13):
$$\frac{F_0}{t_2}=\frac{1}{2i\pi }\omega (z)\mathrm{ln}(z\lambda )\text{d}z$$
After integration by part, the pole in $`(z\lambda )`$ picks a residue, and the result is a primitive of $`\omega (\lambda )`$:
$$\frac{F_0}{t_2}=_{\lambda _0}^\lambda \omega (z)\text{d}z$$
(4.17)
The lower bound of integration $`\lambda _0`$ is to be chosen such that $`\mathrm{e}^{n_2F_0}_\lambda \mathrm{}\lambda ^n`$. i.e.
$$\mathrm{ln}\lambda _0=_{\lambda _0}^{\mathrm{}}(\omega (z)\frac{1}{z})\text{d}z$$
(4.18)
In order to compute the second derivatives $`_{12}F_0`$ and $`_{22}F_0`$, we will need to differentiate $`\omega (z)`$ with respect to $`t_1`$ and $`t_2`$.
#### 4.2.2 Derivatives of $`\omega (z)`$ with respect to $`t_1`$ and $`t_2`$
The resolvent $`\omega (z)`$ computed for the potential $`𝒱(z)`$ takes the form:
$$\omega (z)=\frac{1}{2}\left(𝒱^{}(z)M(z)\sqrt{\sigma (z)}\right)$$
(4.19)
where $`M(z)`$ is analytic. Notice that when $`𝒱^{}(z)`$ has a pole in $`z=\lambda `$, $`M(z)`$ may have a pole too. $`\omega (z)`$ obeys a linear equation:
$$\omega (z+i0)+\omega (zi0)=𝒱^{}(z)\mathrm{for}z[a,b][c,d]$$
(4.20)
Thus its derivatives obey linear equations as well:
$$_1\omega (z+i0)+_1\omega (zi0)=\delta V_1^{}(z)=V^{}(z)$$
(4.21)
$$_2\omega (z+i0)+_2\omega (zi0)=\delta V_2^{}(z)=\frac{1}{z\lambda }$$
(4.22)
$``$ $`\omega _1`$: The solution of (4.21) is:
$$_1\omega (z)=\omega (z)\frac{f(z)}{\sqrt{\sigma (z)}}$$
(4.23)
where $`f(z)`$ is analytic in $`z`$. The boundary conditions 2.18 imply that $`f(z)z`$ when $`z\mathrm{}`$ and $`f`$ has no pole, thus $`f(z)`$ is a polynomial of degree $`1`$:
$$\frac{\omega (z)}{t_1}=\omega (z)\frac{zz_0}{\sqrt{\sigma (z)}}$$
(4.24)
$`z_0`$ is determined as a function of $`a,b,c,d`$ by the derivative of 2.30 with respect to $`t_1`$:
$$_b^c\text{d}z\frac{zz_0}{\sqrt{\sigma (z)}}=0$$
(4.25)
It can be checked that in term of elliptic theta functions we have (see Appendix A, or ):
$$\frac{zz_0}{\sqrt{\sigma (z)}}=\frac{\text{d}}{\text{d}z}\mathrm{ln}\frac{\theta _1(u(z)+u_{\mathrm{}})}{\theta _1(u(z)u_{\mathrm{}})}$$
(4.26)
and thus:
$$\frac{\omega (z)}{t_1}=\omega (z)\frac{\text{d}}{\text{d}z}\mathrm{ln}\frac{\theta _1(u(z)+u_{\mathrm{}})}{\theta _1(u(z)u_{\mathrm{}})}$$
(4.27)
$``$ $`\omega _2`$: Note that the $`t_2`$ source-term is the primitive of the $`ϵ_\lambda `$ source-term of 3.2, and that
$$\frac{\text{d}}{\text{d}z}\frac{\omega (z)}{t_2}=\frac{\omega (z)}{ϵ_\lambda }=\frac{1}{n^2}\frac{^2F}{ϵ_zϵ_\lambda }=\omega _c(z,\lambda )$$
(4.28)
so, $`\omega /t_2`$ has already been computed in 3.8. The second derivative $`_{22}F`$ corresponds to $`z=\lambda `$.
$$_{22}F_0=\mathrm{ln}\sqrt{\sigma (\lambda )}+2\mathrm{ln}(\theta _1(u(\lambda )u_{\mathrm{}}))$$
(4.29)
#### 4.2.3 Derivatives of $`x`$
Recall that
$$x=_a^b\rho (\lambda )\text{d}\lambda =\frac{1}{2i\pi }_a^bM(z)\sqrt{\sigma (z)}\text{d}z$$
(4.30)
using 4.27 we get:
$$\frac{x}{t_1}=x+\frac{1}{2i\pi }_a^b\frac{zz_0}{\sqrt{\sigma (z)}}\text{d}z=x+2u_{\mathrm{}}$$
(4.31)
Similarly, from 4.28 and 3.18 or 3.31 (or less tediously, taking the primitive of 3.14), we get:
$$\frac{x}{t_2}=u(\lambda )+u_{\mathrm{}}$$
(4.32)
### 4.3 Final result
#### 4.3.1 Case $`\lambda /[a,b][c,d]`$
Eventually, inserting 4.17,4.27,4.29,4.31, 4.32 into 4.16 we get:
$$𝒫_n(\lambda )\underset{\lambda [a,b][c,d]}{=}\sqrt{u^{}(\lambda )}p_n(u(\lambda ))\mathrm{e}^{N_{\lambda _0}^\lambda \omega }$$
(4.33)
where
$`p_n(u)=C_n{\displaystyle \frac{\theta _3(Nx+2(Nn)u_{\mathrm{}}+uu_{\mathrm{}})\theta _1(2u_{\mathrm{}})}{\theta _3(Nx+2(Nn)u_{\mathrm{}})\theta _1(uu_{\mathrm{}})}}\left({\displaystyle \frac{\theta _1(u+u_{\mathrm{}})}{\theta _1(uu_{\mathrm{}})}}\right)^{nN}`$ (4.34)
$`C_n`$ is a normalization such that $`𝒫_n\lambda ^n`$ for $`\lambda \mathrm{}`$.
$$C_n=\sqrt{2K}A^{nN+1}\mathrm{with}A=\frac{1}{2K}\frac{\theta _1^{}(0)}{\theta _1(2u_{\mathrm{}})}\mathrm{and}K=_c^b\frac{\text{d}z}{\sqrt{\sigma (z)}}$$
(4.35)
$$A=\frac{1}{4}|dac+b|\frac{\theta _3(0)}{\theta _3(2u_{\mathrm{}})}$$
Note that 4.33 is unchanged under $`uu+1`$ and $`uu+\tau `$. Indeed, a shift $`uu+\tau `$ amounts to a nontrivial circle around the cut $`[c,d]`$. Thus $`\omega `$ is shifted by $`2i\pi _c^d\rho =2i\pi (1x)`$, and $`\mathrm{e}^{N{\scriptscriptstyle \omega }}`$ receives a phase $`\mathrm{e}^{2i\pi Nx}`$. In the same time, the $`\theta `$ functions receive phase factors: $`\theta (v+\tau )=\theta (v)\mathrm{e}^{2i\pi (v+\tau /2)}`$. One can easily check that the total phase shift is $`0`$.
#### 4.3.2 Case $`\lambda [a,b][c,d]`$
Expression 4.34 has been derived by a saddle point approximation of 4.12 when $`\lambda `$ does not belong to $`[a,b][c,d]`$. When $`\lambda `$ lies on the cut $`[a,b][c,d]`$, 4.12 actually has two saddle points, contributing to the same order. They correspond to the two determinations of the square root $`\pm \sqrt{\sigma (\lambda )}`$. The asymptotic expression for the orthogonal polynomial is then given by a sum of two terms:
$$𝒫_n(\lambda )\underset{\lambda [a,b][c,d]}{=}C\sqrt{u^{}}\left[p_n(u)\mathrm{e}^{iN\pi \zeta (\lambda )}+ip_n(u)\mathrm{e}^{iN\pi \zeta (\lambda )}\right]\mathrm{e}^{\frac{N}{2}V(\lambda )}$$
(4.36)
where $`\zeta (\lambda )=_d^\lambda \rho (z)\text{d}z`$ and
$$C=\mathrm{e}^{\frac{N}{2}\left(V(\lambda _0)+_{\lambda _0}^dM(z)\sqrt{\sigma (z)}\text{d}z\right)}=d^N\mathrm{e}^{\frac{N}{2}V(d)}\mathrm{e}^{N_d^{\mathrm{}}(\omega (z)\frac{1}{z})\text{d}z}$$
(4.37)
To summarize: When $`\lambda [a,b][c,d]`$, the wave function $`\psi _n(\lambda )=𝒫_n(\lambda )\mathrm{e}^{NV/2}`$ decays exponentially, and within the support $`[a,b][c,d]`$, it oscillates at a frequency of order $`N`$.
#### 4.3.3 Check of the orthogonality
For completeness, let us check that the functions (4.33) are indeed orthogonal (at leading order in $`N^1`$). Let us compute the integral:
$$_{\mathrm{}}^{\mathrm{}}\text{d}\lambda 𝒫_n(\lambda )𝒫_m(\lambda )\mathrm{e}^{NV(\lambda )}$$
The contributions of the integral along $`]\mathrm{},a][b,c][d,\mathrm{}[`$ are exponentially small and do not contribute at leading order.
Along $`[a,b][c,d]`$, we use expression 4.36, and get a sum of four terms:
$$C^2\text{d}\lambda u^{}(\lambda )\left(\begin{array}{cc}ip_n(u)p_m(u)\hfill & +ip_n(u)p_m(u)\hfill \\ +p_n(u)p_m(u)\mathrm{e}^{2iN\pi \zeta (\lambda )}\hfill & p_n(u)p_m(u)\mathrm{e}^{2iN\pi \zeta (\lambda )}\hfill \end{array}\right)$$
(4.38)
Since the two last terms have fast oscillations of frequency $`N`$, they are suppressed as $`O(1/N)`$.
The leading contribution is thus given by the two first terms of 4.38, which can be rewritten as integrals in the $`u`$ plane along the contour depicted on fig.3.a:
$`\begin{array}{cc}{\displaystyle 𝒫_n𝒫_m\mathrm{e}^{NV}}=\hfill & ic_{nm}{\displaystyle \text{d}u\frac{\theta _3(x_n+uu_{\mathrm{}})\theta _3(x_muu_{\mathrm{}})}{\theta _3(x_n)\theta _3(x_m)\theta _1(uu_{\mathrm{}})\theta _1(u+u_{\mathrm{}})}\left(\frac{\theta _1(u+u_{\mathrm{}})}{\theta _1(uu_{\mathrm{}})}\right)^{nm}}\hfill \\ & +(uu)\hfill \end{array}`$ (4.41)
where $`x_n`$ and $`c_{nm}`$ are short notations for:
$$x_n=Nx+2(Nn)u_{\mathrm{}}\mathrm{and}c_{nm}=C^2C_nC_m\theta _1^2(2u_{\mathrm{}})$$
(4.42)
If $`n>m`$ we may deform the contour to a circle around the point $`u_{\mathrm{}}`$ (fig.3.b), and the integral vanishes since there is no pole, while if $`m>n`$ we deform the contour to a circle around $`+u_{\mathrm{}}`$ (fig.3.c). Therefore, the integral vanishes for $`nm`$.
When $`n=m`$, the integral picks a residue:
$$\text{d}\lambda 𝒫_n𝒫_m\mathrm{e}^{NV(\lambda )}=h_n\delta _{nm}$$
(4.43)
with ($`C`$, $`c_{nm}`$, $`A`$, $`x_n`$ are defined in 4.35, 4.42, 4.37):
$$h_n=c_{nn}\frac{4\pi }{\theta _1(2u_{\mathrm{}})\theta _1^{}(0)}\frac{\theta _3(x_{n+1})}{\theta _3(x_n)}=4\pi C^2\frac{\theta _3(x_{n+1})}{\theta _3(x_n)}A^{2(nN+1/2)}$$
(4.44)
#### 4.3.4 Recurrence equation
It is well known that the orthogonal polynomials satisfy a recurrence equation of the form :
$$\lambda 𝒫_n(\lambda )=𝒫_{n+1}(\lambda )+\beta _n𝒫_n(\lambda )+\alpha _n𝒫_{n1}(\lambda )$$
(4.45)
Here, we find that (divide 4.45 by $`𝒫_n`$, and match the poles on both sides):
$$\alpha _n=\frac{h_n}{h_{n1}}=A^2\frac{\theta _3(x_{n+1})\theta _3(x_{n1})}{\theta _3^2(x_n)}$$
(4.46)
which can be rewritten more compactly as
$$\alpha _n=\frac{1}{16}\left(((da)(cb))^2+4(da)(cb)\mathrm{cn}^2(x_n+1/2)\right)$$
(4.47)
And by taking $`u=0`$ in 4.45 we get $`\beta _n`$:
$$\beta _nd=A\left[\frac{\theta _3(x_{n+3/2})\theta _3(x_n)}{\theta _3(x_{n+1})\theta _3(x_{n+1/2})}+\frac{\theta _3(x_{n+1})\theta _3(x_{n1/2})}{\theta _3(x_n)\theta _3(x_{n+1/2})}\right]$$
(4.48)
which can be rewritten more compactly as:
$$\beta _n=\frac{a+d+(cb)}{2}(cb)\frac{db}{cb+\frac{dc}{\mathrm{cn}^2(x_nu_{\mathrm{}}+1/2)}}$$
(4.49)
The sequences $`\alpha _n`$ and $`\beta _n`$ are thus quasi-periodic in $`n`$. It is interesting to recall that the behavior of these coefficients has been extensively studied (mainly by numerical methods) by several authors in the early 90’s . The general conclusion was that in the multi-cut case the general behavior of the recursion coefficients was “chaotic” in $`n`$ (and regular or quasi-periodic only in some special cases). It is clear from our expressions that in the two-cut case the behavior is always periodic or quasi-periodic and never chaotic (in the mathematical sense). This is in fact true even if the number of cuts is larger than 2 (see appendix C).
In the symmetric case, $`x=\frac{1}{2}`$ and $`u_{\mathrm{}}=\frac{1}{4}`$, we have $`x_n=n/2\mathrm{mod}\mathrm{\hspace{0.17em}1}`$, so that we recover $`\beta _n=0`$ and $`\alpha _n=\frac{1}{4}(a(1)^nb)^2`$.
In the general case, $`\alpha _n`$ and $`\beta _n`$ vary along a periodic curve, between two extrema, given by:
$$\frac{(da(cb))^2}{16}\alpha _n\frac{(da+(cb))^2}{16}$$
$$\frac{d+a}{2}\frac{cb}{2}\beta _n\frac{d+a}{2}+\frac{cb}{2}$$
Similarly to the one-cut case, one may relate $`\alpha _n`$ to square width of the distribution of eigenvalues, and $`\beta _n`$ to the center of the distribution.
### 4.4 The kernel $`K(\lambda ,\mu )`$
We can now evaluate the kernel $`K(\lambda ,\mu )`$ according to (4.10). Let us note $`u=u(\lambda )`$ and $`v=u(\mu )`$ and we assume $`\lambda ,\mu [a,b][c,d]`$:
$$K(\lambda ,\mu )\frac{C^2\sqrt{u^{}v^{}}}{Nh_{N1}}\underset{ϵ,\eta =\pm 1}{}\frac{\sqrt{ϵ\eta }p_N(ϵu)p_{N1}(\eta v)\mathrm{e}^{ϵNi\pi \zeta (\lambda )}\mathrm{e}^{\eta Ni\pi \zeta (\mu )}(uv)}{(\lambda \mu )}$$
(4.50)
which can be rewritten as a sum of eight terms:
$$\begin{array}{cc}K(\lambda ,\mu )=& \frac{c_{N,N1}}{h_{N1}\theta _3(x_N)\theta _3(x_{N1})}\frac{\sqrt{u^{}v^{}}}{N(\lambda \mu )}\times \\ & \underset{ϵ,\eta ,\kappa =\pm 1}{}\kappa \sqrt{ϵ\eta }\frac{\theta _3(Nx+ϵu\kappa u_{\mathrm{}})\theta _3(Nx+\eta v+\kappa u_{\mathrm{}})}{\theta _1(ϵu\kappa u_{\mathrm{}})\theta _1(\eta v+\kappa u_{\mathrm{}})}\mathrm{e}^{Ni\pi (ϵ\zeta (\lambda )+\eta \zeta (\mu ))}\end{array}$$
(4.51)
We will see below that not all the terms contribute to the same order.
#### 4.4.1 Regime $`|\lambda \mu |O(1/N)`$
The eight terms of (4.51) can be rewritten as four combinations of the type:
$$\mathrm{sin}\left(N\pi (\zeta (\lambda )\pm \zeta (\mu ))\right)\frac{f(u,v)f(v,u)}{N(\lambda \mu )}\text{and}\mathrm{cos}\left(N\pi (\zeta (\lambda )\pm \zeta (\mu ))\right)\frac{g(u,v)g(v,u)}{N(\lambda \mu )}$$
In the limit $`|\lambda \mu |`$ small, i.e. $`|uv|`$ small, the terms with a cosine will be proportional to derivatives of $`g(u,v)`$, and there will be an overall $`\frac{1}{N}`$ factor. Similarly, the term with a sine and a $`+`$ sign will be proportional to a derivative of $`f(u,v)`$ and will be of order $`1/N`$. Only the term proportional to $`\mathrm{sin}N\pi _\lambda ^\mu \rho (z)\text{d}z`$ can balance the $`1/N`$ factor, and is dominant in the short range regime. After calculation we get:
$$K(\lambda ,\mu )\underset{|\lambda \mu |O(1/N)}{}\frac{\mathrm{sin}N\pi _\lambda ^\mu \rho (z)\text{d}z}{N\pi (\lambda \mu )}$$
(4.52)
As expected we have
$$K(\lambda ,\lambda )=\rho (\lambda )$$
(4.53)
and we recover the universal short range correlation function:
$$\rho (\lambda ,\mu )\rho (\lambda )\rho (\mu )\left(1\left(\frac{\mathrm{sin}N\pi \rho (\lambda )(\lambda \mu )}{N\pi \rho (\lambda )(\lambda \mu )}\right)^2\right)$$
(4.54)
#### 4.4.2 Long range regime, smoothed oscillations
When $`|\lambda \mu |O(1)`$, $`K(\lambda ,\mu )`$ has high frequency oscillations, and only a smoothed correlation function obtained by averaging the oscillations can be observed.
Recall that the connected 2-point correlation function is related to $`K^2`$ by (4.9):
$$\rho _{2c}(\lambda ,\mu )=K(\lambda ,\mu )^2$$
(4.55)
with $`K(\lambda ,\mu )`$ given by 4.51.
Smoothing out the oscillations amounts to kill all terms containing some $`\mathrm{e}^{iN\pi \zeta }`$ in the square of eq 4.51, we thus have:
$$\begin{array}{c}\overline{K(\lambda ,\mu )^2}=\frac{2u^{}v^{}c_{NN1}^2}{h_{N1}^2\theta _3^2(x_N)\theta _3^2(x_{N1})N^2(\lambda \mu )^2}\underset{ϵ,\eta ,\kappa _1,\kappa _2=\pm 1}{}\kappa _1\kappa _2\hfill \\ \frac{\theta _3(Nx+ϵu\kappa _1u_{\mathrm{}})\theta _3(Nxϵu\kappa _2u_{\mathrm{}})\theta _3(Nx+\eta v+\kappa _1u_{\mathrm{}})\theta _3(Nx\eta v+\kappa _2u_{\mathrm{}})}{\theta _1(ϵu\kappa _1u_{\mathrm{}})\theta _1(ϵu\kappa _2u_{\mathrm{}})\theta _1(\eta v+\kappa _1u_{\mathrm{}})\theta _1(\eta v+\kappa _2u_{\mathrm{}})}\hfill \end{array}$$
(4.56)
Using that (see appendix A, and ):
$$\lambda \mu =2A\frac{\theta _1(uv)\theta _1(u+v)\theta _1^2(2u_{\mathrm{}})}{\theta _1(uu_{\mathrm{}})\theta _1(u+u_{\mathrm{}})\theta _1(vu_{\mathrm{}})\theta _1(v+u_{\mathrm{}})}$$
(4.57)
we get ($`ϵ_{ij}=ϵ_{ji}=\pm 1`$):
$$\begin{array}{c}\overline{K(\lambda ,\mu )^2}=\frac{2u^{}v^{}\theta _1^2(0)}{N^24\pi ^2\theta _3^4(Nx)\theta _1^2(2u_{\mathrm{}})}\hfill \\ \left(\frac{\theta _1(uu_{\mathrm{}})\theta _1(u+u_{\mathrm{}})\theta _1(vu_{\mathrm{}})\theta _1(v+u_{\mathrm{}})}{\theta _1(uv)\theta _1(u+v)}\right)^2\underset{i,j,k,l=\pm 1}{}(ϵ_{ij}ϵ_{kl}+ϵ_{il}ϵ_{kj})\hfill \\ \frac{\theta _3(Nx+uiu_{\mathrm{}})\theta _3(Nxuku_{\mathrm{}})\theta _3(Nx+vju_{\mathrm{}})\theta _3(Nxvlu_{\mathrm{}})}{\theta _1(uiu_{\mathrm{}})\theta _1(uku_{\mathrm{}})\theta _1(vju_{\mathrm{}})\theta _1(vlu_{\mathrm{}})}\hfill \end{array}$$
(4.58)
We see that 4.58 has no pole when $`u=\pm u_{\mathrm{}}`$, it can have (double) poles only when $`u=\pm v`$. Thus, 4.58 can be rewritten in terms of Weirstrass functions of $`u+v`$ and $`uv`$:
$$\overline{K(\lambda ,\mu )^2}=\frac{1}{2N^2\pi ^2}u^{}v^{}\left(C_1\mathrm{}(uv)+C_2\mathrm{}(u+v)2S\right)$$
Taking $`u=v`$ and $`u=v`$ in 4.58, we find that the residues are $`C_1=C_2=1`$, and taking a particular value of $`u`$ and $`v`$, we find the constant $`S`$, equal to what we had in 3.31:
$$\overline{K(\lambda ,\mu )^2}=\frac{1}{2N^2\pi ^2}u^{}v^{}\left(\mathrm{}(uv)+\mathrm{}(u+v)2\mathrm{}(Nx+\frac{\tau }{2})\right)$$
(4.59)
and we recover the result 3.31 found in section 3.6.
## 5 Conclusions
In this article, we have solved the puzzle raised by and understood why the naive mean-field method and the orthogonal polynomial ansatz approach used in the symmetric case disagree.
We have proven here that this effect has nothing to do with a $`_2`$ symmetry breaking, as it was sometimes assumed , it is general as soon as the support of the density is not connected.
The apparent paradox comes from the fact that when the support of eigenvalues is not connected, the free energy admits no large $`N`$ expansion in powers of $`1/N^2`$ (topological expansion ). This means that the free energy in the multi-cut case is not given by a topological expansion, i.e. the sum of diagrams with a weight $`N^\chi `$ ($`\chi `$=Euler Characteristic of the diagram).
The explanation lies in the discreteness of the number of eigenvalues. For instance in the symmetric 2-cut case, the classical approach assumes that the minimum of the free energy is reached when one half ($`x=1/2`$) of the eigenvalues are in each cut. Obviously, this minimum is never reached when the total number of eigenvalues is odd, and in general, the result depends on the fractional part of $`Nx`$.
At leading order in $`N`$ only, the free energy is correctly given by the classical saddle point limit , but the first order in $`N`$ is not sufficient to determine the 2-point (or higher) correlation function.
Here we have computed explicitly the two-point connected correlation function. It contains a universal part depending only on the number of cuts, which was obtained by , and contains in addition, a non universal term quasiperiodic in $`N`$ .
Let us stress that our calculation holds for any potential, not necessarily symmetric, and it can also be generalized to a potential with complex coefficients (appendix. B), and to an arbitrary number of cuts (appendix. C).
We have also reobtained directly the asymptotic expressions for the orthogonal polynomials , which allows in principle through the Darboux-Christoffel theorem (eq. 4.10) to compute any correlation function of any number of eigenvalues in the short or long range domain (and one can smooth it afterwards).
The orthogonal polynomial approach may in turn be used for other random matrix ensembles, and it would be interesting to apply our results to orthogonal or symplectic ensembles .
The authors are thankful to K. Mallick for useful discussions, and to the Eurogrid European Network HPRN-CT-1999-00161 for supporting part of the work. They also thank E. Kanzieper, O. Lechtenfeld and G. Akemann for their interest and for pointing some missing references.
## Appendix AA few useful identities on elliptic functions
Here we collect a few useful identities on elliptic functions used through the paper. For details see . We start from
$$\sigma (\lambda )=(\lambda a)(\lambda b)(\lambda c)(\lambda d),a<b<c<d$$
(A.1)
and the map from the complex plane to the torus
$$u(\lambda )=\frac{1}{2K}_d^\lambda \frac{\text{d}z}{\sqrt{\sigma (z)}}$$
(A.2)
where the half period $`K`$ is
$$K=_b^c\frac{\text{d}z}{\sqrt{|\sigma (z)|}}=\frac{2}{\sqrt{(ca)(db)}}K[m]=\frac{\pi \theta _3^2(0|\tau )}{\sqrt{(ca)(db)}}$$
(A.3)
$`K[m]`$ is the standard complete elliptic integral , with the modulus $`m`$ equal to the biratio of the four points $`a`$, $`b`$, $`c`$, $`d`$:
$$m=\frac{(da)(cb)}{(db)(ca)}$$
(A.4)
$`m`$ is related to the modular parameter $`\tau `$ of the torus by:
$$m=\mathrm{e}^{i\pi \tau }\frac{\theta _3^4(\frac{\tau }{2}|\tau )}{\theta _3^4(0|\tau )}\mathrm{and}\mathrm{conversely}\tau =\mathrm{i}\frac{K[1m]}{K[m]}$$
(A.5)
where we have used the Jacobi theta functions:
$$\theta _1(z|\tau )=\theta _1(z)=\mathrm{i}\underset{r+1/2}{}(1)^rq^{r^2}\mathrm{e}^{2\mathrm{i}\pi rz}\mathrm{with}q=\mathrm{e}^{\mathrm{i}\pi \tau }$$
(A.6)
$$\mathrm{and}\theta _3(z|\tau )=\theta _3(z)=q^{\frac{1}{4}}\mathrm{e}^{i\pi z}\theta _1(z+\frac{1}{2}+\frac{\tau }{2}|\tau )$$
(A.7)
With this mapping $`u(\lambda )`$ between the $`\lambda `$ complex plane and the periodic rectangle of sides ($`1,\tau `$), we have:
$$u(d)=0,u(a)=\frac{1}{2},u(b)=\frac{1+\tau }{2},u(c)=\frac{\tau }{2},u(\mathrm{})=u_{\mathrm{}}$$
(A.8)
The inverse mapping can be written in terms of theta functions:
$$\lambda d=\frac{\theta _1^{}(0)}{2K}\frac{\theta _1^2(u)\theta _1(2u_{\mathrm{}})}{\theta _1(u+u_{\mathrm{}})\theta _1(uu_{\mathrm{}})\theta _1^2(u_{\mathrm{}})}$$
(A.9)
$$\sqrt{\sigma (\lambda )}=\frac{\theta _1^2(0)}{4K^2}\frac{\theta _1(2u)\theta _1(2u_{\mathrm{}})}{\theta _1^2(uu_{\mathrm{}})\theta _1^2(u+u_{\mathrm{}})}$$
(A.10)
and in terms of the usual trigonometric elliptic functions $`\mathrm{sn}`$, $`\mathrm{cn}`$, $`\mathrm{dn}`$ that we normalize to have periods $`1`$ and $`\tau `$, i.e.
$$\mathrm{sn}(u)=\mathrm{sn}(2K[m]u|m),\mathrm{dn}(u)=\mathrm{dn}(2K[m]u|m),\mathrm{}$$
(A.11)
one has
$$\frac{\mathrm{sn}^2(u)}{\mathrm{sn}^2(u_{\mathrm{}})}=\frac{\lambda d}{\lambda c},\frac{\mathrm{cn}^2(u)}{\mathrm{cn}^2(u_{\mathrm{}})}=\frac{\lambda a}{\lambda c},\frac{\mathrm{dn}^2(u)}{\mathrm{dn}^2(u_{\mathrm{}})}=\frac{\lambda b}{\lambda c}$$
$$\lambda c=(dc)\frac{1}{1\frac{\mathrm{sn}^2(u)}{\mathrm{sn}^2(u_{\mathrm{}})}},\lambda b=(db)\frac{\mathrm{dn}^2(u)}{1\frac{\mathrm{sn}^2(u)}{\mathrm{sn}^2(u_{\mathrm{}})}}$$
$$\lambda a=(da)\frac{\mathrm{cn}^2(u)}{1\frac{\mathrm{sn}^2(u)}{\mathrm{sn}^2(u_{\mathrm{}})}},\lambda d=(dc)\frac{da}{ca}\frac{\mathrm{sn}^2(u)}{1\frac{\mathrm{sn}^2(u)}{\mathrm{sn}^2(u_{\mathrm{}})}}$$
$$\sqrt{\sigma (\lambda )}=(dc)(da)\sqrt{\frac{db}{ca}}\frac{\mathrm{sn}(u)\mathrm{cn}(u)\mathrm{dn}(u)}{\left(1\frac{\mathrm{sn}^2(u)}{\mathrm{sn}^2(u_{\mathrm{}})}\right)^2}$$
and $`u_{\mathrm{}}`$ is related to $`a,b,c,d`$ by any of the following relations:
$$\mathrm{sn}^2(u_{\mathrm{}})=\frac{ca}{da},\mathrm{cn}^2(u_{\mathrm{}})=\frac{dc}{da},\mathrm{dn}^2(u_{\mathrm{}})=\frac{dc}{db}$$
## Appendix BComplex potentials
The case of complex potentials, that is to say of a polynomial potential $`V(\lambda )`$ with complex coefficients, is interesting for some applications of the matrix models to 2 dimensional gravity and when studying their connections with integrable hierarchies. In this case, the mean field large $`N`$ solution is known to be given by a continuous distribution of the eigenvalues along arcs in the complex plane .
In this appendix we show that our results are only slightly modified in this case.
In the two-cut case, we can repeat the analysis of sect.2. We fix $`x=n_1/N`$ (the proportion of e.v. in the first cut). The resolvent is still of the form 2.28, with the polynomial $`M`$ and the end points $`a,b,c,d`$ determined by the constraints 2.18 and 2.29, but they are no more real in general, as well as the resulting mean-field free energy $`F_0(x)`$.
If we now repeat the calculation of sect.2.2.3 we cannot use a saddle-point approximation for the sum over $`n_1`$ by expanding $`F_0(x)`$ around the saddle point $`x_0`$ which is the true extremum of $`F_0`$.
$$\frac{F_0}{x}(x_0)=\mathrm{\hspace{0.17em}0}$$
(B.1)
Indeed this extremum is at a finite non-zero distance of the real axis, i.e. $`\mathrm{Im}(x_0)=𝒪(1)`$, while the method of sect.2.2.3 is valid only if $`\mathrm{Im}(x_0)=𝒪(1/N)`$. However, since $`N`$ is integer, we can expand $`F_0`$ around any $`x_k`$, provided that
$$F_0^{}(x_k)=\mathrm{\hspace{0.17em}2}\mathrm{i}\pi \frac{k}{N},k$$
(B.2)
since the dangerous oscillating term $`\mathrm{e}^{N^2(xx_k)F_0^{}(x_c)}`$ is then a constant for $`x=n/N`$, $`n`$.
Therefore, as in we have to consider the real pseudo saddle-point $`x_c`$ such that
$$\mathrm{Re}(F_0^{}(x_c))=0\mathrm{with}\mathrm{Im}(x_c)=\mathrm{\hspace{0.17em}0}$$
(B.3)
and denote
$$\mathrm{\Delta }_c=\frac{1}{2\mathrm{i}\pi }F_0^{}(x_c)=\frac{1}{2\pi }\mathrm{Im}(F_0^{}(x_c))$$
(B.4)
We expand $`F_0`$ around some $`x_k`$ defined by Eq. (B.2) and such that
$$x_kx_c=𝒪(N^1)$$
(B.5)
and we get for the total free energy (by exactly the same calculation as in sect.2.2.3)
$`\begin{array}{ccc}\hfill F& =& N^2F_0(x_k)\mathrm{ln}\left(\theta _3(Nx_k)\right)\hfill \\ & & \\ & & +F_1(x_c)+\frac{1}{2}\mathrm{ln}\left(2\pi F_0^{\prime \prime }(x_c)\right)+𝒪(N^2)\hfill \end{array}`$ (B.9)
where $`\theta _3`$ is the theta function with modular parameter
$$\tau =\frac{2\mathrm{i}\pi }{F_0^{\prime \prime }(x_c)}$$
(B.10)
Only the first two terms are important for calculating the two-point functions and the orthogonal polynomials in the large $`N`$ limit. This leading term does not depend on $`k`$. Indeed we have
$$F_0^{}(x_k)F_0^{}(x_c)=(x_kx_c)F_0^{\prime \prime }(x_c)+𝒪(N^1)$$
(B.11)
hence
$$(x_kx_c)=\frac{1}{N}\tau (kN\mathrm{\Delta }_c)+𝒪(N^2)$$
(B.12)
and using the periodicity relations of $`\theta _3`$ we can rewrite the leading term for the free energy as
$$N^2F_0(x_k)\mathrm{ln}\left(\theta _3(Nx_k)\right)=N^2F_0(n_c/N)\frac{\mathrm{i}\pi }{\tau }u_c^2\mathrm{ln}\left(\theta _3(u_c)\right)$$
(B.13)
with
$$n_c=\mathrm{E}[Nx_c],u_c=[Nx_c]\tau [N\mathrm{\Delta }_c]$$
(B.14)
where $`\mathrm{E}[u]`$ is the integer part of $`u`$ (largest integer smaller than $`u`$) and $`[u]=u\mathrm{E}[u]`$ is the fractional part of $`u`$. This does not depend on $`k`$ up to negligible terms of order $`𝒪(N^1)`$ (provided that condition B.5 for $`k`$ holds).
One can now repeat the calculation of sect.3 for the 2-point function. Nothing is changed but we simply have to replace $`x_c`$ by $`x_k`$ in the intermediate steps and to use Eq. (B.12) at the end of the calculation. This amounts to replace the $`x_c`$ in the elliptic function $`\mathrm{sn}^2`$ by $`x_c\tau \mathrm{\Delta }_c`$. The final result for the two-point resolvent is
$`\omega ^c(\lambda ,\mu )`$ $`=`$ $`{\displaystyle \frac{1}{4(\lambda \mu )^2}}[(1\sqrt{{\displaystyle \frac{(\lambda a)(\lambda b)(\mu c)(\mu d)}{(\mu a)(\mu b)(\lambda c)(\lambda d)}}})+(\lambda \mu )]`$ (B.15)
$`{\displaystyle \frac{(ca)(db)}{4\sqrt{\sigma (\lambda )}\sqrt{\sigma (\mu )}}}\mathrm{sn}^2(N(x_c\tau \mathrm{\Delta }_c)+{\scriptscriptstyle \frac{1}{2}})`$
Similar results holds for the orthogonal polynomials. We simply have to consider the end-points $`a,b,c,d`$ for the mean-field real parameter $`x_c`$ and to make the replacement
$$Nx_cN(x_c+\tau \mathrm{\Delta }_c)$$
(B.16)
in the elliptic functions involving $`Nx_c`$. In any case these terms depend only on the fractional parts of $`Nx_c`$ and of $`N\mathrm{\Delta }_c`$.
A final interesting remark on the periodicity properties of the non-universal term $`\mathrm{sn}^2(N(x_c+\tau \mathrm{\Delta }_c))`$ can be made. From the definition B.4 $`\mathrm{\Delta }_c`$ corresponds to a “phase shift” between the two arcs where the density of e.v. is non zero.
$$\mathrm{\Delta }_c=\frac{1}{2\mathrm{i}\pi }F_0^{}(x_c)=\frac{\mathrm{\Gamma }_1\mathrm{\Gamma }_2}{2\mathrm{i}\pi }$$
(B.17)
where $`N\mathrm{\Gamma }_\alpha `$ is the (constant) effective potential on the arc $`\alpha `$. The two periods $`1`$ and $`\tau `$ of the $`\mathrm{sn}^2`$ function correspond respectively in term of eigenvalues to (i) transfer a single e.v. from the first arc to the second one ($`\delta Nx=\pm 1`$), (ii) or to shift the phase between the two arcs by $`2\pi `$ ($`\delta N\mathrm{\Delta }_c=\pm 1`$).
## Appendix CMulticut Case
Consider now a support of eigenvalues split into $`s`$ intervals:
$$𝒞=𝒞_1\mathrm{}𝒞_s$$
(C.1)
Let $`n_i`$ be the number of eigenvalues in each $`𝒞_i`$, and $`x_i=n_i/N`$ the occupation ratio, which we denote collectively as a vector:
$$x_i=_{𝒞_i}\text{d}\lambda \rho (\lambda ),\stackrel{}{x}=(x_1,\mathrm{},x_{s1})$$
(C.2)
Note that only $`s1`$ of them are independent since $`x_1+\mathrm{}+x_s=1`$.
As in the two-cut case (eq. 2.36), the free energy at fixed $`\stackrel{}{n}`$ admits a topological large $`N`$ expansion:
$$F[V;\stackrel{}{n}]=N^2F_0[V;\stackrel{}{x}]+N^0F_1[V;\stackrel{}{x}]+𝒪(1/N^2)$$
(C.3)
and as in 2.38, the partition function can be written as a sum over $`\stackrel{}{n}`$:
$$Z=\mathrm{e}^F=\underset{\stackrel{}{n}}{}\mathrm{e}^{F[V,\stackrel{}{n}]}$$
(C.4)
The sum is dominated by the vicinity of the extremum $`\stackrel{}{x}_c`$ of $`F_0[V;\stackrel{}{x}]`$:
$$Z\underset{\stackrel{}{n}}{}\mathrm{e}^{N^2\left(F_0[V,\stackrel{}{x_c}]+i\pi (\frac{\stackrel{}{n}}{N}\stackrel{}{x}_c).\tau ^1(\frac{\stackrel{}{n}}{N}\stackrel{}{x}_c)\right)}\mathrm{where}\frac{}{\stackrel{}{x}}F_0(\stackrel{}{x})|_{\stackrel{}{x}=\stackrel{}{x}_c}=\stackrel{}{0}$$
(C.5)
$`\tau `$ is the $`s1\times s1`$ matrix defined by:
$$\tau _{ij}^1=\frac{1}{2i\pi }\frac{^2F_0}{x_ix_j}|_{\stackrel{}{x}=\stackrel{}{x}_c}$$
(C.6)
Then, the summation over $`\stackrel{}{n}`$ yields:
$$Z\mathrm{e}^{N^2F_0[V,\stackrel{}{x}_c]}\theta (N\stackrel{}{x}_c|\tau )$$
(C.7)
where $`\theta (\stackrel{}{u}|\tau )`$ is Riemann’s theta function in genus $`s1`$:
$$\theta (\stackrel{}{u}|\tau )=\theta (\stackrel{}{u})=\underset{\stackrel{}{n}}{}\mathrm{e}^{i\pi (\stackrel{}{n}\stackrel{}{u}).\tau ^1(\stackrel{}{n}\stackrel{}{u})}=\underset{\stackrel{}{n}}{}\mathrm{e}^{i\pi \stackrel{}{n}.\tau \stackrel{}{n}}\mathrm{e}^{2i\pi \stackrel{}{n}.\stackrel{}{u}}$$
(C.8)
where $`\tau `$ is a $`s1\times s1`$ matrix, $`\stackrel{}{u}`$ is a $`s1`$ component vector, and $`\stackrel{}{n}`$ is a vector with integer coordinates.
The $`\theta `$ function obeys the relations ($`\stackrel{}{k}`$ being an arbitrary integer vector):
$$\theta (\stackrel{}{u}+\stackrel{}{k})=\theta (\stackrel{}{u}),\theta (\stackrel{}{u}+\tau \stackrel{}{k})=\mathrm{e}^{i\pi (2\stackrel{}{u}.\stackrel{}{k}+\stackrel{}{k}.\tau \stackrel{}{k})}\theta (\stackrel{}{u}),\theta (\stackrel{}{u})=\theta (\stackrel{}{u})$$
(C.9)
Eventually the free energy at leading orders in $`N`$ is:
$$FN^2\left[F_0[V,\stackrel{}{x}_c]\frac{1}{N^2}\mathrm{ln}\theta (N\stackrel{}{x_c}|\tau )+\frac{1}{N^2}F_1(\stackrel{}{x}_c)+\mathrm{}\right]$$
(C.10)
It is now straightforward but lengthy to rederive the 2-point correlation function and the orthogonal polynomials from C.10. One needs to differentiate C.10 with respect to variations of the potential as in 3.3 or 4.16, and express the hyperelliptical functions involved in the calculation through prime forms (hyperelliptical generalization of the $`\theta _1`$ function) . One should thus obtain expressions similar to those of . |
warning/0003/hep-th0003083.html | ar5iv | text | # Searching for 𝑆-duality in GravitationInvited talk at the Third Workshop on Gravitation and Mathematical-Physics, Nov. 28-Dec. 3 1999, León Gto. México.
## I Introduction
Strong/weak coupling duality ($`S`$-duality) in superstring and supersymmetric gauge theories in various dimensions has been, in the last five years, the major tool to study the strong coupling dynamics of these theories. Much of these results require supersymmetry through the notion of BPS state. These states describe the physical spectrum and they are protected of quantum corrections leaving the strong/coupling duality under control to extract physical information. In the non-supersymmetric case there are no BPS states and the situation is much more involved. This latter case is an open question and it is still under current investigation.
In the specific case of non-supersymmetric gauge theories in four dimensions, the subject has been explored recently in the Abelian as well as in the non-Abelian cases (for a review see ). In the Abelian case, one considers $`CP`$ non-conserving Maxwell theory on a curved compact four-manifold $`X`$ with Euclidean signature or, in other words, U(1) gauge theory with a $`\theta `$ vacuum coupled to four-dimensional gravity. The manifold $`X`$ is basically described by its associated classical topological invariants: the Euler characteristic $`\chi (X)=\frac{1}{16\pi ^2}_X\mathrm{tr}R\stackrel{~}{R}`$ and the signature $`\sigma (X)=\frac{1}{24\pi ^2}_X\mathrm{tr}RR`$. In the Maxwell theory, the partition function $`Z(\tau )`$ transforms as a modular form under a finite index subgroup $`\mathrm{\Gamma }_0(2)`$ of SL$`(2,𝐙)`$ , $`Z(1/\tau )=\tau ^u\overline{\tau }^vZ(\tau )`$, with the modular weight $`(u,v)=(\frac{1}{4}(\chi +\sigma ),\frac{1}{4}(\chi \sigma ))`$. In the above formula $`\tau =\frac{\theta }{2\pi }+\frac{4\pi i}{g^2}`$, where $`g`$ is the U(1) electromagnetic coupling constant and $`\theta `$ is the usual theta angle.
In order to cancel the modular anomaly in Abelian theories, it is known that one has to choose certain holomorphic couplings $`B(\tau )`$ and $`C(\tau )`$ in the topological gravitational (non-dynamical) sector, through the action
$$I^{TOP}=_X\left(B(\tau )\mathrm{tr}R\stackrel{~}{R}+C(\tau )\mathrm{tr}RR\right),$$
(1)
i.e., which is proportional to the appropriate sum of the Euler characteristic $`\chi (X)`$ and the signature $`\sigma (X)`$.
## II $`S`$-Duality in MacDowell-Mansouri Gauge Theory of Gravity
Let us briefly review the MacDowell-Mansouri (MM) proposal . The starting point for the construction of this theory is to consider an SO(3,2) gauge theory with a Lie algebra-valued gauge potential $`A_\mu ^{AB}`$, where the indices $`\mu =0,1,2,3`$ are space-time indices and the indices $`A,B=0,1,2,3,4`$. From the gauge potential $`A_\mu ^{AB}`$ we may introduce the corresponding field strength $`F_{\mu \nu }^{AB}=_\mu A_\nu ^{AB}_\nu A_\mu ^{AB}+\frac{1}{2}f_{CDEF}^{AB}A_\mu ^{CD}A_\nu ^{EF},`$ where $`f_{CDEF}^{AB}`$ are the structure constants of SO(3,2). MM choose $`F_{\mu \nu }^{a4}0`$ and as an action
$$S_{MM}=d^4xϵ^{\mu \nu \alpha \beta }ϵ_{abcd}F_{\mu \nu }^{ab}F_{\alpha \beta }^{cd},$$
(2)
where $`a,b,\mathrm{}\mathrm{etc}.=0,1,2,3.`$
On the other hand, by considering the self-dual (or anti-self-dual) part of the connection, a generalization has been proposed . The extension to the supergravity case is considered in .
One can then search whether the construction of a linear combination of the corresponding self-dual and anti-self-dual parts of the MacDowell-Mansouri action can be reduced to the standard MM action plus a kind of $`\mathrm{\Theta }`$-term and, moreover, if by this means one can find the “dual-theory” associated with the MM theory. This was showed in and the corresponding extension to supergravity is given at . In what follows we follow Ref. . Let us consider the action
$$S=d^4xϵ^{\mu \nu \alpha \beta }ϵ_{abcd}\left({}_{}{}^{+}\tau {}_{}{}^{+}F_{\mu \nu }^{ab}{}_{}{}^{+}F_{\alpha \beta }^{cd}{}_{}{}^{}\tau {}_{}{}^{}F_{\mu \nu }^{ab}{}_{}{}^{}F_{\alpha \beta }^{cd}\right),$$
(3)
where $`{}_{}{}^{\pm }F_{\mu \nu }^{ab}=\frac{1}{2}\left(F_{\mu \nu }^{ab}\pm \stackrel{~}{F}_{\mu \nu }^{ad}\right)`$ and $`\stackrel{~}{F}_{\mu \nu }^{ab}=\frac{1}{2}iϵ_{cd}^{ab}F_{\mu \nu }^{cd}`$. It can be easily shown , that this action can be rewritten as
$$S=\frac{1}{2}d^4xϵ^{\mu \nu \alpha \beta }ϵ_{abcd}\left[({}_{}{}^{+}\tau {}_{}{}^{}\tau )F_{\mu \nu }^{ab}F_{\alpha \beta }^{cd}+({}_{}{}^{+}\tau +{}_{}{}^{}\tau )F_{\mu \nu }^{ab}\stackrel{~}{F}_{\alpha \beta }^{cd}\right].$$
(4)
In their original paper, MM have shown that the first term in this action reduces to the Euler topological term plus the Einstein-Hilbert action with a cosmological term. This was achieved after identifying the components of the gauge field $`A_\mu ^{AB}`$ with the Ricci rotation coefficients and the vierbein. Similarly, the second term can be shown to be equal to $`i\theta P`$, where $`P`$ is the Pontrjagin topological term . Thus, it is a genuine $`\theta `$ term, with $`\theta `$ given by the sum $`{}_{}{}^{+}\tau +{}_{}{}^{}\tau `$.
Our second task is to find the “dual theory”, following the same scheme as for Yang-Mills theories . For that purpose we consider the parent action
$$I=d^4xϵ^{\mu \nu \alpha \beta }ϵ_{abcd}\left(c_1{}_{}{}^{+}G_{\mu \nu }^{ab}{}_{}{}^{+}G_{\alpha \beta }^{cd}+c_2{}_{}{}^{}G_{\mu \nu }^{ab}{}_{}{}^{}G_{\alpha \beta }^{cd}+c_3{}_{}{}^{+}F_{\mu \nu }^{ab}{}_{}{}^{+}G_{\alpha \beta }^{cd}+c_4{}_{}{}^{}F_{\mu \nu }^{ab}{}_{}{}^{}G_{\alpha \beta }^{cd}\right).$$
(5)
From which the action (3) can be recovered after integration on $`{}_{}{}^{+}G`$ and $`{}_{}{}^{}G`$.
In order to get the “dual theory” one should start with the partition function
$$Z=𝒟{}_{}{}^{+}G𝒟{}_{}{}^{}G𝒟Ae^I.$$
(6)
To proceed with the integration over the gauge fields we observe that $`F_{\mu \nu }^{ab}=_\mu A_\nu ^{ab}_\nu A_\mu ^{ab}+\frac{1}{2}f_{CDEF}^{ab}A_\mu ^{CD}A_\nu ^{EF}.`$ Taking into account the explicit expression for the structure constants, the second term of $`F_{\mu \nu }^{ab}`$ will naturally split in four terms given by $`A_\mu ^{ad}A_{\nu d}^b`$ $`A_\nu ^{ad}A_{\mu d}^b`$ $`\lambda ^2\left(A_\mu ^{a4}A_\nu ^{b4}A_\nu ^{a4}A_\mu ^{b4}\right).`$ The integration over the components $`A_\mu ^{a4}`$ is given by a Gaussian integral, which turns out to be $`det𝐆^{1/2}`$, where G is a matrix given by $`𝐆_{ab}^{\mu \nu }`$ $`=8i\lambda ^2ϵ^{\mu \nu \alpha \beta }\left(c_3{}_{}{}^{+}G_{\alpha \beta ab}^{}c_4{}_{}{}^{}G_{\alpha \beta ab}^{}\right).`$
Thus, the partition function (6) can be written as
$$Z=𝒟{}_{}{}^{+}G𝒟{}_{}{}^{}G𝒟A_\mu ^{ab}𝑑et𝐆^{1/2}e^{II},$$
(7)
where
$$II=2id^4xϵ^{\mu \nu \alpha \beta }\left[c_1{}_{}{}^{+}G_{\mu \nu }^{ab}{}_{}{}^{+}G_{\alpha \beta ab}^{}c_2{}_{}{}^{}G_{\mu \nu }^{ab}{}_{}{}^{}G_{\alpha \beta ab}^{}+2H_{\mu \nu }^{ab}(c_3{}_{}{}^{+}G_{\alpha \beta ab}^{}c_4{}_{}{}^{}G_{\alpha \beta ab}^{})\right],$$
and $`H_{\mu \nu }^{ab}=_\mu A_\nu ^{ab}_\nu A_\mu ^{ab}+\frac{1}{2}f_{cdef}^{ab}A_\mu ^{cd}A_\nu ^{ef}`$ is the SO(3,1) field strength.
Our last step to get the dual action is to integrate over $`A_\mu ^{ab}`$. This kind of integration is well known and has been performed in previous works . The result is
$$Z=𝒟{}_{}{}^{+}G𝒟{}_{}{}^{}G𝑑et𝐆^{1/2}𝑑et({}_{}{}^{+}M)^{1/2}𝑑et({}_{}{}^{}M)^{1/2}e^{{\scriptscriptstyle d^4x\stackrel{~}{L}}},$$
(8)
with
$$\begin{array}{cc}\stackrel{~}{L}\hfill & =ϵ^{\mu \nu \rho \sigma }[\frac{1}{4{}_{}{}^{+}\tau }{}_{}{}^{+}G_{\mu \nu }^{ab}{}_{}{}^{+}G_{\rho \sigma ab}^{}+\frac{1}{4{}_{}{}^{}\tau }{}_{}{}^{}G_{\mu \nu }^{ab}{}_{}{}^{}G_{\rho \sigma ab}^{}+2_\nu {}_{}{}^{+}G_{\rho \sigma ab}^{}({}_{}{}^{+}M)_{\mu \lambda }^{1abcd}ϵ^{\lambda \theta \alpha \beta }_\theta {}_{}{}^{+}G_{\alpha \beta cd}^{}\hfill \\ & 2_\nu {}_{}{}^{}G_{\rho \sigma ab}^{}({}_{}{}^{}M)_{\mu \lambda }^{1abcd}ϵ^{\lambda \theta \alpha \beta }_\theta G_{\alpha \beta cd}^{}],\hfill \end{array}$$
(9)
where $`{}_{}{}^{\pm }M_{ab}^{\mu \nu cd}=\frac{1}{2}ϵ^{\mu \nu \alpha \beta }\left(\delta _a^c{}_{}{}^{\pm }G_{\alpha \beta b}^{d}+\delta _b^c{}_{}{}^{\pm }G_{\alpha \beta a}^{d}+\delta _a^d{}_{}{}^{\pm }G_{\alpha \beta b}^{c}\delta _b^d{}_{}{}^{\pm }G_{\alpha \beta a}^{c}\right)`$ and $`{}_{}{}^{+}\tau =\frac{1}{4c_1},`$ $`{}_{}{}^{}\tau =\frac{1}{4c_2}`$, $`c_3=c_4=1.`$
The non-dynamical model considered in a previous work results in a kind of non-linear sigma model of the type considered by Freedman and Townsend , as in the usual Yang-Mills dual models. The dual to the dynamical gravitational model (9) considered here, results in a Lagrangian of the same structure. However, it differs from the non-dynamical case by the features discussed above.
## III (Anti)Self-duality of the Three-dimensional Chern-Simons Gravity
It is well known that the $`2+1`$ Einstein-Hilbert action with nonvanishing cosmological constant $`\lambda `$ is given by the “standard” and “exotic” Einstein actions . It is well known that for $`\lambda >0`$ (and $`\lambda <0`$), these actions are equivalent to a Chern-Simons actions in $`2+1`$ dimensions with gauge group $`𝒢`$ to be SO(3,1) (and SO(2,2)).
In this section we will work out the Chern-Simons Lagrangian for (anti)self-dual gauge connection with respect to duality transformations of the internal indices of the gauge group $`𝒢`$, in the same philosophy of MM , and that of
$$L{}_{}{}^{\pm }{}_{CS}{}^{}=_{}\epsilon ^{ijk}({}_{}{}^{\pm }A_{i}^{AB}_j{}_{}{}^{\pm }A_{kAB}^{}+\frac{2}{3}{}_{}{}^{\pm }A_{iA}^{B}{}_{}{}^{\pm }A_{jB}^{C}{}_{}{}^{\pm }A_{kC}^{A}),$$
(10)
where $`A,B,C,D=0,1,2,3,`$ $`\eta _{AB}=diag(1,+1,+1,+1)`$ and the complex (anti) self-dual connections are $`{}_{}{}^{\pm }A_{i}^{AB}=\frac{1}{2}(A_i^{AB}\frac{i}{2}\epsilon _{CD}^{AB}A_i^{CD})`$, which satisfy the relation $`\epsilon _{CD}^{AB}{}_{}{}^{\pm }A_{i}^{CD}=\pm i{}_{}{}^{\pm }A_{}^{AB}.`$
Thus using the above equations, the action (10) can be rewritten as
$$L{}_{}{}^{\pm }{}_{CS}{}^{}=_{}\frac{1}{2}\epsilon ^{ijk}(A_i^{AB}_jA_{kAB}+\frac{2}{3}A_{iA}^BA_{jB}^CA_{kC}^A)\frac{i}{4}\epsilon ^{ijk}\epsilon ^{ABCD}(A_{iAB}_jA_{kCD}+\frac{2}{3}A_{i\text{ }A}^EA_{jEB}A_{kCD}).$$
(11)
In this expression the first term is the Chern-Simons action for the gauge group $`𝒢`$, while the second term appears as its corresponding “$`\theta `$-term”. The same result was obtained in $`3+1`$ dimensions when we considered the (anti)self-dual MM action , or the (anti)self-dual $`3+1`$ pure topological gravitational action .
One should remark that the two terms in the action (13) are the Chern-Simons and the corresponding “$`\theta `$-term” for the gauge group $`𝒢`$ under consideration. After imposing the particular identification $`A_i^{AB}=(A_i^{ab},A_i^{3a})=(\omega _i^{ab},\sqrt{\lambda }e_i^a)`$ and $`\omega _i^{ab}=\epsilon ^{abc}\omega _{ic}`$, the “exotic” and “standard” actions for the gauge group SO(3,1) are given respectively by
$`L_{CS}{}_{}{}^{\pm }={\displaystyle _X}{\displaystyle \frac{1}{2}}\epsilon ^{ijk}(\omega _i^a(_j\omega _{ka}_k\omega _{ja})+{\displaystyle \frac{2}{3}}\epsilon _{abc}\omega _i^a\omega _j^b\omega _k^c+\lambda e_i^a(_je_{ka}_ke_{ia})2\lambda \epsilon _{abc}e_i^ae_j^b\omega _k^c)`$
$$\pm i\sqrt{\lambda }\epsilon ^{ijk}\left(e_i^a(_j\omega _{ka}_k\omega _{ja})\epsilon _{abc}e_i^a\omega _j^b\omega _k^c+\frac{1}{3}\lambda \epsilon _{abc}e_i^ae_j^be_k^c\right),$$
(12)
plus surface terms. It is interesting to note that the above action (12) can be obtained from action (1) (for a suitable choice of $`B(\tau )`$ and $`C(\tau )`$) by dimensional reduction from $`X`$ to its boundary $`=X`$. Thus the “standard” action come from the Euler characteristic $`\chi (X)`$, while the “exotic” action come from the signature $`\sigma (X).`$
## IV Chern-Simons Gravity Dual Action in Three Dimensions
This section is devoted to show that a “dual” action to the Chern-Simons gravity action can be constructed following . Essentially we will repeat the procedure to find the the “dual” action to MM gauge theory given in Sec. II.
We begin from the original non-Abelian Chern-Simons action given by
$$L=_{}d^3x\frac{g}{4\pi }\epsilon ^{ijk}A_i^{AB}\left(_jA_{kAB}+\frac{1}{3}f_{ABCDEF}A_j^{CD}A_k^{EF}\right).$$
(13)
Now, as usual we propose a parent action in order to derive the dual action
$$L_D=_{}d^3x\epsilon ^{ijk}\left(aB_i^{AB}H_{jkAB}+bA_i^{AB}G_{jkAB}+cB_i^{AB}G_{jkAB}\right),$$
(14)
where $`H_{jkAB}=_jA_{kAB}+\frac{1}{3}f_{ABCDEF}A_j^{CD}A_k^{EF}`$ and $`B_i^{AB}`$ and $`G_{ij}^{AB}`$ are vector and tensor fields on $``$. It is a very easy matter to show that the action (13) can be derived from this parent action after integration of $`G`$ fields
The “dual” action $`L_D^{}`$ can be computed as usually in the Euclidean partition function, by integrating first out with respect to the physical degrees of freedom $`A_i^{AB}.`$ The resulting action is of the Gaussian type in the variable $`A`$ and thus, after some computations, it is easy to find the “dual” action
$$L_D^{}=_{}d^3x\epsilon ^{ijk}\left\{\frac{3}{4a}(a_iB_{jAB}+bG_{ijAB})[𝐑^1]_{kn}^{ABCD}\epsilon ^{lmn}(a_lB_{mCD}+bG_{lmCD})+c\alpha _i^{AB}G_{jkAB}\right\},$$
(15)
where $`[𝐑]`$ is given by $`[𝐑]_{ABCD}^{ij}=\epsilon ^{ijk}f_{ABCD}^{EF}B_{kEF}`$ whose inverse is defined by $`[𝐑]_{ABCD}^{ij}[𝐑^1]_{jk}^{CDEF}=\delta _k^i\delta _{AB}^{EF}.`$
The partition function is finally given by
$$Z=𝒟G𝒟B\sqrt{det(𝐌^1)}exp\left(L_D^{}\right).$$
(16)
In this “dual action” the $`G`$ field is not dynamical and can be integrated out. The integration of this auxiliary field gives
$$L_D^{}=_{}d^3x\frac{4\pi }{g}\epsilon ^{lmn}\left(B_l^{AB}_mB_{nAB}\frac{4\pi }{g}f_{ABCDEF}B_l^{AB}B_m^{CD}B_n^{EF}\right).$$
(17)
The fields $`B`$ cannot be rescaled if we impose “periodicity” conditions on them. Thus, this dual action has inverted coupling with respect to the original one (compare with for the Abelian case).
Acknowledgments
The results of sections 3 and 4 were obtained in collaboration with Miguel Sabido. We are very grateful to him for a critical reading of this manuscript. |
warning/0003/math0003082.html | ar5iv | text | # Notes for a Quantum Index Theorem
## 0 Introduction.
These notes are a natural outgrowth of our previous work on a local holomorphic formula for the dimension of a superselection sector and were motivated by the purpose to give a geometrical picture to aspects of local quantum physics related to the superselection structure. They may be read from different points of view, in particular, guided by a similarity of the statistical dimension with the Fredholm index and a possible index theorem, already suggested in , we shall regard the DHR localized endomorphisms as quantum analogs of elliptic differential operators.
A heuristic preamble. We begin to give a heuristic, but elementary, motivation for our dimension formula, postponing for the moment the specification of the underlying structure. Let the selfadjoint operator $`H_0`$ be a reference Hamiltonian for a Quantum Statistical Mechanics system, generating the evolution on an operator algebra $`𝔄`$. $`H_0`$ may be thought to correspond to the Laplacian $`\mathrm{\Delta }`$ on a compact Riemann manifold. We write any other Hamiltonian $`H`$ as a perturbation $`H=H_0+P`$ of $`H_0`$. We assume a supersymmetric structure, essentially the existence of a Dirac operator $`D`$ (an odd square root of $`H`$) that implements an odd derivation of $`𝔄`$.
The McKean–Singer lemma then shows that
$$\text{Tr}_s(e^{\beta H})=\text{Fredholm index of }D$$
(1)
for any $`\beta >0`$, where $`\text{Tr}_s`$ denotes the super-trace. In particular $`\text{Tr}_s(e^{\beta H})`$ (the Witten index) is an integer.
Now $`\omega ^{(\beta )}=\text{Tr}_s(e^{\beta H_0})/\text{Tr}_s(e^{\beta H_0})`$ is the normalized super-Gibbs functional at inverse temperature $`\beta `$ for the dynamics generated by $`H_0`$ and if we consider the unitary cocycle
$$u_P(t)e^{itH}e^{itH_0}$$
(2)
relating the two evolutions, and that belongs to $`𝔄`$ if $`P𝔄`$, we may write the following formula in terms of $`\omega ^{(\beta )}`$ and $`u_P`$:
$$\underset{ti\beta }{\text{anal.cont. }}\omega ^{(\beta )}(u_P(t))=\frac{\text{Tr}_s(e^{\beta H})}{\text{Tr}_s(e^{\beta H_0})},$$
(3)
provided the latter makes sense. We will thus regard the above expression as a multiplicative relative index between $`H_0`$ and $`H`$.<sup>1</sup><sup>1</sup>1As a counterpart, an additive relative index generalizing $`\text{Tr}_s(e^{\beta H}e^{\beta H_0})`$ appeared in .
As shown in this index is invariant under deformations, in particular $`\text{Index}(u_P)=\text{Index}(u_I)`$ if $`P`$ is an odd element of $`𝔄`$ in the domain of the superderivation and may be obtained by evaluating at the identity the JLO Chern character associated with $`\omega ^{(\beta )}`$.
At infinite volume however, the integrality of the index is not evident.
In the case of infinite volume systems, likewise for the Laplacian on a non-compact manifolds, the Hamiltonian has not any longer discrete spectrum and may exist only as a derivation. What survives after the thermodynamical limit is the time evolution, a one-parameter automorphism group $`\alpha `$ of a C-algebra $`𝔄`$. The equilibrium states $`\omega ^{(\beta )}`$ are characterized by the KMS condition . However the cocycle (2) may well exist and belong to $`𝔄`$, so that formula (3) may be generalized to define the index of a cocycle, if the analytic continuation exists.
In this paper we shall consider the case of Relativistic Quantum Statistical Mechanics, by which we mean the consideration of KMS functionals for the time evolution in Quantum Field Theory. The basic relativistic property, locality or finite propagation of speed of light, will select the appropriate class of cocycles and will imply the integrality of the index as will be explained. For reasons that will partly be clarified, the supersymmetric structure will not play a direct role in our formulae for the dimension.
Ingredients for QFT analysis. Let us discuss now some fundamental aspects of Physics and Analysis. Quantum Field Theory can be considered at the same time as a generalization of two Physical theories of very different nature: Classical (Lagrangian) Field Theory and Quantum Mechanics. Both of them extend Classical Mechanics, but point in apparently divergent directions. In the first case one goes from finitely many to infinitely many degrees of freedom, but remains in the classical framework. In the second case one replaces classical variables by quantum variables (operators), but remains within finitely many degree of freedoms. Quantum Field Theory inherits the richness of the two theories by treating infinitely many quantum variables and enhances them further, in particular by interaction, particle creation/annihilation, special relativity.
There are thus two paths from the finite-dimensional classical calculus to QFT according to the following diagram:
$$\begin{array}{ccc}\text{Classical, finite dim.}& & \text{Variational calculus}\\ & & & & \\ \text{Quantum, finite dim.}& & \text{Quantum Field Theory}\end{array}$$
(4)
Note now that the passage from ordinary manifolds to variational calculus did not require a new calculus; for example the notion of derivative still make sense replacing points by functions.
On the other hand, the passage from classical to quantum mechanics does require a new structure (non-commutativity) a new calculus. The standard quantization procedure replaces functions by selfadjoint operators and Poisson brackets by commutators. In this correspondence $`x_hP_h`$ and $`i\frac{}{x_h}Q_h`$ give position and momentum operators that satisfy the Heisenberg commutation relations $`[P_h,Q_k]=i\delta _{hk}I`$.
A quantized, finite-dimensional, calculus has been developed in recent times by A. Connes (see ); a sample dictionary is here below:
| CLASSICAL | QUANTUM |
| --- | --- |
| Variable | Operator |
| Differential | $`[F,]`$ |
| Integral | $``$ (Dixmier trace) |
| Infinitesimal | Compact operator |
| $`\mathrm{}`$ | $`\mathrm{}`$ |
Concerning Quantum Field Theory, more or less implicit suggestions concerning a “second quantized” or QFT calculus can be found in . In particular we consider Jones subfactors and endomorphisms or Connes correspondences to be basic objects in this setting. The underlying structure at each level is illustrated in the following table:
| CLASSICAL | Classical variables Differential forms Chern classes | Variational calculus Infinite dimensional manifolds Functions spaces Wiener measure |
| --- | --- | --- |
| QUANTUM | Quantum geometry Fredholm operators Index Cyclic cohomology | Subfactors Correspondences, Endomorphisms Multiplicative index Supersymmetric QFT, $`(𝔄,,Q)`$ |
Note that there is a non-trivial map from
$$\overline{)\text{points}}\overline{)\text{fields}},$$
horizontally in the diagram (4), that further enriches the structure. At the quantum level this is the second quantization functor; this is partly at the basis of the multiplicative structure of the index (cf. for an example).
We mention a first result, due to Connes , that may be read within the context of QFT analysis: the index map
$$\text{Ind}_Q:K_0(A(𝒪))$$
is not polynomial and the K-theory group $`K_0(A(𝒪))`$ is of infinite rank. Here $`A(𝒪)`$ is a “smooth” local Bosonic algebra associated with a free massive supersymmetric field on the cylinder.
After these premises, let us discuss the basic objects of our analysis, superselection sectors.
Superselection sectors as QFT analogs of elliptic operator. The celebrated Atiyah–Singer index theorem equates the analytic index of an elliptic operator to a geometric–topological index. The analytic index is the Fredholm index, which is manifestly an integer. The geometric index is intrinsically invariant under deformations. A major consequence of the index theorem is then the integrality of the geometical index.
As discussed, Operator Algebras provide the proper quantization (non commutative setting) for measure theory, topology and geometry. In particular, extensions of the index theorem by means of noncommutative K-theory and cyclic cohomology occur naturally in Noncommutative Geometry . These results pertain to Connes quantized calculus.
Here we shall deal with a quantum field theory analog of the index theorem. The role of the Fredholm linear operators is now played by the endomorphisms of an infinite factor with finite Jones index . In Quantum Field Theory the index-statistics theorem equates the DHR statistical dimension with the square root of the Jones index or, more precisely, with the minimal dimension, whence the integrality of the index is immediate by the integrality of the statistical dimension . We then look for possible geometric counterparts of the statistical dimension.
Our framework is Quantum Field Theory and its superselection structure , in the Doplicher-Haag-Roberts framework . The local observable algebra $`𝔄(𝒪)`$ associated with region $`𝒪`$ of the spacetime provides a noncommutative version of the algebra of functions with support in $`𝒪`$, and the localized endomorphisms with finite statistics are analog to the elliptic differential operators, as suggested by S. Doplicher . Note indeed that an endomorphism $`\rho `$ localized in the double cone $`𝒪_0`$ is local in the sense that
$$\rho (𝔄(𝒪))𝔄(𝒪),𝒪𝒪_0,$$
similarly to the locality property that characterizes the differential operators in the classical setting . More is true, the correspondence $`𝒪𝔄(𝒪)`$ is endowed with a natural, but not manifest, sheaf structure with respect to which covariant localized endomorphisms are sheaf maps, a fact that will be used only implicitly (it gives the automatic covariance used in Sect. 3.2).
It is now clear that, in our context, geometric information may be contained in the classical geometry of the spacetime and in the net
$$𝒪𝔄(𝒪).$$
A geometrical description of the superselection structure of $`𝔄`$ as been given by Roberts , who defined a non-abelian cohomology ring $`H_R^1(𝔄)`$ whose elements correspond to the superselection sectors of $`𝔄`$. In his formalism, however, it is unclear how to integrate a cohomological class to obtain an invariant, the dimension.
A conceptually different cohomological structure appears in by considering unitary cocycles associated with a dynamics. As we shall see in Sect. 4, if we consider only localized unitary cocycles associated with translations, we obtain a cohomology ring $`H_\tau ^1(𝔄)`$ which describes the covariant superselection sectors. Denoting by $`𝔖_{KMS}`$ the set of extremal KMS states for the time evolution, at inverse temperature $`\beta `$, satisfying Haag duality, we have indeed a pairing
$$𝔖_{KMS}\times H_\tau ^1(𝔄)\phi \times [u]\phi ,[u]=u(i\beta )\text{d}\phi $$
(5)
that we shall below describe in equations (6,7). We have used, only here, the notation $`u\text{d}\phi \phi (u)`$ in order to provide resemblance with the classical context.
It is now useful to compare in a table the context in the Atyiah-Singer theorem and in Quantum Field Theory, including the role possibly played by supersymmetry (see below).
Atiyah-Singer context QFT context sheaf structure functions on manifold $`𝒱`$ net (sheaf) of C-algebras on $`𝒱`$ smooth structure smooth functions net of dense -algebras differential operator sheaf map localized endomorphism elliptic operator Fredholm opearator finite index endomorphism analitical index Fredholm index minimal dimension intergrality Fredholm index $``$ statistical dimension $``$ geometric index associated with $`(D,𝒱)`$ associated with $`(\rho ,𝔄,𝒱)`$ cohomology De Rham Roberts; cyclic deformation invariance intrinsic perturbation invariance Chern character Chern character pairing (5); JLO cyclic cocycle Hamiltonian Laplacian Killing Hamiltonian grading Dirac operator supersymmetry spectral formula heat kernel (super)-Gibbs state
In this table we have considered finite volume spaces, although we will deal with non-compact spacetimes. As is known, Gibbs states becomes KMS states at infinite volume . The analogue replacement of super-Gibbs functionals by super-KMS functionals is far less obvious, see the outlook. However, as mentioned, for many purposes concerning the index, one may work with ordinary KMS states.
Black holes, conformal symmetries and holography. There is a first important setting where the above programme may be implemented, with a formula for the geometric index involving (classical) spacetime geometry, namely the analysis of charge addition for a quantum black hole in a thermal state.
We introduce here another piece of structure, namely we consider Quantum Field Theory on a curved spacetime. This physical theory combines General Relativity and Quantum Field Theory, but treats the gravitational field as a background field and therefore disregards effects occurring at the Planck length. Yet important effects, as the Hawking effect , pertain to this context.
We start with a black hole described by a globally hyperbolic spacetime with bifurcate Killing horizon, for example the Schwarzschild-Kruskal spacetime, and we consider quantum effects on this gravitational background. We then consider the incremental entropy due to the addition of a localized charge. The case of the Rindler spacetime has been previously described by means of a local analogue of a Kac-Wakimoto formula .
Here however we use a different point of view and conceptual scheme. The first basic point is that the restriction of the net $`𝔄`$ to each of the two horizon components $`𝔥_+`$ and $`𝔥_{}`$ gives a conformal net on $`S^1`$, a general fact that is obtained by applying Wiesbrock’s characterization of conformal nets , a structure already discussed in . It may appear analogous to the holography on the anti-de Sitter spacetime that independently appeared in the Maldacena-Witten conjecture<sup>2</sup><sup>2</sup>2We thank M. Konsevich for pointing out this similarity to us. , proved by Rehren . Yet their context differs inasmuch as the anti-de Sitter spacetime is not globally hyperbolic and the holography there is a peculiarity of that spacetime, rather than a general phenomenal.
With these conformal nets at one hand, and assuming the due duality properties, we may consider endomorphisms $`\rho `$ and $`\sigma `$ of $`𝔄`$ that are localizable on $`𝔥_+`$ or $`𝔥_{}`$. In particular, we will have the formula for the difference of the logarithm of dimensions:
$$\mathrm{log}d(\rho )\mathrm{log}d(\sigma )=\frac{2\pi }{\kappa (𝒱)}(F(\phi _\rho |\phi _\sigma )+F(\phi _{\overline{\rho }}|\phi _{\overline{\sigma }})),$$
where $`F`$-terms represent the incremental free energy between the thermal equilibrium states with the charges $`\rho `$ and $`\sigma `$ or the conjugate charges $`\overline{\rho }`$ and $`\overline{\sigma }`$. Here $`\phi `$ is the Hartle-Hawking state and $`\phi _\rho `$ is the corresponding equilibrium state in presence of the charge $`\rho `$. The geometry appears here in the surface gravity $`\kappa (𝒱)`$ associated to the spacetime manifold $`𝒱`$, see . As a consequence the right hand side is the difference of the logarithm of two integers.
We will consider also different temperature states. For a finer analysis of this context we refer to our Sect. 5. We shall need to study the chemical potential, as we are going to explain.
Chemical potential. Relativistic case. The chemical potential is a label that each charge sets on different equilibrium states at the same temperature. Its structure in Quantum Statistical Mechanics has been explained in the work of Araki, Haag, Kastler and Takesaki , see also . The labels appear by considering the extensions of these states from the observable algebra to KMS states of the field algebra (with the time evolution modified by one-parameter subgroups of the gauge group).
Our aim is to consider temperature states in Quantum Field Theory. This may be motivated by wish to study extreme physical contexts such as the early universe (not discussed in this paper) or black holes.
In the context of Quantum Field Theory on the $`d+1`$-dimensional Minkowski spacetime with $`d2`$, Doplicher and Roberts have constructed the field algebra associated with local observables and short range interaction charges.
It is not difficult to apply the AHKT analysis, originally made in the context of C-dynamics, to the case of Quantum Field Theory on the Minkowski spacetime as above. It turns out the every covariant irreducible localized endomorphism $`\rho `$ with finite dimension extends to the weak closure in the GNS representation $`\pi _\phi `$ associated with a KMS state $`\phi `$, a fact that should be expected on physical grounds because the addition of a single charge should not lead to an inequivalent representation for an infinite system.
If $`\phi `$ is extremal KMS, we are then led to the context of endomorphisms of factors with finite index. Assuming $`\pi _\phi `$ to satisfy Haag duality, the extension of $`\rho `$ to the weak closure is still irreducible. If $`u`$ is the time covariance cocycle for $`\rho `$, then $`u`$ is a Connes Radon-Nikodym cocycle up to phase, hence it satisfies certain holomorphic properties. We have the formula, that generalizes ,
$$\mathrm{log}d_\phi (u)=\mathrm{log}d(\rho )+\beta \mu _\rho (\phi ),$$
(6)
where the holomorphic dimension is defined by
$$d_\phi (u)\underset{ti\beta }{\text{anal.cont. }}\phi (u(t)).$$
(7)
If there is a canonical choice for $`u`$, for example if the $`\rho `$ is Poincaré covariant in the vacuum sector, then the above formulae define the chemical potential $`\mu _\rho (\phi )`$ of $`\phi `$ corresponding to the charge $`\rho `$. This extends an analogous expression in the AHKT work in the case of abelian charges ($`d(\rho )=1`$). In general only the difference $`\mu _\rho (\phi )\mu _\rho (\psi )`$ is intrinsic.
Now the localized endomorphisms form a C-tensor category and in a C-tensor category there exists a natural anti-linear conjugation on the arrows between finite-dimensional objects :
$$T(\rho ,\sigma )T^{}(\overline{\rho },\overline{\sigma }).$$
Therefore, even if in general $`u`$ is defined only up to phase, $`u^{}`$ is a covariance cocycle for the conjugate charge $`\overline{\rho }`$ with opposite chemical potential
$$\mu _{\overline{\rho }}(\phi )=\mu _\rho (\phi )$$
(8)
and we obtain an expression for the intrinsic dimension $`d(\rho )`$ as the geometric mean
$$d(\rho )=\sqrt{d_\phi (u)d_\phi (u^{})},$$
which is independent of the phase fixing. We regard the right hand side of this expression as a geometric dimension, according to what was explained before, being candidate for geometric interpretation.
The quantity $`\beta ^1\mathrm{log}d_\phi (u)`$ represents the incremental free energy (adding the charge $`\rho `$). It has then a canonical decomposition as a sum of an intrinsic part $`\beta ^1\mathrm{log}d(\rho )`$, which is independent of $`\phi `$, where $`\mathrm{log}d(\rho )`$ is half the incremental entropy associated with $`\rho `$ (cf. ), and the chemical potential part $`\mu _\rho (\phi )`$ characterized by the asymmetry with respect to charge conjugation in (8).
Chemical potential. Low dimensional case. Motivated by black hole thermodynamics and the associated conformal nets on the black hole horizon, among other considerations, one is led to the analysis of the chemical potential in one-dimensional Quantum Field Theory. In this context the field algebra does not any longer exist and the AHKT work is not applicable.
The notion of holomorphic dimension still makes sense and is the basis of our analysis. The crucial point however is to show that localized endomorphisms are normal in the representation $`\pi _\phi `$ associated to a thermal state $`\phi `$. Based on Wiesbrock characterization of conformal nets on $`S^1`$ , we shall see that there is a conformal net associated with $`\phi `$, the conformal thermal completion. If $`\pi _\phi `$ satisfies duality for half-lines, we prove that the thermal completion automatically satisfies Haag duality, namely it is strongly additive .
Then, a localized endomorphism $`\rho `$ with finite dimension of the original net gives rise to a transportable localized endomorphism of the thermal completion. By a result in $`\rho `$ is automatically conformally covariant. We finally use this conformal covariance to show the normality of $`\rho `$ in the representation $`\pi _\phi `$.
For completeness, we extend the work of AHKT also in regard to describing the chemical potential in terms of extensions of KMS states. This will be achieved by considering extensions to the quantum double, a C-algebraic version of a construction in , that is a substitute for the non-existing field algebra.
At this point the analysis of the chemical potential goes through as in the higher dimensional case with the corresponding formulae for the incremental free energy. In the case of globally hyperbolic spacetimes with bifurcate Killing horizon, these results allow one to treat thermal states for the Killing evolution and charges localizable on the horizon, as explained.
The expected role of supersymmetry. Connes cyclic cohomology enters in Supersymmetric Quantum Field Theory via the work of Jaffe, Lesnieski and Osterwalder , see also . There is a noncommutative Chern character associated with a thermal equilibrium state, i.e. an entire cyclic cocycle associated to a supersymmetric KMS functional for the time evolution. In this context an index formula appears, which essentially coincides with our previous formula (7). There, however, the index formula acquires the geometric meaning of evaluating JLO cycle at the identity. Indeed such a formula was used to show the deformation invariance of the index.
On the other hand, in our context, our formula for the dimension is independent of an underlying supersymmetric structure. When we consider unitary cocycles associated with charges localizable in a bounded region, the dimension varies, but remains an integer.
If we consider localized endomorphisms and super-KMS functional, namely if we consider the index associated with a covariance cocycle and a super-KMS functional, we should expect our index formula to be read geometrically by the JLO cyclic cocycle. We will explain this point at the end of this paper, as it can certainly give further insight. But the present picture is too primitive to be directly applicable because super-KMS functionals fail to exist when the spacetime is non-compact in the most natural situation when the functional is translation invariant and space translations act in an asymptotically abelian fashion . This drawback is entirely caused by the assumption that the super-KMS functional is bounded. The structure associated with unbounded super-KMS functional is presently under investigation.
## 1 First properties of holomorphic cocycles.
In this section we begin to study holomorphic cocycles and give first formulae for the dimension.
### 1.1 Index formulae. Factor case.
Here we recall and extend results in , Section 1. We shall show how to obtain a holomorphic formula for the dimension of a sector that assumes neither a perfect symmetry group nor a PCT anti-automorphism. The reader is however assumed to have read the above quoted reference.
Let $``$ be an infinite factor and denote by End$`()`$ the (injective, normal, unital)<sup>3</sup><sup>3</sup>3The ‘index’ or ‘dimension’ terminology is here interchangeable; analogously, the Fredholm index of an isometry is the dimension of its cokernel. endomorphisms of $``$ with finite Jones index and Sect$`()`$ the sectors of $``$, namely the equivalence classes of End$`()`$ modulo inner automorphisms of $``$. End$`()`$ is a tensor category where the tensor product is the composition, cf. . The intertwiner space $`(\rho ,\sigma )`$ between objects $`\rho ,\sigma \mathrm{End}()`$ is defined as
$$(\rho ,\sigma )\{T:T\rho (X)=\sigma (X)T,X\}.$$
Sect$`()`$ is endowed with a natural conjugation. $`\overline{\rho }\mathrm{End}()`$ is a conjugate of $`\rho `$ iff the conjugate equation holds true: there exist multiples of isometries $`R(\iota ,\overline{\rho }\rho )`$ and $`\overline{R}(\iota ,\rho \overline{\rho })`$, that we normalize with $`R=\overline{R}`$, such that
$$R^{}\overline{\rho }(\overline{R})=1,\overline{R}^{}\rho (R)=1.$$
(9)
The minimum
$$d(\rho )\mathrm{min}R\overline{R}$$
over all possible choices of $`R`$ and $`\overline{R}`$ is the intrinsic dimension of $`\rho `$. A pair $`R_\rho `$, $`\overline{R}_\rho `$ where the minimum is attained exists and coincides with a standard solution as discussed in . It turns out that $`d(\rho )=d_{an}(\rho )`$, the analytical dimension defined as $`d_{an}(\rho )=\sqrt{[:\rho ()]}`$, the square root of the minimal index $`[:\rho ()]`$ (Jones-Kosaki index with respect to the minimal expectation) .
For each $`T(\rho _1,\rho _2)`$, the conjugate arrow $`T^{}(\overline{\rho }_1,\overline{\rho }_2)`$ is defined by
$$T^{}=\overline{\rho }_2(\overline{R}_{\rho _1}^{}T^{})R_{\rho _2}$$
where $`R_{\rho _i}`$ and $`\overline{R}_{\rho _i}`$ give a standard solution for the conjugate equation defining the conjugate $`\overline{\rho }_i`$. The map $`TT^{}`$ is anti-linear and satisfies natural properties, see .
Fix a normal faithful state $`\phi `$ of $``$ and let $`\sigma ^\phi `$ be the modular group of $`\phi `$. As is well known, $`\sigma ^\phi `$ satisfies the KMS condition with respect to $`\phi `$ at inverse temperature $`1`$, namely, setting $`\alpha _t=\sigma _t^\phi `$, the relation (12) holds with $`\beta =1`$. For a fixed $`\rho `$, let $`u(\rho ,)`$ be a unitary $`\sigma ^\phi `$-cocycle (i.e. $`u(\rho ,t)`$ is a unitary in $``$ and $`u(\rho ,t+s)=u(\rho ,t)\sigma _t^\phi (u(\rho ,s))`$, $`t,s`$ ) such that
$$\mathrm{Ad}u(\rho ,t)\sigma _t^\phi \rho =\rho \sigma _t^\phi ,$$
that is $`u(\rho ,t)(\rho _t,\rho )`$, where $`\rho _t\sigma _t^\phi \rho \sigma _t^\phi `$. Note that $`u(\rho ,)`$ is not not assumed to be continuous. Once standard $`R_\rho `$ and $`\overline{R}_\rho `$ are given, we assume the corresponding operators for the conjugate equation for $`\rho _t`$ and $`\overline{\rho }_t\sigma _t^\phi \overline{\rho }\sigma _t^\phi `$ to be given by $`R_{\rho _t}=\sigma _t^\phi (R_\rho )`$ and $`\overline{R}_{\rho _t}=\sigma _t^\phi (\overline{R}_\rho )`$.
Then $`u(\rho ,t)^{}(\overline{\rho }_t,\overline{\rho })`$ is given by
$$u(\rho ,t)^{}=\overline{\rho }(\overline{R}_{\rho _t}^{}u(\rho ,t)^{})R_\rho =\overline{\rho }(\sigma _t^\phi (\overline{R}_\rho ^{})u(\rho ,t)^{})R_\rho .$$
If $`\rho `$ is irreducible, the choice of $`R_\rho `$ and $`\overline{R}_\rho `$ is unique up to a phase, therefore $`u(\rho ,t)^{}`$ is uniquely defined. This holds in more generality, if $`\rho `$ is reducible.
###### Proposition 1.1.
$`u(\rho ,t)^{}`$ is well defined, namely it does not depend on the choice of $`R_\rho `$ and $`\overline{R}_\rho `$ giving a solution for the conjugate equation for $`\rho `$ and $`\overline{\rho }`$.
###### Proof.
Let $`R_\rho `$ and $`\overline{R}_\rho `$ be a standard solution. If $`R_\rho ^{}`$ and $`\overline{R}_\rho ^{}`$ is another solution of the conjugate equation, then $`R_\rho ^{}=\overline{\rho }(v)R_\rho `$ and $`\overline{R}_\rho ^{}=v^1\overline{R}_\rho `$ for some invertible $`v(\rho ,\rho )`$ , hence the conjugate of $`u(\rho ,t)`$ with respect to $`R_\rho ^{}`$ and $`\overline{R}_\rho ^{}`$ is given by
$$\begin{array}{c}\overline{\rho }(\sigma _{}^{\phi }{}_{t}{}^{}(\overline{R}_\rho ^{})\sigma _t^\phi (v^1)u(\rho ,t)^{})\overline{\rho }(v)R_\rho =\overline{\rho }(\sigma _t^\phi (R_\rho ^{})u(\rho ,t)^{}\sigma _t^{\phi _\rho }(v^1))\overline{\rho }(v)R_\rho \hfill \\ \hfill =\overline{\rho }(\sigma _t^\phi (R_\rho ^{})u(\rho ,t)^{}v^1)\overline{\rho }(v)R_\rho =\overline{\rho }(\sigma _t^\phi (\overline{R}_\rho {}_{}{}^{}))u(\rho ,t)^{})R_\rho =u(\rho ,t)^{},\end{array}$$
(10)
where $`\sigma ^{\phi _\rho }`$ is the modular group of $`\phi \mathrm{\Phi }_\rho `$ , so that $`\sigma _t^{\phi _\rho }(v)=v`$ because the minimal left inverse $`\mathrm{\Phi }_\rho `$ of $`\rho `$ is tracial on $`(\overline{\rho },\overline{\rho })`$. $`\mathrm{}`$
###### Proposition 1.2.
Let $`\rho `$ be irreducible. We have
$$d(\rho )^2=\underset{ti}{\text{anal.cont. }}\phi (u(\rho ,t))\phi (u(\rho ,t)^{}).$$
If $`u(\rho ,)`$ is weakly continuous, then $`u(\rho ,)`$ and $`u(\rho ,)^{}`$ are holomorphic (see below) in the state $`\phi `$, and
$$d(\rho )^2=\underset{ti}{\text{anal.cont. }}\phi (u(\rho ,t))\underset{ti}{\text{anal.cont. }}\phi (u(\rho ,t)^{}).$$
###### Proof.
We may suppose that $`d(\rho )<\mathrm{}`$. Set $`\phi _\rho \phi \mathrm{\Phi }_\rho `$, where $`\mathrm{\Phi }_\rho `$ is the minimal left inverse of $`\rho `$ as before. Since the Connes Radon-Nikodym cocycle $`(D\phi _\rho :D\phi )`$ is a covariance cocycle for $`\rho `$ , there exists a one-dimensional character $`\chi `$ of $``$ such that
$$u(\rho ,t)=\chi (t)d(\rho )^{it}(D\phi _\rho :D\phi )_t,$$
(both cocycles intertwine $`\sigma ^\phi `$ and $`\sigma ^{\phi _\rho }`$). As shown in , $`(d(\rho )^{it}(D\phi _\rho :D\phi )_t)^{}=d(\rho )^{it}(D\phi _{\overline{\rho }}:D\phi )_t`$
$$u(\rho ,t)^{}=\overline{\chi (t)}(d(\rho )^{it}(D\phi _\rho :D\phi )_t)^{}=\overline{\chi (t)}d(\rho )^{it}(D\phi _{\overline{\rho }}:D\phi )_t,$$
Moreover the Connes cocycle is holomorphic and $`\underset{ti}{\text{anal.cont. }}\phi ((D\phi _\rho :D\phi )_t)=\phi _\rho (1)=1`$ , therefore
$$\begin{array}{c}\underset{ti}{\text{anal.cont. }}\phi (u(\rho ,t))\phi (u(\rho ,t)^{})\hfill \\ \hfill =\underset{ti}{\text{anal.cont. }}d(\rho )^{2it}\phi ((D\phi _\rho :D\phi )_t)\phi ((D\phi _{\overline{\rho }}:D\phi )_t)=d(\rho )^2.\end{array}$$
(11)
If $`u(\rho ,)`$ is continuous, then $`\chi `$ is continuous. So $`\chi `$ extends to an entire function, hence $`u(\rho ,)`$ is holomorphic, and the second formula in the statement follows from the first one. $`\mathrm{}`$
Remark. Let $`𝒯\text{End}()`$ be a C tensor category as before, and $`u(\rho ,t)`$ a two-variable cocycle. Setting $`d_\phi (\rho )=d(u_\rho )\underset{ti}{\text{anal.cont. }}\phi (u_\rho (t))`$ (see below) the arguments in show that
$`d_\phi (\rho _1\rho _2)`$ $`=d_\phi (\rho _1)+d_\phi (\rho _2),`$
$`d_\phi (\rho _1\rho _2)`$ $`=d_\phi (\rho _1)d_\phi (\rho _2),\rho _1,\rho _2𝒯.`$
In particular, if $`𝒯`$ is rational, namely there are only finitely many inequivalent irreducible objects, the usual application of the Perron-Frobenious theorem entails $`d(\rho )=d_\phi (\rho )`$ for all objects of $`𝒯`$ (no chemical potential, see Sect. 2).
#### 1.1.1 Case of a non-full C tensor sub-category.
Now let $`𝒯`$ be a C tensor category with conjugates contained in End$`()`$, thus the objects of $`𝒯`$ are finite-index endomorphisms of $``$, but we do not assume that $`𝒯`$ is a full sub-category of End$`()`$, namely the intertwiner spaces $`(\rho ,\sigma )`$ in $`𝒯`$ can be strictly contained in the corresponding intertwiner spaces in End$`()`$.
Assume that the modular group $`\sigma ^\phi `$ gives an action of $``$ on $`𝒯`$, that is $`\rho _t𝒯`$, for all $`t`$, $`\rho 𝒯`$, and $`\sigma _t^\phi ((\rho ,\rho ^{}))=(\rho _t,\rho _t^{})`$ if $`\rho ,\rho ^{}𝒯`$ are objects of $`𝒯`$.
Recall that $`u`$ is a two-variable unitary cocycle for the above action if, for each fixed object $`\rho 𝒯`$, $`u(\rho ,)`$ is a unitary $`\sigma ^\phi `$-cocycle as above and, for each fixed $`t`$, $`u(\rho ,)^{}`$ is cocycle with respect to $`\sigma _t^\phi `$, namely
$$u(\rho \sigma ,t)=\rho (u(\sigma ,t))u(\rho ,t),\rho ,\sigma 𝒯,$$
and
$$Tu(\rho ,t)=u(\sigma ,t)\sigma _t^\phi (T),\rho ,\sigma 𝒯,T(\rho ,\sigma ).$$
Note that if $`u`$ is a two variable cocycle also for the full tensor subcategory of End$`()`$ with the same objects of $`𝒯`$, then $`u(\rho ,t)^{}`$ defined there coincides with $`u(\rho ,t)^{}`$ defined in $`𝒯`$ cf. , Propositions 1.5 and A.2.
###### Corollary 1.3.
With the above notations, if $`u(\rho ,t)`$ is a weakly continuous unitary two-variable cocycle for the action of $``$ on $`𝒯`$ given by $`\sigma ^\phi `$, then
$$d(\rho )\sqrt{\underset{ti}{\text{anal.cont. }}\phi (u(\rho ,t))\phi (u(\overline{\rho },t))}$$
for all irreducible $`\rho 𝒯`$ . Here $`d(\rho )`$ is the intrinsic dimension of $`\rho `$ as an object of $`𝒯`$.
###### Proof.
We have $`u(\rho ,t)^{}=u(\overline{\rho },t)`$ (the conjugate map is the one associated with $`𝒯`$), hence the above Proposition 1.1 applies, provided $`\rho `$ is irreducible in End$`()`$. If $`\rho `$ is reducible in End$`()`$ we have $`u(\rho ,t)=z(t)d_{an}(\rho )^{it}(D\phi _\rho :D\phi )_t`$ with $`z(t)(\rho ,\rho )`$. By using the tracial property of $`\mathrm{\Phi }_\rho `$ it is easy to check that $`z`$ is a one-parameter group of unitaries in the finite-dimensional algebra $`(\rho ,\rho )`$, and therefore can be diagonalized. If $`\rho =_i\rho _i`$ is a decomposition of $`\rho `$ into irreducibles (with eigen-projections of $`z`$) in End$`()`$, we have a corresponding decomposition of $`d_{an}(\rho )^{it}(D\phi _\rho :D\phi )_t`$ as direct sum of the $`d_{an}(\rho _i)^{it}(D\phi _{\rho _i}:D\phi )_t`$ thus $`d_\phi (u(\rho _i,))=\mathrm{}_id_{an}(\rho _i)`$ with $`\mathrm{}_i>0`$, hence
$$d_\phi (u_\rho )d_\phi (u_{\overline{\rho }})=\underset{i,j}{}\mathrm{}_id_{an}(\rho _i)\mathrm{}_jd_{an}(\rho _j)\underset{i,j}{}d_{an}(\rho _i)d_{an}(\rho _j)=d_{an}(\rho )^2d(\rho )^2.$$
$`\mathrm{}`$
We now recall that the following holds.
###### Proposition 1.4.
. In the setting of Prop. 1.1, if there exists an $`\phi `$-preserving anti-automorphism $`j`$ of $``$ inducing an anti-automorphism of $`𝒯`$ such that $`j\rho j=\overline{\rho }`$ and $`j(u(\rho ,t))=u(\overline{\rho },t)`$, then $`d(\rho )=\underset{ti}{\text{anal.cont. }}\phi (u(\rho ,t))`$ for all objects $`\rho 𝒯`$.
###### Proof.
See , Prop.1.7. $`\mathrm{}`$
### 1.2 The holomorphic dimension in the C-case.
In this section we give a first look at the structure that emerges in the C context, in analogy to what studied in the previous section in the setting of factors. Here we assume from the start that a holomorphic dimension is definable, postponing the more relevant derivation of the holomorphic property and the analysis of the chemical potential to subsequent sections.
Let $`𝔄`$ be a unital C-algebra and $`\alpha `$ a one-parameter automorphism group of $`𝔄`$. A linear functional $`\phi 𝔄^{}`$ is said to be a KMS functional with respect to $`\alpha `$ at inverse temperature $`\beta >0`$ if for any given $`a,b𝔄`$ there exists a function $`F_{a,b}A(S_\beta )`$ such that
$`F_{a,b}(t)`$ $`=\phi (\alpha _t(a)b)`$ (12)
$`F_{a,b}(t+i\beta )`$ $`=\phi (b\alpha _t(a))`$ (13)
Here $`S_\beta `$ is the strip $`\{0<\text{Im}z<\beta \}`$ and $`A(S_\beta )`$ is the algebra of functions analytic in $`S_\beta `$, bounded and continuous on the closure of $`S_\beta `$. We do not assume $`\alpha `$ to be pointwise norm continuous, nonetheless a weaker continuity property follows from the KMS condition. Note that a KMS functional $`\phi `$ is $`\alpha `$-invariant.
Let now $`u𝔄`$ be a unitary cocycle with respect to $`\alpha `$, namely $`tu(t)𝔄`$ is a map taking values in the unitaries of $`𝔄`$ satisfying the equation
$$u(t+s)=u(t)\alpha _t(u(s)).$$
We shall say that the cocycle $`u`$ is holomorphic, in the functional $`\phi `$, if the function $`t\phi (u(t))`$ is the boundary value of a function in $`A(S_\beta )`$.
If $`u`$ is holomorphic in the state $`\phi `$, we define the holomorphic dimension of the cocycle $`u`$ (with respect to $`\phi `$) by
$$d_\phi (u)=\underset{ti\beta }{\text{anal.cont. }}\phi (u(t)).$$
As we shall see, in our context $`d_\phi (u)`$ will be a positive number related to a (noncommutative) relative index.
Clearly, for a given dynamics $`\alpha `$, $`d_\phi (u)`$ may depend on the KMS state $`\phi `$. We shall sometimes rescale the “time parameter” to make the inverse temperature $`\beta =1`$. If $`u`$ is not holomorphic we write $`d_\phi (u)=+\mathrm{}`$.
Let $`\rho `$ be an endomorphism of $`𝔄`$. We shall say that $`\rho `$ is covariant if there exists a $`\alpha `$-cocycle of unitaries $`u(\rho ,t)𝔄`$ such that
$$\alpha _t\rho \alpha _t=\text{Ad}u(\rho ,t)^{}\rho .$$
(14)
We shall say that $`\rho `$ has finite holomorphic dimension (with respect to the KMS state $`\phi `$) if it is covariant and there exists a covariance cocycle $`u(\rho ,)`$ as above with finite holomorphic dimension. Note that, if $`𝔄`$ has trivial centre and $`\rho `$ is irreducible, $`u(\rho ,)`$ is unique up to a phase, that doesn’t alter the finite-dimensional property of $`u(\rho ,)`$, provided such a phase is chosen to be a continuous character of $``$.
In the rest of this section we study whether endomorphisms of $`𝔄`$ with finite holomorphic dimension extend to $`\pi _\phi (𝔄)^{\prime \prime }`$.
In the following we identify $`𝔄`$ with its image $`\pi _\phi (𝔄)`$ and suppress the suffix $`\phi `$.
###### Lemma 1.5.
Let $`𝔄`$ be a C-algebra acting on a Hilbert space $``$ and $`\rho `$ an endomorphism of $`𝔄`$. Let $`\xi `$ be a cyclic separating vector for $`=𝔄^{\prime \prime }`$ and $`\phi =(\xi ,\xi )`$. If $`\phi \rho `$ extends to a normal faithful positive functional of $``$, then $`\rho `$ extends to a normal endomorphism of $``$.
###### Proof.
Let $`\eta `$ be a cyclic separating vector such that $`(\rho (a)\xi ,\xi )=(a\eta ,\eta )`$, $`a𝔄`$, and let $`V`$ be the isometry of $``$ given by
$$Va\eta =\rho (a)\xi ,a𝔄.$$
(15)
The final projection of $`V`$ is given by
$$e=VV^{}=\overline{\rho (𝔄)\xi }\rho (𝔄)^{}$$
(16)
thus $`xVxV^{}`$ is a homomorphism of $``$ onto $`\rho (𝔄)^{\prime \prime }e`$ and
$$VaV^{}=\rho (a)e,a𝔄.$$
(17)
Now the central support of $`e`$ in $`\rho (𝔄)^{\prime \prime }`$ is $`1`$ as $`\overline{\rho (𝔄)^{}\xi }\overline{^{}\xi }=`$, hence if $`x`$ there exists a unique $`\rho (x)\rho (𝔄)^{\prime \prime }`$ such that $`\rho (x)e=VxV^{}`$, providing an extension of $`\rho `$ to $``$. As $`\eta `$ is separating, $`\rho `$ is an isomorphism. $`\mathrm{}`$
Now, as in the factor case, End$`(𝔄)`$ is the tensor category whose objects $`\rho ,\sigma ,\mathrm{}`$ are the endomorphisms of $`𝔄`$: the monoidal product $`\rho \sigma =\rho \sigma `$ is given by the composition of maps, while the intertwiner space $`(\rho ,\sigma )`$ is given by $`\{T𝔄:T\rho (a)=\sigma (a)T,a𝔄\}`$. The tensor product of intertwiners is also defined in a natural fashion, see e.g. . The conjugate equation is defined as in the previous section. The *intrinsic dimension* $`d(\rho )`$, and the conjugation on arrows are defined as well, in fact all these notions make sense for a a general tensor C-category, . The following proposition is a special case of results in . We state it in the particular case needed for our applications.
###### Proposition 1.6.
(). Let $`𝔄`$ be a unital C-algebra, acting on a Hilbert space, with $`=𝔄^{\prime \prime }`$ a factor. Let $`𝒯`$ be a tensor category with conjugates and subobjects of endomorphisms of $`𝔄`$ admitting a unitary braid group symmetry. Suppose that every endomorphism $`\rho 𝒯`$ extends to a normal endomorphism of $`\widehat{\rho }`$. Then
$$d(\widehat{\rho })=d(\rho )$$
where $`d(\widehat{\rho })=d_{an}(\widehat{\rho })`$, i.e. $`d(\widehat{\rho })^2`$ is the minimal index $`[:\widehat{\rho }()]`$, and $`d(\rho )`$ is the intrinsic dimension of $`\rho `$ in $`𝒯`$.
###### Proof.
The map $`\rho \widehat{\rho }`$ is a functor of C tensor categories from $`𝒯`$ to a sub-tensor category of End$`()`$, hence $`d(\widehat{\rho })d(\rho )`$. As $`𝒯`$ has a unitary braiding, every real object $`\sigma 𝒯`$ is amenable Th. 5.31, thus $`d(\sigma )=m^\sigma `$, where $`m^\sigma `$ is the $`\mathrm{}^2`$ norm of the fusion matrix $`m^\sigma `$ associated with $`\sigma `$, therefore $`d(\sigma )=m^\sigma m^{\widehat{\sigma }}d(\widehat{\sigma })d(\sigma )`$, thus $`d(\widehat{\sigma })=d(\sigma )`$. For any $`\rho 𝒯`$, the object $`\sigma =\rho \overline{\rho }`$ is real hence, by the multiplicativity of the dimension, $`d(\widehat{\rho })=d(\rho )`$. $`\mathrm{}`$
#### 1.2.1 Case of a unique KMS state.
We now restrict our attention to the case of a unique KMS functional. This is done more with an illustrative intent, rather than for later applications, where we shall treat a more general context.
###### Proposition 1.7.
Suppose $`\phi `$ is the unique KMS functional for $`\alpha `$. Then a covariant endomorphism $`\rho `$ of $`𝔄`$ with finite holomorphic dimension has a normal extension to $`𝔄^{\prime \prime }`$.
###### Proof.
Let $`u=u(\rho ,)`$ a holomorphic unitary covariance $`\alpha `$-cocycle. As $`t\phi (u(t))`$ is continuous, the map $`tu(t)`$ is strongly continuous. To check this, note that by cocycle property it is enough to verify the continuity at $`t=0`$ because then the strong limit
$$u(s+t)=u(s)\alpha _s(u(t))0,\text{as}t0$$
(18)
due to the normality of $`\alpha _s`$.
Let then $`x`$ be a weak limit point of $`u(t)`$ as $`t0`$. Then $`x1`$ and
$$(x\xi _\phi ,\xi _\phi )=\underset{i\mathrm{}}{lim}(u(t_i)\xi _\phi ,\xi _\phi )=\underset{i\mathrm{}}{lim}\phi (u(t_i))=\phi (1)=(\xi _\phi ,\xi _\phi )$$
(19)
for some sequence $`t_i0`$. Thus $`x\xi _\phi =\xi _\phi `$ by the limit case of the Schwartz inequality, thus $`x=1`$ because $`\xi _\phi `$ is separating.
Therefore by Connes’ theorem there exists a normal faithful semifinite weight $`\phi _\rho `$ on $``$ with $`(D\phi _\rho ,D\phi )_t=u(t)`$ and, by the finite holomorphic dimension assumption, $`\phi _\rho (1)=d_\phi (u)\mathrm{}`$ so that $`\phi `$ is indeed a positive linear functional. Then $`\phi _\rho `$ is a KMS functional for its modular group $`\alpha _t^\rho =\text{Ad}u(t)\alpha _t`$ (we are setting $`\beta =1`$ here).
The functional on $`𝔄`$
$$\phi ^{}(a)=\phi _\rho (\rho (a)),a𝔄$$
(20)
is KMS with respect to $`\alpha `$
$$\begin{array}{c}\underset{ti}{\text{anal.cont. }}\phi ^{}(a\alpha _t(b))=\underset{ti}{\text{anal.cont. }}\phi _\rho (\rho (a)\rho (\alpha _t(b)))\hfill \\ \hfill =\underset{ti}{\text{anal.cont. }}\phi _\rho (\rho (a)\alpha _t^\rho (\rho (b)))=\phi _\rho (\rho (b)\rho (a))=\phi ^{}(ba).\end{array}$$
(21)
Hence, by the uniqueness of the KMS state,
$$\phi ^{}=\lambda \phi $$
(22)
on $`𝔄`$, for some $`\lambda >0`$. Thus $`\phi _\rho `$ is a normal faithful functional and $`\phi _\rho \rho `$ is normal too. Therefore the proof is completed by Lemma 1.5. $`\mathrm{}`$
To shorten notation, we shall often set
$$u_\rho =u(\rho ,).$$
###### Proposition 1.8.
Let $`\phi `$ be the unique KMS state as in Proposition 1.7. If $`\rho `$ and $`\overline{\rho }`$ are conjugate and $`d_\phi (u_\rho )<\mathrm{}`$, $`d_\phi (u_{\overline{\rho }})<\mathrm{}`$, then the extension of $`\rho `$ to $``$ has finite Jones index.
###### Proof.
By assumption and Proposition 1.7 both $`\rho `$ and $`\overline{\rho }`$ extend to $``$. As $`R_\rho `$ and $`\overline{R}_\rho `$ are also intertwiners on $``$ by weak continuity and the conjugate equation for $`\rho `$ and $`\overline{\rho }`$ is obviously satisfied on $``$, the extension of $`\rho `$ to $``$ has finite index. $`\mathrm{}`$
###### Proposition 1.9.
Suppose that $`\phi `$ is faithful and the unique KMS state for $`\alpha `$, as in Prop. 1.8. Let $`𝒯`$ be a tensor category with conjugates of endomorphisms of $`𝔄`$ and $`u(\rho ,t)𝔄`$ a unitary two-variable cocycle. If $`u_\rho `$ has finite holomorphic dimension for all irreducible objects $`\rho `$, then the intrinsic dimension $`d(\rho )`$ of $`\rho `$ is bounded by
$$d(\rho )\sqrt{d_\phi (u_\rho )d_\phi (u_{\overline{\rho }})}.$$
(23)
and equality holds if $`\rho `$ is irreducible and extending to $``$ is a full functor (thus $`\rho `$ is irreducible on $``$).
###### Proof.
By Lemma 1.8 $`\rho `$ extends to $``$ and has finite index. We are then in the case covered by Prop. 1.2 and Corollary 1.3. $`\mathrm{}`$
Motivated by the above Proposition, we define the geometric dimension $`d_{geo}(\rho )`$ as
$$d_{geo}(\rho )\sqrt{d_\phi (u_\rho )d_\phi (u_\rho ^{})}.$$
If $`\rho `$ is irreducible, $`d_{geo}(\rho )`$ does not depend on the choice of the covariance cocycle $`u(\rho ,t)`$, because, if we multiply $`u(\rho ,t)`$ by a phase $`\chi (t)`$, then $`u(\rho ,t)^{}`$ has to be replaced by $`\overline{\chi (t)}u(\rho ,t)^{}`$. Of course, since for a two-variable cocycle $`u(\rho ,t)`$ we have
$$u(\rho ,t)^{}=u(\overline{\rho },t),$$
in this case we also have
$$d_{geo}(\rho )=\sqrt{d_\phi (u_\rho )d_\phi (u_{\overline{\rho }})}.$$
Also, since $`u(\rho \overline{\rho },t)=\rho (u(\overline{\rho },t))u(\rho ,t)`$, we have
$$d_{geo}(\rho )=\sqrt{d_\phi (u_{\rho \overline{\rho }})}$$
if $`\rho `$ and $`\overline{\rho }`$ extend to $``$.
A priori $`d_{geo}(\rho )`$ might depend on the KMS functional $`\phi `$, but, as we shall see, in most interesting cases it will actually be independent of $`\phi `$.
#### 1.2.2 Graded KMS functionals: reduction to ordinary KMS states.
The above results extend to graded KMS functionals. Indeed the analysis of these functionals can be reduced to the case of ordinary KMS states.
Let $`𝔄`$ be a $`_2`$-graded unital C-algebra, namely $`𝔄`$ is a unital C-algebra equipped with an involutive automorphism $`\gamma `$ <sup>4</sup><sup>4</sup>4In this paper morphisms always commute with the -mapping and preserve the unit..
Given a graded one-parameter automorphism group $`\alpha `$ of $`𝔄`$ (i.e. one commuting with $`\gamma `$), a linear functional $`\phi 𝔄^{}`$ is said to be a graded KMS functional with respect to $`\alpha `$ at inverse temperature $`\beta >0`$ if for any given $`a,b𝔄`$ there exists a function $`F_{a,b}A(S_\beta )`$ such that
$`F_{a,b}(t)`$ $`=\phi (\alpha _t(a)b)`$ (24)
$`F_{a,b}(t+i\beta )`$ $`=\phi (\gamma (b)\alpha _t(a))`$ (25)
Note that a graded KMS functional $`\phi `$ is $`\gamma `$-invariant.
We recall the following.
###### Proposition 1.10.
() Let $`\phi `$ be a graded KMS functional of $`𝔄`$ for $`\alpha `$ and let $`\omega =|\phi |`$ be the modulus of $`\phi `$. Then $`\omega `$ is an ordinary KMS positive functional and $`\pi _\omega \gamma `$ extends to an inner automorphism of $`\pi _\omega (𝔄)^{\prime \prime }`$, implemented by a selfadjoint unitary $`\mathrm{\Gamma }`$ in the centralizer $`_\omega `$ of $`\omega `$. Moreover $`\phi `$ is proportional to $`\omega (\mathrm{\Gamma })`$.
###### Corollary 1.11.
If $`\phi `$ is the unique non-zero graded KMS functional (up to a phase) of $`𝔄`$, then $`\omega =|\phi |`$ is extremal KMS, i.e. $`=\pi _\omega (𝔄)^{\prime \prime }`$ is a factor.
###### Proof.
As usual $`𝔄`$ is identified with $`\pi _\omega (𝔄)`$. If $`Z()`$, there exist two non-zero projections $`z_1,z_2Z()`$ with sum $`1`$. Thus $`\phi (z_1),\phi (z_2)`$ are different graded KMS functionals on $`𝔄`$. This can be checked since both the extensions of $`\alpha `$ and $`\gamma `$ act trivially on $`Z()`$, and by usual approximation arguments. By the uniqueness assumption there exists a constant $`\lambda `$ with $`\phi (z_1a)=\lambda \phi (z_2a),a𝔄`$, then by continuity the same equality holds for $`a`$, thus $`\phi (z_1a)=\phi (z_1z_1a)=\lambda \phi (z_2z_1a)=0`$. Analogously $`\phi (z_2)=0`$, thus $`\phi =0`$. $`\mathrm{}`$
For our purposes Proposition 1.10 allows us to consider ordinary KMS states instead of general graded KMS functionals.
Let $`\rho `$ be a graded endomorphism of $`𝔄`$, namely an endomorphism of $`𝔄`$ commuting with $`\gamma `$. In this graded context, we shall say that $`\rho `$ is covariant if there exists a covariance $`\alpha `$-cocycle of unitaries $`u(t)𝔄`$ such that $`\gamma (u(t))=u(t)`$.
In the following we again identify $`𝔄`$ with its image $`\pi _\omega (𝔄)`$.
###### Lemma 1.12.
Let $`\phi `$ be a graded KMS functional for $`\alpha `$ and $`\rho `$ a graded covariant endomorphism of $`𝔄`$ as above. If $``$ is a factor and $`\rho `$ extends to a finite index irreducible endomorphism of $``$, then $`\rho `$ has finite holomorphic dimension, both with respect to $`\omega `$ and $`\phi `$. Indeed, if $`\omega (u())`$ is continuous,
$$d_\phi (u_\rho )=\pm d_\omega (u_\rho ),$$
namely $`\phi (1)d_\omega (u)=\pm \underset{ti\beta }{\text{anal.cont. }}\phi (u(t))`$.
###### Proof.
The holomorphic property of $`u`$ is a direct consequence of the holomorphic property of the Connes cocycle because $`u`$ is indeed a Connes Radon-Nikodym cocycle, up to phase, with respect to two bounded positive normal functionals of $``$ (see and the previous section).
Note now that, since the extension of $`\rho `$ to $``$ (still denoted by $`\rho `$) commutes with $`\gamma =\mathrm{Ad}\mathrm{\Gamma }`$, we have $`\rho (\mathrm{\Gamma })\mathrm{\Gamma }^{}\rho ()^{}=`$, thus $`\rho (\mathrm{\Gamma })=\pm \mathrm{\Gamma }`$ because $`\mathrm{\Gamma }`$ is self-adjoint, hence $`\mathrm{\Phi }_\rho (\mathrm{\Gamma })=\pm \mathrm{\Gamma }`$, where $`\mathrm{\Phi }_\rho `$ is the (unique) left inverse of $`\rho `$.
To check eq. (1.12), recall that, by the holomorphic properties of Connes cocycles, we have (see ):
$$\underset{ti\beta }{\text{anal.cont. }}\omega (Xu(t))=d_\omega (u)\omega \mathrm{\Phi }_\rho (X),X.$$
(26)
As above we have the polar decomposition $`\phi =\omega (\mathrm{\Gamma })`$. Then
$$\underset{ti\beta }{\text{anal.cont. }}\phi (u(t))=\underset{ti\beta }{\text{anal.cont. }}\omega (\mathrm{\Gamma }u(t))=d_\omega (u)\omega \mathrm{\Phi }_\rho (\mathrm{\Gamma })=\pm d_\omega (u)\omega (\mathrm{\Gamma })=\pm \phi (1)d_\omega (u),$$
namely the holomorphic dimension with respect to $`\phi `$ coincides with the holomorphic dimension with respect to $`\omega `$, up to a sign. $`\mathrm{}`$
Because of the above Lemma 1.12, it is more convenient to define the holomorphic dimension $`d_\phi (u)`$ directly with respect to the modulus $`\omega `$ of $`\phi `$.
Before concluding this section, we recall that the interest in (bounded) graded KMS functionals is limited by the following no-go theorem.
###### Proposition 1.13.
*()* Let $`\phi `$ be a graded KMS functional of $`𝔄`$ with respect to $`\alpha `$ as above. If there exists a $`\phi `$-asymptotically abelian sequence $`\beta _n\mathrm{Aut}(𝔄)`$, then $`\gamma =\iota `$ and $`\phi `$ is an ordinary KMS functional.
Here the $`\beta _n`$’s commute with the grading and the $`\phi `$-asymptotically abelianness means that $`\phi \beta _n=\phi `$ and $`\phi (c[\beta _n(a),b])0`$ for all $`a,b,c𝔄`$, where the commutator is a graded commutator.
#### 1.2.3 Table of dimensions.
Before concluding this section we display the following table that summarizes the various notions of dimension we are dealing with.
| Dimension | Definition | Context |
| --- | --- | --- |
| Intrinsic | $`d\left(\rho \right)=\sqrt{R_\rho \overline{R}_\rho }`$, (standard $`R_\rho ,\overline{R}_\rho `$) | Tensor C-categories |
| Analytical | $`d_{an}\left(\rho \right)=\sqrt{\left[:\rho \left(\right)\right]}`$ | Subfactors |
| Statistical | $`d_{DHR}\left(\rho \right)=\left|\mathrm{\Phi }_\rho \left(\epsilon _\rho \right)\right|^1`$, ($`\epsilon _\rho `$ stat. operator) | QFT, localized endomorphisms |
| Holomorphic | $`d_\phi \left(u\right)=\underset{ti\beta }{\text{anal.cont. }}\phi \left(u\left(t\right)\right)`$ | Unitary cocycles |
| Geometric | $`d_{geo}\left(\rho \right)=\sqrt{d_\phi \left(u_\rho \right)d_\phi \left(u_{\overline{\rho }}\right)}`$ | Covariant endomorphisms |
Here $`[:\rho ()]`$ denotes the minimal index, namely the Jones index with respect to the minimal expectation. We have omitted the notion of minimal dimension $`d_{min}(\rho )\text{min}\sqrt{R_\rho \overline{R}_\rho }`$, in the context of C-tensor categories, as it turns out to coincide with the intrinsic dimension .
## 2 The chemical potential in Quantum Field Theory.
In this section we examine certain aspects of the chemical potential for thermal states in Quantum Field Theory. Our discussion will rely on basic results as the description of the chemical potential in terms of extensions of KMS states and the construction of the field net and the gauge group . Together with certain results for tensor categories , our analysis will show a splitting of the incremental free energy into an absolute part, that depends only on the charge and not on the state, and a part which is asymmetric with respect to the charge conjugation; this indeed represents the chemical potential that labels the equilibrium states.
Let $`𝕄`$ be the Minkowski spacetime $`^{d+1}`$, with $`d2`$, and $`𝔄`$ a net of local observable von Neumann algebras on $`𝕄`$, namely we have an inclusion preserving map
$$𝒪𝔄(𝒪)$$
from the set $`𝒦`$ of (open, non-empty) double cones of $`𝕄`$ to von Neumann algebras $`𝔄(𝒪)`$ on a Hilbert space $``$. If $`E𝕄`$ is arbitrary, we set $`𝔄(E)`$ for the C-algebra generated by the von Neumann algebras $`𝔄(𝒪)`$ as $`𝒪𝒦`$ varies $`𝒪E`$ ($`𝔄(E)=`$ if the interior of $`E`$ is empty) and denote by $`𝔄𝔄(𝕄)`$ the quasi-local C-algebra. We denote the quasi-local C-algebra and the net itself by the same symbol, but this should not create confusion. We shall also denote by $`𝒜(E)=𝔄(E)^{\prime \prime }`$ the von Neumann algebra generated by $`𝔄(E)`$ (of course $`𝔄(𝒪)=𝒜(𝒪)`$ if $`𝒪𝒦`$).
We assume the following properties<sup>5</sup><sup>5</sup>5These properties automatically hold, in particular, in a Wightman theory . Here we will also consider the case of a reducible net.:
Additivity: If $`𝒪,𝒪_1,\mathrm{},𝒪_n`$ are double cones and $`𝒪_1\mathrm{}𝒪_n𝒪`$, then $`𝔄(𝒪_1)\mathrm{}𝔄(𝒪_n)𝔄(𝒪)`$.
Here and in the following, the lattice symbol $``$ denotes the von Neumann algebra generated.
Wedge duality (or essential duality): If $`W𝕄`$ is a wedge region (namely a Poincaré translate of the region $`\{x𝕄:x_1>x_0\}`$), then
$$𝒜(W^{})=𝒜(W)^c.$$
Here $`𝒩^c`$ denotes the relative commutant in the von Neumann algebra $`𝔄^{\prime \prime }`$, namely $`𝒩^c𝒩^{}`$, and $`W^{}`$ denotes the spacelike complement of $`W`$.
In particular the net $`𝔄`$ is local, namely $`𝔄(𝒪_1)`$ and $`𝔄(𝒪_2)`$ commute if the double cones $`𝒪_1`$ and $`𝒪_2`$ are space-like separated. As in the case of irreducible nets, one may consider the dual net, here defined as
$$𝔄^d(𝒪)𝔄(𝒪^{})^c,𝒪𝒦,$$
and show the following, :
###### Proposition 2.1.
$`𝔄^d`$ satisfies Haag duality, by which we mean here that
$$𝔄^d(𝒪)=𝔄^d(𝒪^{})^c,𝒪𝒦,$$
where $`𝔄^d(𝒪^{})`$ is the C-algebra associated to $`𝒪^{}`$ in the net $`𝔄^d`$. In particular $`𝔄^d`$ is local.
Moreover $`𝔄`$ and $`𝔄^d`$ have the same weak closure:
$$𝔄^{\prime \prime }=(𝔄^d)^{\prime \prime }=.$$
###### Proof.
Because of additivity and wedge duality one can write
$$𝔄^d(𝒪)=\underset{W𝒪}{}𝒜(W)$$
(27)
(intersection over all wedges containing $`𝒪`$).
If $`𝒪_1𝒦`$ is spacelike separated from $`𝒪𝒦`$ there is a wedge $`W`$ with $`W𝒪`$ and $`W^{}𝒪_1`$, hence $`𝔄^d`$ is local. As $`𝔄^d`$ extends $`𝔄`$ and is local, wedge duality must hold for $`𝔄^d`$ too, therefore $`𝔄^d(W)=𝔄(W)`$ for all wedges $`W`$. In particular the global von Neumann algebra associated with $`𝔄^d`$ coincides with the one associated with $`𝔄`$:
$$(\underset{𝒪}{}𝔄^d(𝒪))^{\prime \prime }=(\underset{W}{}𝔄^d(W))^{\prime \prime }=(\underset{W}{}𝔄(W))^{\prime \prime }=.$$
Analogously, it follows from formula 27 that $`𝔄(𝒪^{})^{\prime \prime }=𝔄^d(𝒪^{})^{\prime \prime }`$, namely $`𝔄^d(𝒪)=𝔄^d(𝒪^{})^c`$, that is to say Haag duality holds for $`𝔄^d`$. $`\mathrm{}`$
Note however that, since $``$ is not a type I factor in general, $`𝔄(𝒪^{})^{\prime \prime }`$ may be non-normal in $``$, namely
$$𝔄^d(𝒪)^c=𝔄(𝒪^{})^{cc}$$
may be strictly larger than $`𝔄(𝒪^{})^{\prime \prime }`$ if $`𝒪𝒦`$.
Translation covariance: There exists a unitary representation $`U`$ of $`^{d+1}`$ making $`𝔄`$ covariant:
$$\tau _x(𝔄(𝒪))=𝔄(𝒪+x),x^{d+1},𝒪𝒦,$$
where we have set $`\tau _x\mathrm{Ad}U(x).`$
Properly infiniteness and Borchers property B <sup>6</sup><sup>6</sup>6In the vacuum representation these properties follows by positivity of the energy, see .: If $`𝒪𝒦`$, then $`𝔄(𝒪)`$ is a properly infinite von Neumann algebra. If $`𝒪,\stackrel{~}{𝒪}`$ are double cones and $`𝒪+x\stackrel{~}{𝒪}`$ for $`x`$ in a neighborhood of $`0`$ in $`^{d+1}`$, then every non-zero projection $`E𝔄(𝒪)`$ is equivalent to $`1`$ in $`𝔄(\stackrel{~}{𝒪})`$.
Factoriality: $`𝔄^{\prime \prime }`$ is a factor.
A localized endomorphism $`\rho `$ of $`𝔄`$ is an endomorphism of $`𝔄`$ such that $`\rho |_{𝔄(𝒪^{})}=\mathrm{id}|_{𝔄(𝒪^{})}`$ for some $`𝒪𝒦`$. Two localized endomorphisms $`\rho ,\rho ^{}`$ are equivalent ($`\rho \rho ^{}`$) if there is a unitary $`u`$ such that $`\rho ^{}=\mathrm{Ad}u\rho `$. In the DHR theory $`=B()`$, namely $`𝔄`$ is irreducible, but most of what we are saying holds in the reducible case as well.
A localized endomorphism $`\rho `$ is translation covariant if there exists a $`\tau `$-cocycle of unitaries $`u(\rho ,x)`$ such that
$$\mathrm{Ad}u(\rho ,x)\rho _x=\rho ,x^{d+1},$$
(28)
where $`\rho _x\tau _x\rho \tau _x`$.
The equivalence classes of irreducible, translation covariant, localized endomorphisms are the superselection sectors of $`𝔄`$.
The translation covariant endomorphisms of $`𝔄`$ form a tensor category. If $`T(\rho ,\rho ^{})`$ is an intertwiner between $`\rho `$, $`\rho ^{}`$, namely $`T`$ and $`T\rho (X)=\rho ^{}(X)T`$ for all $`X𝔄`$, then, as an immediate consequence of the Haag duality property for $`𝔄^d`$, we have $`T𝔄^d(𝒪)`$ if $`𝒪𝒦`$ contains the localization regions of both $`\rho `$ and $`\rho ^{}`$. Now the unitary intertwiners changing the localization region of $`\rho `$ (in particular the unitaries $`u(\rho ,x)`$ above) are used to define Roberts cohomology . In particular the endomorphism $`\rho `$ can be reconstructed from these charge transfers; as they are local operators in $`𝔄^d`$, they also provide an extension of $`\rho `$ to $`𝔄^d`$ with the same localization (only in the case $`d2`$). The superselection structure for $`𝔄`$ and $`𝔄^d`$ coincide (the extension map is a full functor) and, replacing $`𝔄`$ by $`𝔄^d`$, we may thus assume that $`𝔄`$ satisfies Haag duality (in the original representation of $`𝔄`$).
As shown in , attached with any localized endomorphism $`\rho `$ there is a unitary representation of the permutation group $`_{\mathrm{}}`$, the statistics of $`\rho `$, that is classified by a statistics parameter $`\lambda _\rho `$ whose possible values are $`\lambda _\rho =0,\pm 1,\pm \frac{1}{2},\pm \frac{1}{3}\mathrm{}`$. Thus the statistical dimension $`d_{DHR}(\rho )=|\lambda _\rho |^1`$ takes integral values
$$d_{DHR}(\rho )=1,2,3,\mathrm{}+\mathrm{}.$$
By the index-statistics theorem , $`d_{DHR}(\rho )`$ coincides with an analytic dimension, the square root of the Jones index
$$d_{DHR}(\rho )=[𝔄:\rho (𝔄)]^{\frac{1}{2}},$$
(one way to read $`[𝔄:\rho (𝔄)]`$ is $`[𝔄(W):\rho (𝔄(W))]`$, with $`W`$ a wedge region).
We shall denote by $``$ be the tensor category of translation covariant localized endomorphisms of $`𝔄`$ with finite statistics (i.e. with finite dimension). For an object $`\rho `$ of $``$, the intrinsic dimension coincides with the statistical dimension :
$$d(\rho )=d_{DHR}(\rho ).$$
###### Theorem 2.2.
Let $`𝔄`$ be as above, $`\alpha _t=\tau _{x(t)}`$ a one-parameter automorphism group of translations of $`𝔄`$ and $`\phi `$ a translation invariant state which is extremal KMS for $`\alpha `$.
If $`u(\rho ,t)`$ is a $`\alpha `$-covariance cocycle for the irreducible localized endomorphism $`\rho `$ (i.e. eq. (14) holds), then $`u_\rho `$ is holomorphic $`d(\rho )d_{geo}(\rho )`$.
If moreover $`\pi _\phi `$ satisfies Haag duality, then
$$d(\rho )=d_{geo}(\rho ).$$
(29)
In particular the right hand side in the above formula is independent of the KMS state $`\phi `$.
As is known the cyclic vector $`\xi _\phi `$ in the GNS representation of a KMS state $`\phi `$ is separating for $`\pi _\phi (𝔄)^{\prime \prime }`$, namely $`\phi `$ is a separating state. This is crucial for the following theorem of Takesaki and Winnink.
###### Theorem 2.3.
(). A KMS state $`\phi `$ is locally normal, namely $`\phi |_{𝔄(𝒪)}`$ is normal for any double cone $`𝒪𝒦`$.
Note that if $`𝔄`$ is a Poincaré covariant net and $`\rho `$ a Poincaré covariant localized endomorphism, then the covariance cocycle $`u(\rho ,L)(\rho ,L𝒫_+^{})`$ is uniquely fixed because $`𝒫_+^{}`$ has no non-trivial finite-dimensional unitary representation, hence it is a two-variable cocycle. In particular, if $`\rho `$ extends to an irreducible endomorphism of the weak closure, then $`d(\rho )=d_{geo}(\rho )`$ by Prop. 1.3. We shall return on this point in the next section.
If moreover there is a PCT symmetry $`j`$ for $`𝔄`$, or an anti-automorphism $`j`$ of $`𝔄`$ such that $`j^1\rho j=\overline{\rho }`$, and $`\phi j=\phi `$, then $`d(\rho )=d_\phi (u_\rho )`$ by Prop. 1.4.
We have essentially mentioned that the tensor category $``$ has a permutation symmetry (if dim($`𝕄)3)`$, as shown in . Then, by there exists a field net $`𝔉`$ of von Neumann algebras $`𝔉(𝒪)𝔄(𝒪)`$, with normal commutation relations, with $`𝔄=𝔉^G`$ the fixed point of $`𝔉`$ under the action $`\gamma `$ of a compact group $`G`$ of internal symmetries of $`𝔉`$. One has $`𝔄^{}𝔉=`$, where $`𝔉`$ is the quasi-local field C-algebra $`𝔉(𝕄)`$. Every endomorphism $`\rho `$ is implemented by a Hilbert space of isometries in $`𝔉`$.
We now relax the pointwise continuity condition in . We denote by $`\stackrel{~}{\tau }`$ the translation automorphism group on $`𝔉`$ extending $`\tau `$ and by $`\stackrel{~}{\alpha }=\stackrel{~}{\tau }_{x()}`$ the one-parameter automorphism group extending $`\alpha `$.
###### Lemma 2.4.
Let $`\phi `$ be an extremal KMS state of $`𝔄`$ with respect to $`\alpha `$. There exists a locally normal state $`\psi `$ of $`𝔉`$ that extends $`\phi `$ and is extremal KMS with respect to $`\stackrel{~}{\alpha }_t\gamma _{g(t)}`$, with $`tg(t)`$ a one-parameter subgroup of $`G`$.
###### Proof.
Let $`𝔉_c𝔉`$ denote the sub-C-algebra of all elements with pointwise norm continuous orbit under the action of $`\stackrel{~}{\tau }\gamma `$ of $`^{d+1}\times G`$, and set $`𝔉_c(𝒪)𝔉_c𝔉(𝒪)`$. We have
$$𝔉_c(𝒪)^{}𝔉(𝒪_0),𝒪_0𝒦,\overline{𝒪}_0𝒪,$$
where $`𝔉_c(𝒪)^{}`$ is the $`\sigma `$-weak closure of $`𝔉(𝒪)`$. Indeed if $`X𝔉(𝒪_0)`$ and $`j_n`$ is an approximation of the identity in $`^{d+1}`$ by continuous functions with support in a ball of radius $`\frac{1}{n}`$, then
$$X_nj_n(X)\stackrel{~}{\tau }_x(X)dxX$$
($`\sigma `$-weak convergence) and $`X_n`$ has pointwise $`\stackrel{~}{\tau }`$-orbit and, for large $`n`$, belongs to $`(𝒪)`$ (the $`\gamma `$-continuity is checked similarly).
Set $`\phi _c\phi |_{𝔄_c}`$, where $`𝔄_c𝔉_c𝔄`$ and $`\stackrel{~}{\phi }=\phi \epsilon `$, where $`\epsilon =_G\gamma _gdg`$ is the expectation of $`𝔉`$ onto $`𝔄`$ as above. Clearly $`\stackrel{~}{\phi }`$ is a $`\stackrel{~}{\tau }`$-invariant locally normal state of $`𝔉`$ and so is its restriction $`\stackrel{~}{\phi }_c`$ to $`𝔉_c`$. Let $`\psi _c\stackrel{~}{\phi }_c`$ be an extremal $`\stackrel{~}{\alpha }`$-invariant state of $`𝔉_c`$ extending $`\phi _c`$. By the AHKT theorem $`\psi _c`$ is KMS with respect to a one-parameter automorphism group $`t\stackrel{~}{\alpha }_t\gamma _{g(t)}`$ of $`𝔉_c`$. Since $`\stackrel{~}{\phi }_c`$ is locally normal and dominates $`\psi _c`$, also $`\psi _c`$ is locally normal, thus it extends to a locally normal state $`\psi `$ of $`𝔉`$. By usual arguments, $`\psi `$ is a KMS state on $`𝔉`$ with respect to $`\stackrel{~}{\alpha }`$. $`\mathrm{}`$
###### Lemma 2.5.
*( Prop. III.3.2)* Let $`𝔄𝔉`$ be C-algebras, $`\phi `$ a state of $`𝔄`$ and $`\stackrel{~}{\phi }`$ an extension to $`𝔉`$. If $`\stackrel{~}{\phi }`$ is separating, then $`\pi _{\stackrel{~}{\phi }}|_𝔄`$ is quasi-equivalent to $`\pi _\phi `$.
###### Corollary 2.6.
Let $`𝔉`$ be a C-algebra and $`\rho `$ an inner endomorphism of $`𝔉`$. If $`\phi `$ is a separating state of $`𝔉`$, the GNS representations $`\pi _\phi `$ and $`\pi _{\phi \rho }`$ are quasi-equivalent.
###### Proof.
Let $`H𝔉`$ be a Hilbert space of isometries implementing $`\rho `$ on $`𝔉`$ and let $`\{v_i,i=1,\mathrm{},n\}`$ be an orthonormal basis of $`H`$, thus the $`v_i`$’s are isometries in $`𝔉`$ with orthogonal final projections summing up to the identity and $`\rho (X)=_iv_iXv_i^{}`$, $`X𝔉`$. Then
$$\begin{array}{c}(\pi _{\phi \rho }(X)\xi _{\phi \rho },\xi _{\phi \rho })=\phi (\rho (X))=(\pi _\phi (\rho (X))\xi _\phi ,\xi _\phi )\hfill \\ \hfill =(\pi _\phi (X)\xi _i,\xi _i)=(\pi _\phi (X)\mathrm{}\pi _\phi (X)\overline{\xi },\overline{\xi }),X𝔉,\end{array}$$
where $`\xi _i=\pi _\phi (v_i^{})\xi _\phi `$ and $`\overline{\xi }=\xi _i`$. Therefore $`\pi _{\phi \rho }\pi _\phi \mathrm{}\pi _\phi `$. On the other hand $`\overline{\xi }`$ is a separating vector for $`\pi _\phi \mathrm{}\pi _\phi (𝔉)^{\prime \prime }`$; indeed elements of $`\pi _\phi \mathrm{}\pi _\phi (𝔉)^{\prime \prime }`$ have the form $`Y=X\mathrm{}X`$, $`X\pi _\phi (𝔉)^{\prime \prime }`$, thus $`Y\overline{\xi }=0`$ iff $`X\pi _\phi (v_i^{})\xi _\phi =X\xi _i=0,i`$, thus iff $`X\pi _\phi (v_i^{})=0`$ because $`\xi _\phi `$ is separating, which is equivalent to $`X=0`$ because $`v_i^{}v_i=0`$. $`\mathrm{}`$
###### Lemma 2.7.
Let $`\phi `$ be an extremal KMS state of $`𝔄`$. The GNS representations $`\pi _\phi `$ and $`\pi _{\phi \rho }`$ are quasi-equivalent.
###### Proof.
Let $`\psi `$ a KMS state of $`𝔉`$ extending $`\phi `$. As $`\psi `$ is a separating state of $`𝔉`$, also $`\psi \rho _H`$ is also separating as in the proof of Prop. 2.6, where $`\rho _H`$ is the inner endomorphism of $`𝔉`$ implemented by $`H`$. As $`\psi \rho _H`$ extends $`\phi \rho `$ we have by Lemma 2.6 that $`\pi _\phi \pi _\psi |_𝔄\pi _{\psi \rho _H}|_𝔄\pi _{\phi \rho }`$, where the symbol “$``$” denotes quasi-equivalence. $`\mathrm{}`$
###### Corollary 2.8.
Every localized endomorphism $`\rho `$ is normal with respect to $`\phi `$.
###### Proof.
The proof now follows by Lemma 2.7 and Lemma 1.5. $`\mathrm{}`$
Proof of Theorem 2.2 As $`\phi `$ is extremal KMS, $`\pi _\phi (𝔄)^{\prime \prime }`$ is a factor. If $`\rho `$ is a localized endomorphism, by Corollary 2.8 both $`\rho `$ and its conjugate extend to $``$. The conjugate equation then holds on $``$ showing that the extension $`\widehat{\rho }`$ of $`\rho `$ to $``$ has finite dimension. We have $`d_{an}(\widehat{\rho })=d(\rho )`$ (a priori we only have $`d_{an}(\widehat{\rho })d(\rho )`$). Indeed Proposition 1.6 applies.
If Haag duality holds in the representation $`\pi _\phi `$, then by the following Lemma 2.9 $`\widehat{\rho }`$ is irreducible if $`\rho `$ is irreducible, therefore the last part of the statement follows by Prop. 1.9. The rest now follows from the analysis in the previous section. $`\mathrm{}`$
###### Lemma 2.9.
With the above notation, if $`\pi _\phi `$ satisfies Haag duality, then the extension map $`\rho \widehat{\rho }\mathrm{End}()`$ is a full functor.
###### Proof.
With $`\rho `$ irreducible, we have to show that $`\widehat{\rho }`$ is irreducible too. Let $`T(\widehat{\rho },\widehat{\rho })`$, namely $`T`$ and $`T\widehat{\rho }(X)=\widehat{\rho }(X)T`$ for all $`X`$ in $``$. As $`\rho `$ is localized in a double cone $`𝒪`$, $`\rho `$ acts identically on $`𝔄(𝒪^{})`$, hence $`T𝔄(𝒪^{})^{}=𝔄(𝒪)`$, thus $`T(\rho ,\rho )=`$ as desired. $`\mathrm{}`$
### 2.1 The absolute and the relative part of the incremental free energy.
Beside the description of the chemical potential in terms of extensions of KMS states, the AHKT work provides an intrinsic description of the chemical potential within the observable algebra , see also , that was made explicit only in the case of abelian charges (automorphisms).
Let’s recall this point. Let $`𝔄`$ be a unital C-algebra with trivial centre and $`\alpha `$ a one-parameter automorphism group as before and $`\rho `$ a covariant automorphism of $`𝔄`$, thus $`\rho \alpha _t\rho ^1=\text{Ad}u(t)\alpha _t`$ for some $`\alpha `$-cocycle of unitaries $`u(t)𝔄`$. Notice that $`u`$ is unique up to multiplication by a one dimensional character of $``$, that one fixes once for all.
If now $`\phi `$ is an extremal KMS state for $`\alpha `$ such that $`\phi \rho ^1`$ and $`\phi `$ are quasi-equivalent, thus $`\rho `$ extends to the factor $`\pi _\phi (𝔄)^{\prime \prime }`$, then
$$u(t)=e^{i\mu _\rho (\phi )t}(D\phi \rho ^1:D\phi )_{\beta ^1t}$$
(30)
for some $`\mu _\rho (\phi )`$, called the chemical potential of $`\phi `$ (we are assuming that $`\phi (u())`$ is continuous). A relevant observation is that, although $`\mu _\rho (\phi )`$ depends on the initial phase fixing for $`u`$,
$$\mu _\rho (\phi ^{}|\phi )\mu _\rho (\phi ^{})\mu _\rho (\phi )$$
is independent of that and is therefore an intrinsic quantity associated with a pair $`\phi `$, $`\phi ^{}`$ of extremal KMS states. In other words the chemical potential is a label for the different extremal KMS states.
Now the above argument goes true in more generality if $`\rho `$ is an endomorphism of $`𝔄`$ that extends to a finite-index irreducible endomorphism $``$, once we replace $`\phi \rho ^1`$ with $`\phi _\rho \phi \mathrm{\Phi }_\rho `$, where $`\mathrm{\Phi }_\rho `$ is the minimal left inverse of the extension of $`\rho `$. In general, for a given charge $`\rho `$, the chemical potential is only defined with respect to the two thermal states
$$\mu _\rho (\phi ^{}|\phi )\beta ^1\mathrm{log}d_\phi (u)\beta ^1\mathrm{log}d_\phi ^{}(u)=\beta ^1\mathrm{log}d_\phi ^{}((D\phi _\rho :D\phi )).$$
Moreover, if there is a canonical way to choose the cocycle $`u`$, independently of the state $`\phi `$, then
$$\mu _\rho (\phi )\beta ^1\mathrm{log}d_\phi (u)\beta ^1\mathrm{log}d(u).$$
(31)
defines an absolute chemical potential in the state $`\phi `$, associated with the charge $`\rho `$.
The above discussion relies of course on the normality of the endomorphism $`\rho `$ in a thermal state, a deep fact, proved in when the endomorphisms $`\rho `$ are associated to the dual of a compact gauge group, with certain asympotically abelian and cluster properties for the dynamics.
We now apply the above discussion to the case of quantum relativistic statistical mechanics, namely we consider thermal states for the time evolution in a quantum field theory on Minkowski spacetime. In this situation, there is a two variable cocycle for the Poincaré covariant endomorphisms with finite dimension, which is unique because the Poincaré has no non trivial finite dimensional unitary representation. Hence $`\mu _\rho (\phi )`$ can be defined intrinsically by eq. (31).
To be more explicit, let $`𝔄`$ be a net of von Neumann algebras on the Minkowski space, as in the previous section. We further assume that $`𝔄`$ is Poincaré covariant, namely there is a unitary representation $`U`$ of $`𝒫_+^{}`$ on $``$ that acts covariantly on $`𝔄`$ and extends the translation unitary group.
Our endomorphisms $`\rho `$ are assumed to be covariant with respect to the action of $`𝒫_+^{}`$, namely there exists an $`\alpha `$ cocycle of unitaries $`u(\rho ,L)𝔄`$ such that
$$\text{Ad}u(\rho ,L)\rho =\alpha _L\rho \alpha _L^1,L𝒫_+^{}.$$
(32)
where $`\alpha _L=\mathrm{Ad}U(L)`$. This is indeed a two-variable cocycle.
Restricting this cocycle to the subgroup of time translation, we obtain a canonical choice for the unitary cocycle $`u(\rho ,t),t`$, for the one parameter group $`\alpha `$, which is indeed a two-variable cocycle in $`\times `$.
###### Theorem 2.10.
Let $`𝔄`$ be the quasi-local observable C-algebra and $`\alpha `$ a one-parameter (time) translation automorphism group as in the previous section. If $`\rho `$ is a Poincaré covariant irreducible localized endomorphism with finite dimension, we have:
* If $`\phi `$ is an extremal KMS state for $`\alpha `$ satisfying Haag duality, there exists a chemical potential $`\mu _\rho (\phi )`$ associated with $`\rho `$ defined by the canonical splitting
$$\mathrm{log}d_\phi (u_\rho )=\mathrm{log}d(\rho )+\beta \mu _\rho (\phi )$$
and satisfies
$$\mu _\rho (\phi )=\mu _{\overline{\rho }}(\phi ).$$
The intrinsic dimension is thus given by
$$\mathrm{log}d(\rho )=\frac{1}{2}(\mathrm{log}d_\phi (u_\rho )+\mathrm{log}d_\phi (u_{\overline{\rho }}))$$
independently of $`\phi `$.
* If there is a time reversal symmetry as in Prop. 1.4 and $`\phi j=\phi `$, then $`\mu _\rho (\phi )=0`$, namely
$$d_\phi (u_\rho )=\mathrm{log}d(\rho ),$$
independently of $`\phi `$.
###### Proof.
As above noticed, $`u(\rho ,L)`$ is a two-variable cocycle for the action of $`_0\times 𝒫_+^{}`$ on $`𝔄`$, where $`_0`$ is the tensor category with conjugates generated by $`\rho `$ (see ). Hence by
$$u(\rho ,L)^{}=u(\overline{\rho },L).$$
On the other hand, by the results in the previous section, we may extend $`\rho `$ to the weak closure of $`𝔄`$ in the GNS representation of $`\phi `$, and if then compare with the Connes cocycle, we have
$$u(\rho ,t)=e^{i\mu _\rho (\phi )t}d(\rho )^{i\beta ^1t}(D\phi _\rho :D\phi )_{\beta ^1t}.$$
But $`(d(\rho )^{it}(D\phi _\rho :D\phi )_t)^{}=(d(\rho )^{it}(D\phi _{\overline{\rho }}:D\phi )_t)`$, hence
$$u(\overline{\rho },t)=e^{i\mu _\rho (\phi )t}d(\rho )^{i\beta ^1t}(D\phi _{\overline{\rho }}:D\phi )_{\beta ^1t}$$
namely $`\mu _{\overline{\rho }}(\phi )=\mu _\rho (\phi )`$.
The last point is a consequence of Proposition 1.4. $`\mathrm{}`$
Note in particular that by the point $`(ii)`$ in the above theorem the chemical potential vanishes if there is a time-reversal symmetry, therefore a non-trivial chemical potential sets an arrow of time, in accordance with the second principle of thermodynamics.
Now we make contact with the analysis in . Since the covariance cocycle $`u=u_\rho `$ is canonically defined in the above context, once we choose the thermal state $`\phi `$, the Hamiltonian in the state $`\phi _\rho `$ is canonically defined as
$$H_\rho i\frac{\mathrm{d}}{\mathrm{d}t}u(\rho ,t)e^{itH}|_{t=0},$$
where $`H=\beta ^1\mathrm{log}\mathrm{\Delta }_\xi `$ is the Hamiltonian in state $`\phi `$. Here $`\xi `$ is the GNS vector associated with $`\phi `$. The increment of the free energy between the states $`\phi `$ and $`\phi _\rho `$ is then defined (cf. ) as
$$F(\phi |\phi _\rho )\phi _\rho (H_\rho )\beta ^1S(\phi |\phi _\rho )=\beta ^1\mathrm{log}(e^{\beta H_\rho }\xi ,\xi ),$$
where $`S(\phi |\phi _\rho )`$ is Araki’s relative entropy. It is immediate from the last expression that
$$F(\phi |\phi _\rho )=\beta ^1\mathrm{log}d_\phi (u_\rho )=\beta ^1\mathrm{log}d(\rho )\mu _\rho (\phi ).$$
If $`\phi ^{}`$ is another extremal KMS state as above, so that $`\rho `$ is normal with respect to $`\phi ^{}`$, we may now define the increment of the free energy $`F(\phi ^{}|\phi _\rho ^{})=\beta ^1\mathrm{log}d_\phi ^{}(u_\rho )`$, where $`\phi _\rho ^{}=\phi ^{}\mathrm{\Phi }_\rho `$, so we have:
$$F(\phi ^{}|\phi _\rho ^{})F(\phi |\phi _\rho )=\mu _\rho (\phi )\mu _\rho (\phi ^{})=\mu (\phi ^{}|\phi ),$$
moreover
$$S_c(\rho )=\mathrm{log}d(\rho )^2=\beta (F(\phi |\phi _\rho )+F(\phi |\phi _{\overline{\rho }}))$$
is an integer independent of $`\phi `$. Here $`S_c(\rho )`$ is the conditional entropy of $`\rho `$, see .
According to the thermodynamical formula “$`\mathrm{d}F=\mathrm{d}ET\mathrm{d}S`$”, we have obtained the following relation:
$$F(\phi |\phi _\rho )=\mu _\rho (\phi )\frac{1}{2}\beta ^1S_c(\rho ).$$
The quantity $`\mu _\rho (\phi )`$ may be interpreted as part of the energy increment obtained by adding the the charge $`\rho `$ to the identical charge, more specifically the part which is asymmetric with respect to charge conjugation. The total increment of the free energy contains also a part which is symmetric under charge conjugation and independent of the thermal equilibrium state, namely the intrinsic increment of the entropy $`\frac{1}{2}S_c(\rho )`$ multiplied by the temperature $`\beta ^1`$.
The above analysis simply goes through when we consider the increment of the free energy between two thermal states $`\phi _\rho `$ and $`\phi _\sigma `$ (cf. ). We shall make this explicitly in the context of Section 5.
## 3 Chemical potential. Low dimensional case.
We now study the chemical potential structure in the low dimensional case. The higher dimensional methods cannot be applied to this context, but we shall see that an analysis is possible by using Wiesbrock’s characterizations of conformal nets on $`S^1`$ .
### 3.1 KMS states and the generation of conformal nets.
Let $``$ denote the set of all bounded open non-empty intervals of $``$. We shall consider a net $`𝔄`$ of von Neumann algebras on $``$, namely an inclusion preserving map
$$I𝔄(I)$$
from $``$ to von Neumann algebras $`𝔄(I)`$, not necessarily acting on the same Hilbert space. We denote by the same symbol the quasi-local observable C-algebra $`𝔄=_I𝔄(I)^{}`$ (norm closure).
We shall assume the following properties of $`𝔄`$:
a) Translation covariance: There exists a one-parameter automorphism group $`\tau `$ of $`𝔄`$ that corresponds to the translations on $``$,
$$\tau _s(𝔄(I))=𝔄(I+s),I,s.$$
b) Properly infiniteness: For each $`I`$, the von Neumann algebra $`𝔄(I)`$ is properly infinite.<sup>7</sup><sup>7</sup>7This assumption is needed only for the the local normality of the KMS states (above Th. 2.3 from ) and can alternatively be replaced by the factoriality of $`𝔄(I)`$.
Let now $`\phi `$ be a KMS state on $`𝔄`$ with respect to $`\tau `$; for simplicity we set $`\beta =1`$. Let $`(_\phi ,\pi _\phi ,\xi _\phi )`$ be the associated GNS triple and $`V`$ the one-parameter unitary group implementing $`\tau `$:
$$V(s)\pi _\phi (a)\xi _\phi =\pi _\phi (\tau _s(a))\xi _\phi .$$
Note that by the KMS condition $`s\phi (a\tau _s(b))`$ is a continuous map for all $`a,b𝔄`$, hence $`V`$ is strongly continuous.
Recall now that $`𝔄`$ is additive (resp. strongly additive) if
$$𝔄(I)𝔄(I_1)𝔄(I_2),$$
whenever $`I,I_1,I_2`$ and $`II_1I_2`$ (resp. $`I\overline{I}_1\overline{I}_2`$), where the bar denotes the closure.
We now set $`𝒜(I)=\pi _\phi (𝔄(I))`$, $`I`$, which is a von Neumann algebra by Th. 2.3, and $`𝒜(E)\{𝒜(I):I,IE\}`$ for any set $`E`$ (the von Neumann algebras generated). Again we now assume $`\pi _\phi `$ to be one-to-one and identify $`𝔄(I)`$ with $`𝒜(I)`$, namely we consider the net already in its GNS representation.
The following KMS version of the Reeh-Schlieder theorem is known to experts.
###### Proposition 3.1.
$`\xi `$ is cyclic and separating for $`𝒜(I)`$, if $`I`$ is a half-line. If $`𝔄`$ is additive, then $`\xi `$ is cyclic and separating for all $`𝒜(I)`$, $`I`$.
###### Proof.
Assume first that $`I`$ is a half-line and let $`\eta `$ be orthogonal to $`𝒜(I)\xi `$; we have to show that $`\eta =0`$. Indeed if $`I_0`$ is a half-line and $`\overline{I}_0I`$, then for all $`a𝒜(I_0)`$
$$(\eta ,V(s)a\xi )=0,$$
for all $`s`$ such that $`I_0+sI`$. But, because of the KMS property, the function $`s(\eta ,V(s)a\xi )`$ is the boundary value of a function analytic in the strip $`0<\text{Im}z<\frac{1}{2}`$ (as $`V(\frac{i}{2})=\mathrm{\Delta }^{\frac{1}{2}}`$ and Dom$`(\mathrm{\Delta }^{\frac{1}{2}}𝒜(I_0)\xi `$), hence it must vanish everywhere. It follows that $`\eta `$ is orthogonal to $`𝒜(I_0+s)\xi `$ for all $`s`$, hence $`\eta `$ is orthogonal to $`\overline{_s𝒜(I_0+s)\xi }\overline{𝔄\xi }=`$.
Assume now that $`I`$ and $`𝒜`$ is additive. Set $`𝒜_0(I)=\{𝒜(I_0):I_0,\overline{I}_0I\}`$. We shall show that $`\xi `$ is cyclic for $`𝒜_0(I)`$, hence for $`𝒜(I)`$.
By the same argument as above
$$\eta 𝒜_0(I)\xi \eta 𝒜(I_0+s)\xi ,s,\overline{I}_0I,$$
namely the orthogonal projection $`P`$ onto $`\overline{𝒜_0(I+s)\xi }`$ is independent of $`s`$ and thus belongs to $`_s𝒜_0(I+s)^{}=(_s𝒜_0(I+s))^{}=𝔄^{}`$ (by additivity). As $`\xi `$ is separating for $`𝔄^{}`$ and $`P\xi =\xi `$ it follows that $`P=1`$, namely $`\overline{𝒜_0(I)\xi }=`$. $`\mathrm{}`$
Now $`\tau `$ extends to the rescaled modular group of $`𝒜()`$ with respect to $`\phi `$ and $`\tau _s(𝒩)=𝒜(s,\mathrm{})𝒩,s>0`$, where $`𝒩(0,\mathrm{})`$, namely $`(𝒩,\xi )`$ is a half-sided modular inclusion of von Neumann algebras and by Wiesbrock’s theorem there exists a $`\xi `$-fixing one-parameter unitary group $`U`$ on $``$ with positive generator such that
$$V(s)U(t)V(s)=U(e^st)$$
(33)
$$U(1)U(1)=𝒩$$
(34)
Setting $`(a,b)𝒜(\mathrm{log}a,\mathrm{log}b),b>a>0`$, we have a net $``$ on the intervals of $`(0,\mathrm{})`$ whose closure is contained in $`(0,\mathrm{})`$. We have the following, compare with .
###### Proposition 3.2.
Let $`𝔄`$ be an additive net as above and $`\phi `$ a KMS state. There exists a net $``$ on the intervals of $`(0,\mathrm{})`$ such that $`(a,b)=\pi _\phi (𝔄(\mathrm{log}a,\mathrm{log}b))`$ if $`b>a>0`$. $``$ is dilation covariant and $`V`$ is the dilation one parameter group. $``$ is also translation covariant with positive energy on half-lines, namely there is a one-parameter $`\xi `$-fixing unitary group $`U`$ with positive generator such that $`\text{Ad}U(t)(a,\mathrm{})=(a+t,\mathrm{})`$, where $`(a,\mathrm{})=_{b>a}(a,b)`$.
###### Proof.
Clearly $`V(s)(a,b)V(s)=(e^sa,e^sb)`$ for positive $`a,b`$. Setting $`\overline{}(a,\mathrm{})\text{Ad}U(a)`$ we have
$$\text{Ad}U(t)\overline{}(a,\mathrm{})=\overline{}(t+a,\mathrm{}),$$
therefore, by using the relation $`V(s)U(t)V(s)=U(e^st)`$, it follows that
$$\text{Ad}V(s)\overline{}(a,\mathrm{})=\overline{}(e^sa,\mathrm{}).$$
On the other hand $`\overline{}(1,\mathrm{})=\text{Ad}U(1)=𝒩=(1,\mathrm{})`$ hence
$$(e^s,\mathrm{})=\text{Ad}V(s)(1,\mathrm{})=\text{Ad}V(s)U(1)=\text{Ad}U(e^s)=\overline{}(e^s,\mathrm{}),$$
showing the last part of the statement. $`\mathrm{}`$
We shall call the net $``$ the thermal completion of $`𝔄`$ with respect to $`\phi `$. Note that the translation unitary group $`V`$ for $`𝒜`$ becomes the dilation unitary group for $``$.
Proposition 3.2 does not give the translation covariance of the net $``$ on the bounded intervals ($`(a,b)`$ is not even defined if $`a<0`$).
Further insight in the structure of the thermal completion net may be obtained by considering a local net $`𝔄`$, namely assuming the locality condition
$$[𝔄(I_1),𝔄(I_2)]=\{0\}\text{if}I_1I_2=\mathrm{}.$$
To construct a translation covariant net we define the following von Neumann algebras:
$$\stackrel{~}{}(0,1)=\underset{s0}{}\text{Ad}V_1(s)(0,1),$$
(35)
$$\stackrel{~}{}(0,a)=\text{Ad}V(a)\stackrel{~}{}(0,1),a,$$
(36)
$$\stackrel{~}{}(a,b)=\text{Ad}U(a)\stackrel{~}{}(0,ba),a,b.$$
(37)
Here we have set $`V_1(s)U(1)V(s)U(1)`$, the one-parameter unitary group associated with the dilations with respect to the point $`1`$.
From now on the net $`𝔄`$ will be assumed to be local.
###### Theorem 3.3.
Let the net $`𝔄`$ on the intervals of $``$ be translation covariant, local, and additive and $`\phi `$ a KMS state. With the above notations, $`\stackrel{~}{}`$ defines a conformal net, indeed $`\stackrel{~}{}`$ has a conformal extension to $`S^1`$.
As a consequence the dual net of $`\stackrel{~}{}`$ on $``$ is strongly additive and conformal.
We shall see that
$$\stackrel{~}{}(a,\mathrm{})=(a,\mathrm{}),a,$$
hence $`\stackrel{~}{}`$ is an extension of $``$ on the positive half-lines and is conformal. We shall call $`\stackrel{~}{}`$ the conformal thermal completion of $`𝔄`$ with respect to $`\phi `$.
###### Proof.
We first show that the triple $`\{(0,\mathrm{})^{},\stackrel{~}{}(0,1),(1,\mathrm{}),\xi \}`$ is a +hsm factorization with respect to $`\xi `$ in the sense of , namely these three algebras mutually commute and $`(\stackrel{~}{}(0,1)(0,\mathrm{}),\xi )`$, $`((1,\mathrm{})\stackrel{~}{}(0,1)^{},\xi )`$ and $`((0,\mathrm{})^{}(1,\mathrm{})^{},\xi )`$ are +half-sided modular inclusions.
Now $`(0,\mathrm{})^{}`$ and $`(1,\mathrm{})`$ commute by the isotony of $`𝒜`$; for the same reason $`(1,\mathrm{})`$ commute with $`\stackrel{~}{}(0,1)`$, indeed $`(1,\mathrm{})`$ is $`\text{Ad}V_1`$-invariant where, as above, $`V_1(s)=U(1)V(s)U(1)`$. Again $`(0,\mathrm{})(0,1)`$, hence $`(0,\mathrm{})\stackrel{~}{}(0,1)`$ because $`V_1(s)(0,\mathrm{})V_1(s)(0,\mathrm{})`$ if $`s0`$ by translation-dilation covariance of $``$ on positive half-lines (Prop. 3.2).
Concerning the hsm properties, the only non-trivial verification is that $`(\stackrel{~}{}(0,1)(0,\mathrm{}),\xi )`$ is a +hsm inclusion, namely that
$$\text{Ad}V(s)\stackrel{~}{}(0,1)\stackrel{~}{}(0,1),s<0.$$
We thus need to show that for any fixed $`t<0`$ we have
$$\text{Ad}V(s)V_1(t)(0,1)\stackrel{~}{}(0,1),s<0.$$
Indeed if $`s<0`$ and $`t<0`$, there exist $`s^{}<0`$ and $`t^{}<0`$ such that $`V(s)V_1(t)=V_1(t^{})V(s^{})`$, as follows immediately by the corresponding relation in the “$`ax+b`$” group. Therefore
$$\text{Ad}V(s)V_1(t)(0,1)=\text{Ad}V_1(t^{})V(s^{})(0,1)\text{Ad}V_1(t^{})(0,1)\stackrel{~}{}(0,1)$$
as desired.
By a result in there exists a conformal net $`\stackrel{~}{}`$ on $``$ such that the local von Neumann algebras associated to $`(\mathrm{},0)`$, $`(0,1)`$ and $`(1,\mathrm{})`$ are respectively $`(0,\mathrm{})^{}`$, $`\stackrel{~}{}(0,1)`$ and $`(1,\mathrm{})`$ and having $`U`$ and $`V`$ as translation and dilation unitary groups.
By translation-dilation covariance, $`\stackrel{~}{}`$ is then conformal thermal completion of $`𝒜`$. $`\mathrm{}`$
We may also directly define the dual net $`^d`$ of $`\stackrel{~}{}`$ as the one associated with the half-sided modular factorization $`((0,\mathrm{})^{},(1,\mathrm{})^{}(0,\mathrm{}),(1,\mathrm{}),\xi )`$. This net is conformal, strongly additive and
$$^d(a,b)=(a,\mathrm{})(b,\mathrm{})^{}.$$
This is due to the equivalence between strong additivity and Haag duality on the real line for a conformal net, see . Clearly we have
$$(a,b)\stackrel{~}{}(a,b)^d(a,b).$$
thus $`^d`$ is the dual net of $`\stackrel{~}{}`$ and
$$^d=\stackrel{~}{}\stackrel{~}{}\text{ is strongly additive,}$$
###### Corollary 3.4.
If $`𝒜`$ is strongly additive
$$(1,\mathrm{})^{}(0,\mathrm{})=\underset{s<0}{}V(s)(0,1)V(s).$$
###### Proof.
If $`𝒜`$ is strongly additive, then $``$ is strongly additive (on the intervals of $`(0,\mathrm{})`$), hence $`\stackrel{~}{}`$ is strongly additive and the above comment applies. $`\mathrm{}`$
More directly, Corollary 3.4 states that the relative commutant $`𝒜(0,\mathrm{})^{}`$ is the smallest von Neumann algebra containing $`𝒜(\mathrm{},0)`$ which is mapped into itself by $`\text{Ad}\mathrm{\Delta }^{it}`$, $`t>0`$, where $`\mathrm{\Delta }`$ is the modular operator associated with $`(𝒜(0,\mathrm{}),\xi )`$.
We shall say that the state $`\phi `$ of $`𝔄`$ satisfies essential duality if
$$𝒜(0,\mathrm{})^{}=𝒜(\mathrm{},0).$$
We have:
###### Proposition 3.5.
The following are equivalent:
* $`\phi `$ satisfies essential duality,
* For some (hence for all) $`0<a<b`$ we have $`𝒜(b,\mathrm{})^{}𝒜(a,\mathrm{})=𝒜(a,b)`$,
* $`𝒜`$ is strongly additive and $`\text{Ad}\mathrm{\Delta }^{it}𝒜(\mathrm{},0)𝒜(\mathrm{},0)`$ for all $`t>0`$, where $`\mathrm{\Delta }`$ is the modular operator of $`(𝒜(0,\mathrm{}),\xi )`$.
###### Proof.
$`(i)(iii)`$ follows by the above comments. On the other hand $`(i)(ii)`$ because they are equivalent to the relative commutant property $`(a,b)=(b,\mathrm{})^{}(a,\mathrm{})`$ for $`b>a`$ and either $`a=0`$ or $`a>0`$, which are indeed equivalent conditions in the conformal case . $`\mathrm{}`$
###### Corollary 3.6.
If $`\phi `$ satisfies essential duality, then $`\phi `$ satisfies Haag duality, namely
$$𝒜(a,b)=(𝒜(\mathrm{},a)𝒜(b,\mathrm{}))^c,a<b,$$
where $`^c`$ denotes the relative commutant in $``$.
###### Proof.
If $`\phi `$ satisfies essential duality then, since $`𝒜(b,\mathrm{})^{}𝒜(a,\mathrm{})`$ by $`(ii)`$ of the above proposition, we have
$$𝒜(a,b)=𝒜(b,\mathrm{})^{}𝒜(a,\mathrm{})=𝒜(b,\mathrm{})^c𝒜(a,\mathrm{})$$
for $`b>a>0`$. On the other hand, by essential duality, we have $`𝒜(a,\mathrm{})=𝒜(\mathrm{},a)^c`$, hence
$$𝒜(a,b)=𝒜(\mathrm{},a)^c𝒜(b,\mathrm{})^c=(𝒜(\mathrm{},a)𝒜(b,\mathrm{}))^c$$
as desired. The case of arbitrary $`a<b`$ is obtained by translation covariance. $`\mathrm{}`$
Hence essential duality in a thermal state can occur only if the original net is strongly additive. It would interesting to see if the converse holds true, namely if all KMS states on a strongly additive net satisfy essential duality.
Note also the, in contrast to the situation occurring in the vacuum representation, the equality $`𝒜(a,b)^{}=𝒜(\mathrm{},a)𝒜(b,\mathrm{})`$ cannot hold in any thermal state, unless the superselection structure is trivial (this would be equivalent to the triviality of the 2-interval inclusion for the net $``$).
### 3.2 Normality of superselection sectors in temperature states.
In this section $`𝔄`$ will denote a local net of von Neumann algebras on the intervals $``$ of $``$ satisfying the properties $`a)`$ and $`b)`$ in the previous section.
With $`\tau `$ the translation automorphism group of $`𝔄`$, we shall say that an endomorphism $`\rho `$ of the quasi-local C-algebra $`𝔄`$ is a localized in the interval $`I`$ if $`\rho `$ acts identically on $`𝔄(I^{})`$, where $`I^{}I`$ and, for any open set $`E`$, $`𝔄(E)`$ denotes as before the C-algebra generated by the $`\{𝔄(I):I,IE\}`$. We have also set $`𝒜(E)\pi _\phi (𝔄(E))^{\prime \prime }`$.
As above, $`\rho `$ is translation covariant if there exists a unitary $`\tau `$-cocycle of unitaries $`u(s)𝔄`$ such that $`\text{Ad}u(s)\tau _s\rho \tau _s=\rho `$.
Let $`\phi `$ be a KMS state of $`𝔄`$ with respect to the translation group. In the following $`\rho `$ is a translation covariant endomorphism of $`𝔄`$ localized in an interval $`I`$. By translation covariance we may assume that $`I(0,\mathrm{})`$. Our main result in this section is the following.
###### Theorem 3.7.
Let $`𝔄`$ be a translation covariant net on $``$ and $`\phi `$ a KMS state of $`𝔄`$ satisfying essential duality. If $`\rho `$ is a translation covariant localized endomorphism of $`𝔄`$ with finite dimension $`d(\rho )`$, then $`\rho `$ is normal with respect to $`\phi `$, namely $`\rho `$ extends to a normal endomorphism of $`=\pi _\phi (𝔄)^{\prime \prime }`$.
If $`\phi `$ is an extremal KMS state, i.e. $``$ is a factor, then the extension of $`\rho `$ to $``$ has the same dimension $`d(\rho )`$.
Assuming that $`𝔄`$ acts on $`_\phi `$, as above, we have:
###### Lemma 3.8.
$`\rho |_{𝔄(a,\mathrm{})}`$ extends to a normal endomorphism of $`𝒜(a,\mathrm{})`$ for any $`a0`$.
###### Proof.
By translation covariance there exists a unitary $`u`$ such that $`\rho ^{}\text{Ad}u\rho `$ is localized in an interval contained in $`(\mathrm{},a)`$, thus $`\rho =\text{Ad}u^{}\rho ^{}=\text{Ad}u^{}`$ on $`𝒜(I)`$ for all $`I(a,\mathrm{})`$, $`I`$. It follows that $`\text{Ad}u^{}`$ is a normal extension of $`\rho `$ to $`𝒜(a,\mathrm{})`$. $`\mathrm{}`$
If the endomorphism $`\rho `$ of $`𝒜`$ is localized in the interval $`I`$ then $`\rho (𝒜(I_1))𝒜(I_1)`$ for all intervals $`I_1`$ containing $`I`$ by Haag duality. We shall say that $`\rho `$ has finite dimension if the index $`[𝒜(I_1):\rho (𝒜(I_1))]`$ is finite and independent of $`I_1`$ (the index is here defined for example by the Pimsner-Popa inequality ).
###### Lemma 3.9.
If the endomorphism $`\rho `$ has finite dimension, then the corresponding endomorphism of $`𝒜(0,\mathrm{})`$ given by Lemma 3.8 has finite dimension (i.e. finite index).
###### Proof.
Setting $`_n𝒜(0,n)`$, $`n`$, we have $`\rho (_n)_n`$ for large $`n`$. Moreover $`_n`$ and $`\rho (_n)`$ converge increasingly respectively to $`𝒜(0,\mathrm{})`$ and $`\rho (𝒜(0,\mathrm{}))`$. Thus Prop. 4 of applies and gives
$$[𝒜(0,\mathrm{}):\rho (𝒜(0,\mathrm{}))]\underset{n\mathrm{}}{lim\; inf}[_n:\rho (_n)].$$
$`\mathrm{}`$
Proof of Theorem 3.7. Let $``$ be the thermal completion of $`𝔄`$, thus in particular
$$(a,\mathrm{})=𝒜(\mathrm{log}a,\mathrm{}),a>0,$$
and denote by $`U`$ and $`V`$ the translation and dilation with respect to $``$ as above. We set $`\alpha _s=\text{Ad}V(s)`$ and $`\alpha _s^{(1)}=\text{Ad}U(1)V(s)U(1)`$. By Proposition 3.2 then $`\alpha _s^{(1)}`$ acts on $``$ as a dilation with respect to the point $`1`$, namely
$$\alpha _s^{(1)}((a,\mathrm{}))=(e^sa+1e^s,\mathrm{}),a,s,$$
and $`\alpha ^{(1)}|_{(1,\mathrm{})}`$ is the (rescaled) modular group associated with $`((1,\mathrm{}),\xi )`$.
As $`\rho `$ gives rise to a finite index endomorphism of $`(1,\mathrm{})`$ (Lemma 3.8 and 3.9) there exists a unitary $`\alpha ^{(1)}`$-cocycle $`u^{(1)}(s)(1,\mathrm{})`$ such that
$$\text{Ad}u^{(1)}(s)\alpha _s^{(1)}\rho \alpha _s^{(1)}(X)=\rho (X),X(1,\mathrm{}),s.$$
(38)
Indeed we may take $`u^{(1)}`$ as the Connes Radon-Nikodym cocycle
$$u^{(1)}(s)=(D\phi \mathrm{\Phi }:D\phi )_s,$$
where $`\mathrm{\Phi }`$ is a normal faithful left inverse of $`\rho `$ on $`𝒜(0,\mathrm{})`$ and $`\phi `$ is considered as a state on $`𝒜(0,\mathrm{})`$ .
Now, if $`\epsilon (0,1)`$, $`\rho `$ acts trivially on $`(\epsilon ,1)=\pi _\phi (𝒜(\mathrm{log}\epsilon ,0))`$ and, since $`\phi `$ satisfies essential duality, the conformal thermal completion is strongly additive and
$$(\epsilon ,\mathrm{})=(\epsilon ,1)(1,\mathrm{}).$$
As $`\alpha _s^{(1)}(X)(\epsilon ,1)`$ if $`X(\epsilon ,1)`$ and $`s>0`$, it follows that equation (38) holds true for all $`X(\epsilon ,\mathrm{})`$, $`s>0`$ .
Setting then for a fixed $`s>0`$
$$\rho (X)\text{Ad}u^{(1)}(s)\alpha _s^{(1)}\rho \alpha _s^{(1)}(X),X,$$
this formula does not depend on the choice of $`s>0`$ and provides an extension of $`\rho `$ to $``$ because of formula (38).
Now $``$ is a strongly additive local conformal net on $``$ and (the extension of) $`\rho `$ is a localized endomorphism of $``$ with finite dimension, hence $`\rho `$ is Möbius covariant , therefore the index of $`\rho ((I))(I)`$ is independent of the interval $`I`$, provided $`\rho `$ is localized within $`I`$ . This clearly implies the last part of the statement. $`\mathrm{}`$
Our results then give here a version of Theorem 2.10.
###### Corollary 3.10.
Let $`𝔄`$ be a net on $``$ as above, $`\tau `$ the translation automorphism group and $`\rho `$ a translation covariant localized endomorphism with finite dimension.
Then $`\rho `$ has finite holomorphic dimension in each extremal KMS state $`\phi `$ fulfilling essential duality.
The chemical potential associated with $`\rho `$ is defined by the canonical splitting
$$\beta ^1\mathrm{log}d_\phi (u_\rho )=\beta ^1\mathrm{log}d(\rho )+\mu _\rho (\phi ),$$
and satisfies $`\mu _{\overline{\rho }}(\phi )=\mu _\rho (\phi )`$. In particular
$$d_{geo}(\rho )=d(\rho )$$
for all irreducible $`\rho `$, independently of $`\phi `$.
### 3.3 Extension of KMS states to the quantum double.
The purpose of this subsection is to provide a description of the chemical potential, in a low dimensional theory, in terms of extensions of KMS states, in analogy to what described in in the higher dimensional case. As in our cases charges are not any longer associated to the dual of a compact gauge group, we shall replace the field algebras by a quantum double construction , that we will perform in the C-case in Appendix A.2 for our purposes: the reader is referred to this appendix for the necessary notations and background.
Let then $`𝔄`$ be a unital C-algebra with trivial centre and $`𝒯\mathrm{End}(𝔄)`$ a tensor category of endomorphisms with conjugates and sub-objects. Let $`\alpha `$ be a one-parameter group of automorphisms of $`𝔄`$ and $`u(\rho ,t)`$ unitary covariance cocycle (eq. (14)) which is a two-variable cocycle. We consider $`𝒯^{op}𝔄`$ by setting $`\rho ^{op}=\overline{\rho }`$ and $`(\rho ^{op},\sigma ^{op})=(\rho ,\sigma )^{}`$. Let $`I`$ be an index set so that $`\rho _i,iI`$ is a family of inequivalent irreducible objects of $`𝒯`$, one for each equivalence class. Then $`\stackrel{~}{u}(\stackrel{~}{\rho }_i,t)u(\rho _i,t)u(\rho _i,t)^{}`$, where $`\stackrel{~}{\rho }=\rho \rho `$, extends to a two-variable cocycle for $`\alpha _t\alpha _t`$ and $`𝒯\times 𝒯^{op}`$. Indeed $`\stackrel{~}{u}(\stackrel{~}{\rho }_i,t)`$ is independent of its choice (phase fixing) due to the anti-linearity of the Frobenious map $`TT^{}`$.
We may extend $`\alpha _t\alpha _t`$ to a one-parameter automorphism group $`\stackrel{~}{\alpha }`$ of $`𝔅`$ by setting
$$\stackrel{~}{\alpha }_t(a)=\alpha _t\alpha _t(a),a\stackrel{~}{𝔄}$$
(39)
$$\stackrel{~}{\alpha }_t(R_i)=\stackrel{~}{u}(\stackrel{~}{\rho }_i,t)^{}R_i,$$
(40)
(cf. ). If $`\phi `$ is a KMS state for $`\alpha _t`$ on $`𝔄`$, then $`\phi \phi `$ is a KMS state for $`\alpha _t\alpha _t`$ on $`\stackrel{~}{𝔄}𝔄𝔄`$.
###### Proposition 3.11.
Let $`\phi `$ be an extremal KMS state for $`\alpha _t`$. Then $`\phi \phi \epsilon `$ extends to a KMS state of $`\pi _{\stackrel{~}{\phi }}(𝔅)^{\prime \prime }`$ with respect to $`\stackrel{~}{\alpha }_t\theta _t`$ for some one-parameter automorphism group of $`\pi _{\stackrel{~}{\phi }}(𝔅)^{\prime \prime }`$ leaving $`\stackrel{~}{𝔄}`$ pointwise fixed, if and only if each $`\rho 𝒯`$ is normal with respect to $`\phi `$.
###### Proof.
If $`\stackrel{~}{\phi }\phi \phi \epsilon `$ is KMS state, then $`\stackrel{~}{\phi }`$ is a separating state. If $`R𝔅`$ is a non-zero multiple of an isometry, then also $`\psi =\stackrel{~}{\phi }(RR^{})`$ is a separating positive functional, which is quasi-equivalent to $`\stackrel{~}{\phi }`$ because $`R^{}\xi _{\stackrel{~}{\phi }}`$ is a separating vector for $`\pi _{\stackrel{~}{\phi }}(𝔅)^{\prime \prime }`$.
Take $`R=R_i`$, the element of $`𝔅`$ that implements $`\stackrel{~}{\rho }_i`$. If $`x\stackrel{~}{𝔄}`$ we have
$$\psi (x)=\stackrel{~}{\phi }(R_ixR_i^{})=\stackrel{~}{\phi }(\stackrel{~}{\rho }_i(x)E_i)=\phi \phi (\epsilon (\stackrel{~}{\rho }_i(x)E_i)=\phi \phi \stackrel{~}{\rho }_i(x)$$
where $`E_i=R_iR_i^{}`$, thus $`\epsilon (E_i)=1`$. Thus $`\psi =\phi \phi \rho _i`$, hence
$$\pi _{\phi \phi \stackrel{~}{\rho }_i}\pi _\psi |_{\stackrel{~}{𝔄}}\pi _{\stackrel{~}{\phi }}|_{\stackrel{~}{𝔄}}\pi _{\phi \phi },$$
and this implies that $`\rho _i`$ is normal with respect to $`\phi `$.
Conversely, let us assume that each $`\rho _i`$ is normal with respect to $`\phi `$ and still denote by $`\stackrel{~}{\rho }_i`$ the extension of $`\stackrel{~}{\rho }_i`$ to the von Neumann algebra $`\stackrel{~}{}\pi _{\phi \phi }(\stackrel{~}{𝔄})^{\prime \prime }`$. By rescaling the parameter we may set the inverse temperature $`\beta `$ equal to $`1`$. We may then define the \*-algebra $`𝔅_0`$ as in the appendix with $`𝔄`$ replaced by its weak closure $`\stackrel{~}{}`$ (but the $`C_{ij}^k`$ are still defined with respect to $`\stackrel{~}{𝔄}`$).
Let $`\mathrm{\Psi }_{\stackrel{~}{\rho }_i}=C_{\overline{i}i}^0\overline{\stackrel{~}{\rho }_i}()C_{\overline{i}i}^0`$, thus $`\mathrm{\Psi }_{\stackrel{~}{\rho }_i}`$ is a (not necessarily standard) non normalized left inverse of $`\stackrel{~}{\rho }_i`$. With $`V(\stackrel{~}{\rho }_i,t)=(D\phi \phi \mathrm{\Psi }_{\stackrel{~}{\rho }_i}:D\phi \phi )_t`$, define a one-parameter automorphism group $`\gamma `$ of $`𝔅_0`$ extending $`\alpha \alpha `$ as was done for $`\stackrel{~}{\alpha }`$, but using $`V`$ instead of $`U`$, thus, with obvious notations, $`\gamma _t(R_i)=V(\stackrel{~}{\rho }_i,t)^{}R_i`$. By using the holomorphic properties of the Connes cocycles and the corresponding two-variable cocycle property (that can be checked similarly as in ), it can be seen by elementary calculations that $`\stackrel{~}{\phi }`$ is KMS with respect to $`\gamma `$. Clearly both $`\gamma `$ and $`\stackrel{~}{\alpha }`$ extend to $`𝒩=\pi _{\stackrel{~}{\phi }}(𝔅_0)^{\prime \prime }`$, indeed the extension of $`\gamma `$ is the modular group with respect to $`\stackrel{~}{\phi }`$, thus $`\gamma `$ and $`\stackrel{~}{\alpha }`$ commute and $`\theta _t=\stackrel{~}{\alpha }_t\gamma _t`$ is a one-parameter group of $`𝒩`$ leaving $`\stackrel{~}{}`$ pointwise fixed. $`\mathrm{}`$
Let $`G`$ be the “gauge group”, namely the group of all automorphisms of $`𝔅`$ that leave $`\stackrel{~}{𝔄}`$ pointwise invariant, and $`Z(G)`$ the centre of $`G`$. Any extension $`\widehat{\alpha }_t`$ of $`\alpha _t\alpha _t`$ to a one parameter automorphism group of $`𝔅`$ is clearly given by $`\stackrel{~}{\alpha }_t=\widehat{\alpha }_t\theta _t`$ where $`\theta _tG`$. Reasoning as in the above proof one has the following.
###### Corollary 3.12.
Let $`\phi `$ and $`\psi `$ be extremal KMS states for $`\alpha `$. Assume that all $`\rho 𝒯`$ are normal with respect to both $`\phi `$ and $`\psi `$ and that if $`\rho `$ is irreducible the normal extensions of $`\rho `$ with respect to $`\phi `$ and $`\psi `$ are still irreducible. Then $`\phi \psi \epsilon `$ is a KMS state of $`𝔅`$ with respect to $`\stackrel{~}{\alpha }\theta `$, with $`\theta `$ a one-parameter subgroup of $`Z(G)`$.
In particular this is the case if $`\phi `$ and $`\psi `$ are extremal KMS states with respect to translations of a net of local von Neumann algebras as in Section 3, satisfying essential duality.
The one-parameter group $`\theta `$ is related to the chemical potential associated with the charge $`\rho _i`$ by
$$\theta _t(R_i)=e^{i\mu _{\rho _i}(\phi |\psi )t}R_i.$$
## 4 Roberts and Connes-Takesaki cohomologies.
Roberts has obtained a geometrical description of the theory of superselection sectors by considering a non-abelian local cohomology naturally associated with a net $`𝔄`$. Rather than stating here his formal definitions, that can be obtained from different geometrical pictures, we will recall here the main idea underlying the construction of the first cohomology semiring $`H_R^1(𝔄)`$. Here $`𝔄`$ is a net of von Neumann algebras on the Minkowski spacetime $`𝕄`$, fulfilling the properties stated in Section 2, but the construction may be extented to more general globally hyperbolic spacetimes, see .
Given a representation $`\pi `$ of $`𝔄`$ localizable in all double cones we may, as usual, fix a double cone $`𝒪𝒦`$ and choose an endomorphism $`\rho _𝒪`$ of $`𝔄`$ localized in $`𝒪`$ and equivalent to $`\pi `$. For each $`𝒪_1𝒦`$ we choose a unitary intertwiner $`u_{𝒪,𝒪_1}`$ between $`\rho _𝒪`$ and $`\rho _{𝒪_1}`$ and set
$$z_{𝒪_2,𝒪_1}u_{𝒪_2,𝒪}^{}u_{𝒪,𝒪_1}𝒪_1,𝒪_2𝒦.$$
By duality, $`z_{𝒪_2,𝒪_1}`$ belongs to $`𝔄^d(𝒪_1𝒪_2)`$, in particular $`z_{𝒪_2,𝒪_1}𝔄(\stackrel{~}{𝒪})`$ if $`\stackrel{~}{𝒪}`$ is a double cone containing $`𝒪_1𝒪_2`$, and satisfies the cocycle condition
$$z_{𝒪_3,𝒪_1}=z_{𝒪_3,𝒪_2}z_{𝒪_2,𝒪_1}.$$
The choice of $`\rho _𝒪`$ and $`u_{𝒪,𝒪_1}`$ is not unique, yet different choices gives cohomologous cocycles under a natural equivalence relation.
One can then formalize the definition of $`H_R^1(𝔄)`$. The relevant point is that the map
$$\text{Superselection sectors}H_R^1(𝔄)$$
is invertible, namely each localized cocycle $`z`$ arises from a localized endomorphism $`\rho `$ localized in a given $`𝒪𝒦`$, which is well-defined by the formula
$$\rho (X)=z_{𝒪_1,𝒪}Xz_{𝒪_1,𝒪}^{},X𝔄(\stackrel{~}{𝒪}),$$
where $`\stackrel{~}{𝒪}`$ is a double cone containing $`𝒪`$ and $`𝒪_1`$ is any double cone contained in $`\stackrel{~}{𝒪}^{}`$.
As $`H_R^1(𝔄)`$ is in one-to-one correspondence with the superselection sectors, it is endowed with a ring structure, that can be expressed more directly. Analogusly $`Z_R^1(𝔄)`$ is a tensor C-category.
Let’s now recall the cohomology considered by Connes and Takesaki . This concerns an automorphism group action $`\tau `$ on a C-algebra $`𝔄`$, in our case the translation automorphism group on the quasi-local C-algebra $`𝔄`$. Restricting to unitary cocycles, a map $`z:^4𝔄`$, taking values in the unitaries $`𝔄`$, is a cocycle if satisfies the cocycle equation
$$z(x+y)=z(x)\tau _x(z(y)),x,y^4;$$
(41)
two cocycles $`z`$ and $`z^{}`$ are cohomologous if there exists a unitary $`u𝔄`$ such that
$$z^{}(x)=uz(x)\tau _x(u^{}),x^4.$$
(42)
A cocycle $`z`$ gives rise to a perturbed automorphism group $`\tau ^z`$ defined by $`\tau _x^zz(x)\tau _x()z(x)^{}`$. We are interested in the case $`\tau ^z`$ is a local perturbation, in the sense that, if $`x`$ varies in a bounded set, there exists a double cone $`𝒪𝒦`$ such that
$$\tau _x^z|_{𝔄(𝒪^{})}=\tau _x|_{𝔄(𝒪^{})}.$$
This amounts to define a (unitary) localized cocycle as a unitary map $`z:^4𝔄`$ satisfying the cocycle condition (41) and the locality condition:
$$\delta >0\&𝒪𝒦\text{such that}z(x)𝔄(𝒪),x^4,|x|<\delta .$$
Here $`|x|`$ is the Euclidean modulus of $`x`$. By iterating the cocycle equation (41) it follows immediatly from the locality condition that all the $`z(x)`$ lives in a common double cone as $`x`$ varies in a bounded set. Indeed the following holds.
###### Lemma 4.1.
Let $`z`$ be a localized cocycle. There exists $`r>0`$ such that
$$z(x)𝔄(𝒪_{r+|x|}),x^4,$$
where $`𝒪_r`$ denotes the double cone with basis the ball of radius $`r>0`$, centered at $`0`$, in $`^3`$.
###### Proof.
Let $`\delta ,r>0`$ be such that $`z(x)𝔄(𝒪_r)`$ if $`|x|<\delta `$. If $`(n1)r|x|<nr`$, write $`x=x_1+x_2+\mathrm{}+x_n`$ with $`|x_i|<r`$. By the cocycle equation
$$z(x)=z(x_1)\tau _{x_1}(z(x_2))\tau _{x_1+x_2}(z(x_3))\mathrm{}\tau _{x_1+x_2+\mathrm{}x_{n1}}(z(x_n))$$
we see that $`z(x)_{i=1}^{n1}𝔄(x_1+\mathrm{}+x_i+𝒪_r)𝔄(𝒪_{(n1)r})𝔄(𝒪_{r+|x|})`$. $`\mathrm{}`$
We denote by $`Z_\tau ^1(𝔄)`$ the set of unitary localized cocycles, and by $`H_\tau ^1(𝔄)`$ the quotient of $`Z_\tau ^1(𝔄)`$ modulo the equivalence relation (42), with the further restriction that the unitary $`u`$ is local, namely $`u`$ belongs to $`𝔄(𝒪)`$ for some $`𝒪𝒦`$.
Now a covariant localized endomorphism gives rise to a localized cocycle (formula (28), hence a covariant sectors to an element of $`H_\tau ^1(𝔄)`$. As for the Roberts cohomology, the converse is true.
###### Proposition 4.2.
There is a natural map
$$H_\tau ^1(𝔄)\text{Covariant superselection sectors}.$$
###### Proof.
We shall associate a covariant localized endomorphism $`\rho `$ to a given $`zZ_\tau ^1(𝔄)`$. Similarly as in eq. (4), we set
$$\rho (X)z(x)Xz(x)^{},X𝔄(𝒪_R),|x|>2R,$$
where $`x`$ is a vector in the time-zero hyperplane and $`R>r`$, with $`r`$ the radius of the double cone in Lemma 4.1. We have to show that, for a fixed $`R`$, the above definition is independent of the choice of $`x`$, namely $`z(x)Xz(x)^{}=z(x^{})Xz(x^{})^{}`$ if also $`|x^{}|>2R`$. As $`𝒪_{2R}^{}`$ is connected, by iterating the procedure, it is enough to check this if $`x^{}=x+y`$ with $`|y|<\delta `$. Then, by local commutativity,
$$z(x^{})Xz(x^{})^{}=z(x)\tau _x(z(y))X\tau _x(z(y))^{}z(x)^{}=z(x)Xz(x)^{}$$
because $`\tau _x(z(y))`$ belongs to $`𝔄(𝒪_r+x)`$ and $`𝒪_r+x𝒪_R^{}`$. $`\mathrm{}`$
The map given by Prop. 4.2 is not invertible: two cocycles that differ by multiplication by a one dimensional character give rise to the same localized endomorphism; in the irreducible case this is the only ambiguity, that could be eliminated by considering inner automorphisms rather than unitary operators. Apart from this point, $`H_\tau ^1(𝔄)`$ describes faithfully the covariant superselection sectors. Note also that, if a certain strong additivity assumption holds, all sectors with finite dimension are covariant .
Although the Roberts cohomology describes the superselection sectors, the definition of the statistical dimension is not manifest within $`H_R^1(𝔄)`$. But, passing to $`H_\tau ^1(𝔄)`$, we discover a pairing between factor KMS states $`\phi `$ for the time evolution satisfying duality and $`H_\tau ^1(𝔄)`$
$$\mathrm{log}\phi ,[z]\mathrm{log}d_\phi (z)=\mathrm{log}d(\rho )+\beta \mu _\rho (\phi )$$
hence we have the geometrical description of the holomorphic dimension by the diagram
$$\begin{array}{ccc}\text{Covariant sectors}& & H^1(𝔄)\\ & & & & \\ H_\tau ^1(𝔄)& & \{\mathrm{}\}\end{array}$$
Of course, by considering the involution as before, we obtain an expression for $`d(\rho )`$ as well.
## 5 Increment of black hole entropy as an index.
This section is devoted to the illustration of a physical context where the dimension is expressed by a quantity that includes also classical geometric data of the underlying spacetime. The results here below have been announced in .
We shall consider the increment of the entropy of a quantum black hole which is represented by a globally hyperbolic spacetime with bifurcate Killing horizon.
The case of a Rindler black hole has been studied in , the reader can find the basic ideas and motivations in this reference. The main points here are the following. Firstly we perform our analysis in the case of more realistic black holes spacetimes, in particular we consider Schwarzschild black holes. Secondly we shall consider charges localizable on the horizon, obtaining in this way quantum numbers for the increment of the entropy of the black hole itself, rather than of the outside region as done in the Rindler spacetime.
Indeed the restriction of quantum fields to the horizon will define a conformal quantum field theory on $`S^1`$. This key point has been discussed in with small variations. We shall return on this with further comments. Finally, we shall deal with general KMS states, besides the Hartle-Hawking temperature state. In this context, a non-zero chemical potential can appear.
The extension of the DHR analysis of the superselection structure to a quantum field theory on a curved globally hyperbolic spacetime has been given in and we refer to this paper the necessary background material. We however recall here the construction of conformal symmetries for the observable algebras on the horizon of the black hole.
To be more explicit let $`𝒱`$ be a $`d+1`$ dimensional globally hyperbolic spacetime with a bifurcate Killing horizon. A typical example is given by the Schwarzschild-Kruskal manifold that, by Birkoff theorem, is the only spherically symmetric solution of the Einstein-Hilbert equation; one might first focus on this specific example, as the more general case is treated similarly. We denote by $`𝔥_+`$ and $`𝔥_{}`$ the two codimension 2 submanifolds that constitute the horizon $`𝔥=𝔥_+𝔥_{}`$. We assume that the horizon splits $`𝒱`$ in four connected components, the future, the past and the “left and right wedges” that we denote by $``$ and $``$ (the reader may visualize this also in the analogous case of the Minkowski spacetime, where the Killing flow is a one-parameter group of pure Lorentz transformations).
Let $`\kappa =\kappa (𝒱)`$ be the surface gravity, namely, denoting by $`\chi `$ the Killing vector field, the equation on $`𝔥`$
$$g(\chi ,\chi )=2\kappa \chi ,$$
(43)
with $`g`$ the metric tensor, defines a function $`\kappa `$ on $`𝔥`$, that is actually constant on $`𝔥`$, as can be checked by taking the Lie derivative of both sides of eq. (43) with respect to $`\chi `$ . If $`𝒱`$ is the Schwarzschild-Kruskal manifold, then
$$\kappa (𝒱)=\frac{1}{4M},$$
where $`M`$ is the mass of the black hole. In this case $``$ is the exterior of the Schwarzschild black hole.
Our spacetime is $``$ and we regard $`𝒱`$ as a completion of $``$.
Let $`𝔄(𝒪)`$ be the von Neumann algebra on a Hilbert space $``$ of the observables localized in the bounded diamond $`𝒪`$. We make the assumptions of Haag duality, properly infiniteness of $`𝔄(𝒪)`$, Borchers property B. The Killing flow $`\mathrm{\Lambda }_t`$ of $`𝒱`$ gives rise to a one parameter group of automorphisms $`\alpha `$ of the C-algebra $`𝔄=𝔄()`$ since $`𝒱`$ is a $`\mathrm{\Lambda }`$-invariant region. Here, as before, the C-algebra associated with an unbounded region is the C-algebra generated by the von Neumann algebras associated with bounded diamonds contained in the region.
### 5.1 The conformal structure on the black hole horizon.
We now consider a locally normal $`\alpha `$-invariant state $`\phi `$ on $`𝔄()`$, that restricts to a KMS at inverse temperature $`\beta >0`$ on the horizon algebra, as we will explain. This will be the case in particular if $`\phi `$ is a KMS state on all $`𝔄()`$.
For convenience, we shall assume that the net $`𝔄`$ is already in the GNS representation of $`\phi `$, hence $`\phi `$ is represented by a cyclic vector $`\xi `$. Denote by $`_a`$ the wedge $``$ “shifted by” $`a`$ along, say, $`𝔥_+`$ (see ). If $`I=(a,b)`$ is a bounded interval of $`_+`$, we set
$$𝒞(I)=𝔄(_a)^{\prime \prime }𝔄(_b)^{},0<a<b.$$
We denote by $`(I)`$ ($`I(0,\mathrm{})`$) the C-algebra generated by all $`𝒞(a,b),b>a>0`$, with $`(a,b)I`$. We shall also set $`𝒞(I)=(I)^{\prime \prime }`$.
We obtain in this way a net of von Neumann algebras on the intervals of $`(0,\mathrm{})`$, where the Killing automorphism group $`\alpha `$ acts covariantly by rescaled dilations. (With the due assumptions on duality, $`𝒞(I)`$ turns out to be the intersection of the von Neumann algebras associated with all diamonds containing the “interval $`I`$” of the horizon $`𝔥_+`$ .)
We can now state our assumption: $`\phi |_{(0,\mathrm{})}`$ is a KMS state with respect to $`\alpha `$ at inverse temperature $`\beta >0`$.
The net $`(a,b)𝒞(e^a,e^b)`$ is obviously a net on $``$ where $`\alpha `$ is the translation automorphism group. We are therefore in the setting treated in Section 3, whose results may be now applied.
In particular, by using Wiesbrock theorem , we now show that the restriction of the net to the black hole horizon $`𝔥_+`$ has many more symmetries than the original net.
###### Theorem 5.1.
(). The Hilbert space $`_0=\overline{𝒞(I)\xi }`$ is independent of the bounded open interval $`I`$.
The net $`𝒞`$ extends to a conformal net $``$ of von Neumann algebras acting on $`_0`$, where the Killing flow corresponds to the rescaled dilations.
###### Proof.
Setting $`𝔄(a,b)𝒞(e^a,e^b)`$, $`b>a`$, we trivially obtain a local net on $``$ and $`\phi `$ is a KMS with respect to translations. The conformal extension of $`𝒞`$ is then the conformal completion of $`𝔄`$ given by Th. 3.3; the additivity assumption is here unnecessary since $`_0=\overline{𝒞(I)\xi }`$ is independent of $`I`$. A detailed proof can be found in Prop. 4A.2 of . $`\mathrm{}`$
This theorem says in particular that we may compactify $``$ to the circle $`S^1`$, extend the definition of $`𝒞(I)`$ for all proper intervals $`IS^1`$, find a unitary positive energy representation of the Möbius group $`\text{PSL}(2,)`$ acting covariantly on $`𝒞`$, so that the rescaled dilation subgroup is the Killing automorphism group.
If $`\xi `$ is cyclic for $``$ on $``$, namely $`_0=`$ (as is true in Rindler case for a free field, see ), then the net $`𝒞`$ automatically satisfies Haag duality on $``$. Otherwise one would pass to the dual net $`𝒞^d`$ of $`𝒞`$, which is is automatically conformal and strongly additive . The following proposition gives a condition for $`𝒞`$ itself to be strongly additive.
###### Proposition 5.2.
If $`\phi `$ is KMS on $`𝔄()`$ and the strong additivity for $`𝔄`$ holds in the sense that
$$𝒞(0,1)𝔄(_1)^{\prime \prime }=𝔄()^{\prime \prime },$$
then $`𝒞`$ satisfies Haag duality on $``$.
###### Proof.
Note first that there is a conditional expectation $`\epsilon :𝔄(_a)^{\prime \prime }𝒞(a,\mathrm{})`$, $`a`$, given by $`\epsilon (X)\xi =EX\xi `$, $`X𝔄(_a)^{\prime \prime }`$, where $`E`$ is the orthogonal projection onto $`_0`$.
The case $`a=0`$ is clearly a consequence of Takesaki’s theorem since $`𝒞(0,\mathrm{})`$ is globally invariant under the modular group of $`𝔄()^{\prime \prime }`$. Now the translations on $`𝔄`$ (constructed as in Prop. 3.2) restrict to the translations on $`𝒞`$ and commute with $`\epsilon `$ due to the cyclicity of $`\xi `$ for $`(a,\mathrm{})`$ on $`_0`$. Hence $`\epsilon `$ maps $`𝔄(_a)^{\prime \prime }`$ onto $`𝒞(a,\mathrm{})`$, $`a`$.
We then have
$$\begin{array}{c}𝒞(0,1)𝒞(1,\mathrm{})=𝒞(0,1)\epsilon (𝔄(_1)^{\prime \prime })\hfill \\ \hfill =\epsilon (𝒞(0,1)𝔄(_1)^{\prime \prime })=\epsilon (𝔄()^{\prime \prime })=𝒞(0,\mathrm{}),\end{array}$$
(44)
hence $`𝒞`$ is strongly additive, thus it satisfies Haag duality on $``$ because $`𝒞`$ is conformal. $`\mathrm{}`$
In the following we assume that $`𝔄`$ is strongly additive.
### 5.2 Charges localizable on the horizon.
We now consider an irreducible endomorphism $`\rho `$ with finite dimension of $`𝔄()`$ that is localizable in an interval $`(a,b)`$ of $`𝔥_+`$, $`a>0`$, namely $`\rho `$ acts trivially on $`𝔄(_b)`$ and on $`(0,a)`$, thus it restricts to a localized endomorphism of $`𝒞`$.
This last requirement is necessary to extend $`\rho `$ to a normal endomorphism of $`\pi _\phi (𝔄())^{\prime \prime }`$, with $`\phi `$ as above or a different extremal KMS state.
Remark. If we assume that the net $`𝔄`$ is defined on all $`𝒱`$ and that $`\rho `$ is a transportable localized endomorphism of $`𝔄(𝒱)`$, then $`\rho `$ has a normal extension to the von Neumann algebra of $`𝔄()^{\prime \prime }`$ exists by transportability, cf. , and the strong additivity assumption is unnecessary. This case may be treated just as the case of the Rindler spacetime , the only difference being the possible appearance a non-trivial chemical potential. We omit the detailed discussion of this context.
By a result in , a transportable localized endomorphism with finite dimension $`\rho `$ of $`𝒞`$ is Möbius covariant. In particular we may choose the covariance cocycle so that it verifies the two-variable cocycle property with respect to the the action of the Möbius group, and this uniquely fixes it. Let $`\sigma `$ be another irreducible endomorphism of $`𝔄`$ localized in $`(a,b)𝔥_+`$ and denote by $`\phi _\rho `$ and $`\phi _\sigma `$ the thermal states for the Killing automorphism group in the representation $`\rho `$. As shown in , $`\phi _\rho =\phi \mathrm{\Phi }_\rho `$, where $`\mathrm{\Phi }_\rho `$ is the left inverse of $`\rho `$, and similarly for $`\sigma `$.
We now show that the dimension of $`\rho `$ and of its restriction to $`𝒞`$ coincide. In fact the following is true.
###### Lemma 5.3.
With the above assumptions, if $`\rho `$ is localized in an interval of $`𝔥_+`$, then $`d(\rho |_{(0,\mathrm{})})`$ has a normal extension to the weak closure $`𝒞(0,\mathrm{})`$ with dimension $`d(\rho )`$.
###### Proof.
If $`\rho `$ is localized in the interval $`(a,b)`$ ($`b>a>0`$) of $`𝔥_+`$, then clearly $`\rho `$ restricts to the von Neumann algebra $`𝒞(c,d)`$ for all $`d>b>a>c`$ as it acts trivially on $`𝔄(_c)`$ and by duality for $`𝒞`$. Hence $`\rho `$ restricts to the C-algebra $`(0,\mathrm{})`$, and then it extends to $`𝒞(0,\mathrm{})`$ by Theorem 3.7.
Let $`\overline{\rho }`$ be also localized in the interval $`(a,b)`$. If $`R,\overline{R}`$ are a standard solutions for the conjugate equation of $`\rho `$ and $`\overline{\rho }`$, then $`R,\overline{R}𝔄(_b)^{}𝔄()^{\prime \prime }=𝒞(0,b)`$ and commute with $`𝒞(0,a)`$, hence they belong to $`𝒞(a,b)`$ by strong additivity . Conversely, if $`R,\overline{R}`$ are a standard solution for the conjugate equation of the restriction of $`\rho `$ and $`\overline{\rho }`$, then $`R`$ and $`\overline{R}`$ belongs to $`𝒞(a,b)`$ by the Haag duality for $`𝒞`$, hence the conjugate equation is valid for all elements of $`𝒞(0,b)𝔄(_b)^{\prime \prime }=𝔄()`$. This implies the dimension is the same for $`\rho `$ and its restriction to $`𝒞`$. $`\mathrm{}`$
The increment of the free energy between the thermal equilibrium states $`\phi _\rho `$ and $`\phi _\sigma `$ is expressed as in by
$$F(\phi _\rho |\phi _\sigma )=\phi _\rho (H_{\rho \overline{\sigma }})\beta ^1S(\phi _\rho |\phi _\sigma )$$
Here $`S`$ is the Araki relative entropy and $`H_{\rho \overline{\sigma }}`$ is the Hamiltonian on $`_0`$ corresponding to the composition of the charge $`\rho `$ and the charge conjugate to $`\sigma `$ in $`𝒞(\mathrm{},0)`$ as in ; it is well defined as $`e^{itH_{\rho \overline{\sigma }}}e^{itH_\iota }`$ is the two-variable cocycle. In particular, if $`\sigma `$ is the identity representation, then $`H_{\rho \overline{\sigma }}=H_\rho `$ is the Killing Hamiltonian in the representation $`\rho `$.
The analysis made in Section 2.1 works also in this context, and in fact the formulae there can be written in the case of two different KMS states $`\phi _\rho `$ and $`\phi _\sigma `$. In particular
$$F(\phi _\sigma |\phi _\rho )=\frac{1}{2}\beta ^1(S_c(\sigma )S_c(\rho ))+\mu (\phi _\sigma |\phi _\rho ),$$
(45)
where $`\mu (\phi _\sigma |\phi _\rho )=\beta ^1\mathrm{log}d_\phi (u_\sigma )\beta ^1\mathrm{log}d_\phi (u_\rho )`$. Here $`S_c(\rho )=\mathrm{log}d(\rho )^2`$ is the conditional entropy associated with $`\rho `$ (see ) and the above formula gives a canonical splitting for the incremental free energy.
### 5.3 An index formula.
We now consider the Hartle-Hawking state $`\phi `$, see . In several cases this is the unique KMS state for the Killing evolution. The corresponding Hawking temperature is is related to the surface gravity of $``$:
$$\beta ^1=\frac{\kappa ()}{2\pi }.$$
###### Theorem 5.4.
With $`\phi `$ the Hartle-Hawking state, if $`\rho `$ and $`\sigma `$ are localizable as above, then
$$\mathrm{log}d(\rho )\mathrm{log}d(\sigma )=\frac{\pi }{\kappa ()}(F(\phi _\rho |\phi _\sigma )+F(\phi _{\overline{\rho }}|\phi _{\overline{\sigma }})),$$
in particular, once we fix $`\rho `$, the exponential of the right hand side of the above equation is proportional to an integer.
###### Proof.
Formula (45) states that
$$\mathrm{log}d(\rho )\mathrm{log}d(\sigma )=\beta F(\phi _\rho |\phi _\sigma )\beta \mu (\phi _\sigma |\phi _\rho ).$$
(46)
By the asymmetry of the chemical potential $`\mu (\phi _\sigma |\phi _\rho )+\mu (\phi _{\overline{\sigma }}|\phi _{\overline{\rho }})=0`$, thus
$$\mathrm{log}d(\rho )\mathrm{log}d(\sigma )=\beta F(\phi _{\overline{\rho }}|\phi _{\overline{\sigma }})+\beta \mu (\phi _\sigma |\phi _\rho ).$$
(47)
Summing up equations (46) and (47) and setting $`\beta /2=\pi /\kappa ()`$ we obtain
$$\mathrm{log}d(\rho )\mathrm{log}d(\sigma )=\mathrm{log}d_{geo}(\rho )\mathrm{log}d_{geo}(\sigma )=\frac{\pi }{\kappa ()}(F(\phi _\rho |\phi _\sigma )+F(\phi _{\overline{\rho }}|\phi _{\overline{\sigma }})).$$
$`\mathrm{}`$
We have thus expressed the analytical index $`\mathrm{log}d(\rho )\mathrm{log}d(\sigma )`$ in terms of a physical quantity, the incremental free energy, and the quantity $`\kappa ()`$ associated with the geometry of the spacetime.
###### Corollary 5.5.
$$F(\phi _\sigma |\phi _\rho )+F(\phi _{\overline{\sigma }}|\phi _{\overline{\rho }})=\beta ^1(S_c(\sigma )S_c(\rho ))$$
where $`S_c(\rho )`$ denotes the conditional entropy of the sector $`\rho `$.
###### Proof.
Immediate. $`\mathrm{}`$
The above corollory shows that the part of the incremental free energy which is independent of charge conjugation is proportional to the increment of the conditional entropy.
## 6 On the index of the supercharge operator.
The purpose of this section is to make a few remarks to interpret the statistical dimension of a superselection sector as the Fredholm index of an operator associated with the supercharge operator in a supersymmetric theory.
Let $`𝔉(𝒪)`$ be the von Neumann algebra on a Hilbert space $``$ generated by the fields localized in the region $`𝒪𝒦`$ of a spacetime $`𝕄`$, say the Minkowski spacetime, in a Quantum Field Theory, as in Section 2, in the vacuum representation. Let $`𝔉=_{𝒪𝒦}𝔉(𝒪)^{}`$ be the quasi-local C-algebra and $`\gamma :gG\gamma _g\text{Aut}(𝔉)`$ an action of a compact group of internal symmetries with $`g_0`$ an involutive element of the center of $`G`$ providing a grading automorphism $`\gamma _0=\gamma _{g_0}`$ with normal commutation relations:
$$F_1F_2\pm F_2F_1=0,F_i𝔉(𝒪_i),𝒪_1𝒪_2^{},$$
where $`F_1,F_2𝔉_\pm \{F:\gamma _0(F)=\pm F\}`$, and the $`+`$ sign occurs iff both $`F_1`$ and $`F_2`$ belong to $`𝔉_{}`$.
With $`\mathrm{\Gamma }`$ the canonical selfadjoint unitary on $``$ implementing $`\gamma _0`$, the Hilbert space decomposes according to the eigenvalues of $`\mathrm{\Gamma }`$
$$=_+_{}.$$
(48)
As is known , each irreducible representation $`\pi \widehat{G}`$ gives rise to a superselection sector $`\rho _\pi `$ of the observable algebra $`𝔄𝔉^G`$: given a Hilbert space of isometries $`H_\pi 𝔉(𝒪)`$ carrying the representation $`\pi `$, one has a covariant endomorphism $`\rho _\pi `$ of $`𝔄`$
$$\rho _\pi (X)=\underset{i=1}{\overset{d}{}}v_iXv_i^{},X𝔉,$$
(49)
where $`\{v_1,v_2,\mathrm{}v_d\}`$ is an orthonormal basis of $`_\pi `$ and
$$d=\text{dim}(\pi )=d_{DHR}(\rho _\pi ).$$
The unitary representation of $`G`$ implementing $`\gamma `$ gives raise to a decomposition of $``$
$$=\underset{\pi \widehat{G}}{}_\pi $$
and, denoting by $`\pi _0`$ the identity representation of $`𝔄`$ on $``$, one has the unitary equivalence
$$\pi _0|__\pi \text{dim}(\pi )\pi _0\rho _\pi |__\iota .$$
Let $`H`$ be the Hamiltonian of $`𝔉`$, namely the generator of the time evolution on $``$. We now assume the existence of a supersymmetric structure on $`𝔉`$, namely there exists a supercharge operator $`Q`$, an odd selfadjoint operator with $`Q^2=H`$; corresponding to the decomposition (48) of $``$, $`Q`$ can be written as
$$\left(\begin{array}{cc}0& Q_+\\ Q_{}& 0\end{array}\right)$$
with $`Q_{}=Q_+^{}`$, where $`Q_\pm :_\pm _{}`$. In case (not assumed here) that $`e^{\beta H}`$ is a trace class operator, $`\beta >0`$, one has the well-known formula
$$\text{Index}(Q_+)=\text{Tr}(\mathrm{\Gamma }e^{\beta H}),\beta >0,$$
with $`\text{Index}(Q_+)`$ the Fredholm index of $`Q_+`$, that turns out to coincide with the dimension of the kernel of $`H`$, thus $`\text{Index}(Q_+)=1`$ if there exists a unique vacuum vector.
Now fix an irreducible localized endomorphism $`\rho =\rho _\pi `$ of $`𝔄`$ as above and let $`H_\rho `$ be the Hamiltonian in the representation $`\rho `$, namely the generator of the unitary time evolution in the representation $`\rho `$. The Hamiltonian $`H_\rho `$ is not supersymmetric, in the sense that there exists no odd square root of $`H_\rho `$, as such an operator would interchange Bose and Fermi sectors and cannot map the Hilbert space of $`\rho `$ (which is either contained in $`_+`$ or in $`_{}`$) into itself. Yet, we may define a supercharge operator $`Q_\rho `$ on the global Hilbert space $``$ by setting
$$Q_\rho =\underset{i}{}v_iQv_i^{},$$
with the $`v_i`$’s forming a basis of $`H_\pi `$ as above.
###### Lemma 6.1.
$`Q_\rho ^2|__0=H_\rho `$.
###### Proof.
Indeed
$$Q_\rho ^2=(\underset{i}{}v_iQv_i^{})^2=\underset{i}{}v_iQ^2v_i^{}=\underset{i}{}v_iHv_i^{}.$$
(50)
Because of the expression (49), one checks easily that $`U(t)\rho (X)U(t)=\rho (\alpha _t(X))`$, $`X𝔄`$, where $`U(t)=_iv_ie^{itH}v_i^{}`$ and $`\alpha `$ is the one-parameter automorphism group. Since $`H`$ commutes with $`\gamma _G`$, it follows that $`_0`$ is an invariant subspace for $`_iv_iHv_i^{}`$, thus the latter restricts to $`H_\rho `$ on $`_0`$ and formula (50) implies the statement. $`\mathrm{}`$
We now assume that $`\rho `$ is a Bose sector, thus $`H_\pi 𝔉_+`$, namely each $`v_i`$ commutes with $`\mathrm{\Gamma }`$ and thus preserves the decomposition (48) of the Hilbert space so that
$$Q_\rho =\left(\begin{array}{cc}0& Q_{\rho +}\\ Q_\rho & 0\end{array}\right)$$
where in particular $`Q_{\rho +}=_iv_iQ_+v_i^{}`$. Since the restriction of $`v_iQ_+v_i^{}`$ to $`v_iv_i^{}_+`$ is unitarily equivalent to $`Q_+`$, clearly $`Q_{\rho +}`$ is unitarily equivalent to $`Q_+\mathrm{}Q_+`$ ($`d`$ times) and therefore
$$\text{Index}(Q_{\rho +})=d(\rho )\text{Index}(Q_+).$$
If $`\sigma `$ is another sector as above, we thus have the formula
$$\frac{d(\rho )}{d(\sigma )}=\frac{\text{Index}(Q_{\rho +})}{\text{Index}(Q_{\sigma +})}$$
showing an interpretation of the (statistical) dimension as a multiplicative relative Fredholm index.
It remains to provide a model where the above structure can be realized. To this end let $`𝔄_b`$ and $`𝔄_f`$ be the nets of local algebras associated with the free scalar Bose field and the free Fermi-Dirac field on the Minkowski spacetime, or on the cylinder spacetime. Then $`𝔉=𝔄_b𝔄_f`$ is equipped with a supersymmetric structure (see e.g. ), where the grading unitary is $`1(1)^N`$, with $`(1)^N`$ the even-odd symmetry. We fix a positive integer $`n2`$ and consider
$$\stackrel{~}{𝔉}𝔉𝔉\mathrm{}𝔉,(n\text{factors}),$$
(51)
as our field algebra, with the gauge group $`G=_n`$ acting by permuting the order of the bosonic algebra $`𝔄_b\mathrm{}𝔄_b`$ in (51). The only thing to check is the the existence of a supersymmetric structure, which is given by the following lemma.
###### Lemma 6.2.
The tensor product of two supersymmetric QFT nets is supersymmetric.
###### Proof.
To simplify notations we shall consider the tensor product $`\stackrel{~}{𝔉}=𝔉𝔉`$ of the same net $`𝔉`$ by itself. The decomposition (48) of the Hilbert space $``$ of $`𝔉`$ gives a decomposition of $`\stackrel{~}{}=`$
$`\stackrel{~}{}_+`$ $`=(_+_+)(_{}_{})`$ (52)
$`\stackrel{~}{}_{}`$ $`=(_+_{})(_{}_+),`$ (53)
and we can define the operators $`\stackrel{~}{Q}_\pm :\stackrel{~}{}_\pm _{}`$
$`\stackrel{~}{Q}_+`$ $`=(Q_+1+1Q_+)(Q_{}11Q_{})`$ (54)
$`\stackrel{~}{Q}_{}`$ $`=(Q_+1+1Q_{})(Q_{}11Q_+).`$ (55)
The tensor product Hamiltonian $`\stackrel{~}{H}=(H1)(1H)`$ is equal to $`\stackrel{~}{Q}_{}\stackrel{~}{Q}_+\stackrel{~}{Q}_+\stackrel{~}{Q}_{}`$. $`\mathrm{}`$
It should be noticed that the choice of the tensor product supercharge $`\stackrel{~}{Q}`$ is not canonical: having chosen $`\stackrel{~}{Q}`$ we get another supercharge $`\stackrel{~}{Q}_g`$ by permuting the order of the tensor product factors with a permutation $`g`$.
Remark. If we apply the above procedure to chiral conformal field theory, with $`G`$ a finite group of internal symmetries, the observable algebra $`𝔄=𝔉^G`$ has an extra family of sectors $`\{\sigma _i\}_i`$, beside the above ones $`\{\rho _\pi \}_{\pi \widehat{G}}`$ ; indeed $`_\pi d(\rho _\pi )^2=|G|`$, while
$$\underset{\pi \widehat{G}}{}d(\rho _\pi )^2+\underset{i}{}d(\sigma _i)^2=|G|^2.$$
The dimension $`d(\sigma _i)`$ is not necessarily integral and therefore cannot be related to a Fredholm index. Our formulae, nevertheless, still make sense.
## 7 Outlook. Comparison with the JLO theory.
This section contains a tentative proposal to analyze the superselection sectors by noncommutative geometry. Although it is in a primitive form, we hope it will provide an insight to the structure.
### 7.1 Induction of cyclic cocycles.
Let $`𝔄`$ be a $`_2`$-graded unital pre-C-algebra, with C completion $`\overline{𝔄}`$. Morphisms of $`𝔄`$ will be assumed to be bounded, i.e. to extend to $`\overline{𝔄}`$. We assume that $`𝔄`$ is a Banach -algebra with respect to a norm $`||||||`$ preserved by the grading $`\gamma `$, so that $`𝔄`$ is a graded Banach -algebra.
We briefly recall some basic definitions about Connes entire cyclic cohomology. Let $`𝒞^n(𝔄)`$ be the Banach space of the $`(n+1)`$-linear functionals of $`𝔄`$ with finite norm
$$|f_n|=\underset{|a_i|1}{sup}|f_n(a_0,a_1,\mathrm{},a_n)|$$
and let $`𝒞(𝔄)`$ be the space of the entire cochains, namely the elements of $`𝒞(𝔄)`$ are the sequences $`f=(f_0,f_1,f_2,\mathrm{})`$, $`f_n𝒞^n(𝔄)`$, such that
$$f=\underset{n0}{}\sqrt{n!}|f_n|z^n$$
(56)
is an entire function of $`z`$.
The grading $`\gamma `$ lifts to $`𝒞(𝔄)`$, and let $`𝒞_\pm (𝔄)`$ denote the corresponding splitting of the entire cochains.
Denoting by $`𝒞^e`$ and $`𝒞^o`$ the spaces of the even and odd entire cochains $`(f_0,f_2,f_4,\mathrm{})`$ and $`(f_1,f_3,f_5,\mathrm{})`$, the entire cohomology groups $`H_+^e(𝔄)`$ and $`H_+^o(𝔄)`$ are the ones associated with the complex
$$\mathrm{}𝒞_+^e\stackrel{}{}𝒞_+^o\stackrel{}{}𝒞_+^e\mathrm{}$$
(57)
where the coboundary operator is $`=b+B`$
$$\begin{array}{c}(Bf)_{n1}(a_0,a_1,\mathrm{},a_{n1})=\underset{j=0}{\overset{n1}{}}(1)^{(n1)j}f_n(1,a_{nj}^\gamma ,\mathrm{},a_{n1}^\gamma ,a_0\mathrm{},a_{nj1})\hfill \\ \hfill +(1)^{n1}f_n(a_{nj}^\gamma ,\mathrm{},a_{nj1},1)\end{array}$$
$$(bf)_{n+1}=\underset{j=0}{\overset{n}{}}f_n(a_0,\mathrm{},a_ja_{j+1},\mathrm{},a_{n+1})+(1)^{n+1}f_n(a_{n+1}^\gamma a_0,a_1,\mathrm{},a_n).$$
If $`𝔅`$ is another -Banach algebra, and graded pre-C-algebra, and $`\rho :𝔄𝔅`$ a homomorphism of $`𝔄`$ into $`𝔅`$, any cyclic cocycle $`\tau `$ of $`𝔅`$ has a pull-back to a cyclic cocycle $`(\tau \rho )`$ of $`𝔄`$
$$(\tau \rho )_n(a_0,a_1,\mathrm{},a_n)=\tau _n(\rho (a_0),\rho (a_1),\mathrm{},\rho (a_n)).$$
In particular, let $`𝔅`$ be contained in $`𝔄`$, where both $`\overline{𝔄}`$ and $`\overline{𝔅}`$ are unital with trivial center, and let $`\lambda `$ be a canonical endomorphism of $`𝔄`$ with respect to a $`𝔅`$, namely eq. (61) holds with $`T𝔄`$ and $`S𝔅`$, see the appendix. In this case, if $`(\tau _n)`$ is a cyclic cocycle of $`𝔅`$, its pull-back to $`𝔄`$ via $`\lambda `$ will be called the induction of $`(\tau _n)`$ to $`𝔄`$.
###### Proposition 7.1.
The induction gives rise to a well-defined map from $`H^1(𝔅)`$ to $`H^1(𝔄)`$, independently of the choice of $`\lambda `$.
###### Proof.
Immediate by the Proposition A.2 and the fact that inner automorphisms give the identity map in cyclic cohomology . $`\mathrm{}`$
Let $`𝒯`$ be a tensor category, with conjugates and subobjects, of bounded endomorphisms of $`𝔄`$ and let $`\rho `$ be an object of $`𝒯`$. If $`(\tau _n)`$ is a cyclic cocycle of $`𝔄`$, then we obtain a cyclic cocycle $`(\tau _n^\rho )(\tau _n\rho ^1)`$ on $`𝔅=\rho (𝔄)`$ as pull back by $`\rho ^1`$, hence a cyclic cocycle of $`𝔄`$ by inducting it up to $`𝔄`$.
###### Proposition 7.2.
The induction of $`(\tau _n\rho ^1)`$ from $`\rho (𝔄)`$ to $`𝔄`$ is equivalent to $`(\tau _n\overline{\rho })`$, namely to the pull-back of $`(\tau _n)`$ via the conjugate charge $`\overline{\rho }`$.
###### Proof.
Since $`\lambda =\rho \overline{\rho }`$ is the canonical endomorphism of $`𝔄`$ into $`\rho (𝔄)`$ (see the appendix), the result is immediate. $`\mathrm{}`$
Remark. There is a product on the localized cyclic cocycles (i.e. cocycles of the form $`\tau \rho `$, for some localized endomorphism $`\rho `$), by setting $`(\tau \rho )(\tau \sigma )\tau \rho \sigma `$. This product passes to equivalence classes and gives a well defined product in cohomology. The interest in this point relies on the fact that in the passage from the commutative to the noncommutative case the ring structure in cohomology is usually lost (viewed on the K-theoretical side, there is no noncommutative version of the tensor product of fiber bundles).
### 7.2 Index and super-KMS functionals.
Let $`𝔄_b`$, $`𝔄_f`$ be pre-C-algebras with C-completions $`\overline{𝔄}_b`$, $`\overline{𝔄}_f`$, and let $`𝔄𝔄_b𝔄_f`$ be their tensor product. We assume that $`\overline{𝔄}_f`$ is $`_2`$-graded, while the grading on $`𝔄_b`$ is trivial. Let $`\gamma `$ the involutive automorphism of $`𝔄`$ giving the grading on $`𝔄`$, trivial on $`𝔄_b`$. Let $`\alpha \alpha ^{(b)}\alpha ^{(f)}`$ be a one-parameter group of automorphisms of $`\overline{𝔄}`$ leaving invariant $`𝔄`$ and $`\delta `$ be an unbounded odd derivation of $`\overline{𝔄}`$, with a dense -algebra $`𝔄_\alpha =𝔄_{\alpha ,b}𝔄_{\alpha ,f}𝔄`$ contained in the domain of $`\delta `$, so that $`\delta `$ is a square root of the generator $`D`$ of $`\alpha `$
$$\delta ^2=D\frac{\text{d}}{\text{d}t}\alpha _t|_{t=0}\mathrm{on}𝔄_\alpha ,$$
cf. . We consider now a super-KMS functional $`\phi `$ at inverse temperature $`\beta =1`$, namely $`\phi `$ is a linear functional on $`𝔄`$ that satisfies equation (24) for all $`a,b𝔄`$ and
$$\phi (\delta a)=0,a𝔄_\alpha .$$
We assume that $`\phi `$ is bounded, as is the case for quantum field theory on a compact space, where $`\phi `$ is a super-Gibbs functionals. At the end of this section we shall add comments on the case $`\phi `$ unbounded. Associated with $`\phi `$ there is an entire cyclic cocycle, the JLO cocycle , on the algebra $`𝔄_\alpha `$ normed with the norm $`|a|a+\delta a`$. For $`n`$ even, it is defined as
$$\begin{array}{c}\tau _n(a_0,a_1,\mathrm{},a_n)\hfill \\ \hfill (1)^{\frac{n}{2}}_{0t_1\mathrm{}t_n1}\phi (a_0\alpha _{it_1}(\delta a_1^\gamma )\alpha _{it_2}(\delta a_2)\mathrm{}\alpha _{it_n}(\delta a_n^{\gamma ^n}))dt_1dt_2\mathrm{}dt_n.\end{array}$$
(58)
There is a special class of unitary $`\alpha `$-cocycles with finite holomorphic dimension that correspond to bounded perturbations of the dynamics . If $`q𝔄`$ is even there is a perturbed structure $`(𝔄,\alpha ^q,\delta _q)`$, with $`\delta _q=\delta +[,q]`$ and $`\alpha ^q=\text{Ad}u^q\alpha `$, where $`u^q`$ is unitary cocycle in $`𝔄`$ given by the solution of
$$i\frac{d}{dt}u^q(t)=u^q(t)\alpha _t(h),h=\delta q+q^2;$$
the functional
$$\phi ^q(a)=\underset{ti}{\text{anal.cont. }}\phi (au^q(t))$$
is super-KMS with respect to $`\alpha ^q`$ and the topological index of $`u^q`$ is given by the Chern character, the JLO cocycle $`\tau ^q`$ associated with $`\phi ^q`$ evaluated at the identity
$$d_\phi (u^q)=\phi ^q(1)\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{(2n)!}{n!}\tau _{2n}^q(1,1,\mathrm{},1).$$
Of course, the second equality is trivially true as $`\delta 1=0`$, thus only the first term in the series may be non-zero.
We state the result on the deformation invariance if the topological index in :
###### Theorem 7.3.
(). If $`q𝔄_{}`$ then $`\phi ^q(1)=d_\phi (u^q)=d_\phi (1)=\phi (1)`$.
In the case of Th. 7.3, $`(\tau _n^q)`$ is indeed cohomologous to $`(\tau _n)`$.
We consider now a different class of unitary cocycles. We specialize to the case where $`𝔄_b𝔄_f`$ is dense in the quasi-local C-algebra associated with a supersymmetric quantum field theory (cf. Sect. 6). Let $`𝒯`$ be a tensor category, with conjugate and subobjects, of localized endomorphisms contained in End$`(𝔄_{\alpha ,b})`$. We keep the same symbol for the endomorphism $`\rho `$ of $`𝔄_b`$ to denote the endomorphism $`\rho \iota `$ of $`\overline{𝔄}`$ and assume each $`\rho `$ to be $`\alpha `$-covariant, in the sense that there is a $`\alpha `$-cocycle of unitaries $`u(\rho ,t)𝔄_{\alpha ,b}`$ such that is a eq. (14) holds.
Suppose first that $`\rho `$ is an automorphism. Then, setting $`\delta ^\rho \rho \delta \rho ^1`$, we see that $`\delta ^\rho \delta `$ acts identically where $`\rho `$ acts identically, that is, by Haag duality, $`\delta ^\rho `$ is a perturbation of $`\delta `$ by a derivation localized in a double cone.
In general, when $`\rho `$ is an endomorphism, $`\delta ^\rho `$ is defined on $`\rho (𝔄_\alpha )`$. We shall then define the multilinear form $`\tau _n^\rho `$ on $`\rho (𝔄_\alpha )`$
$$\begin{array}{c}\tau _n^\rho (a_0,a_1,\mathrm{},a_n)(1)^{\frac{n}{2}}_{\mathrm{\Sigma }_n}\phi (a_0u(is_1)\alpha _{is_1}(\delta ^\rho (a_1^\gamma )u(is_2))\alpha _{i(s_1+s_2)}(\delta ^\rho (a_2)u(is_3))\hfill \\ \hfill \mathrm{}u(is_n)\alpha _{i(s_1+\mathrm{}+s_n)}(\delta ^\rho (a_n^{\gamma ^n})u(is_{n+1})))\mathrm{d}s_1\mathrm{}\mathrm{d}s_{n+1}\end{array}$$
(59)
where $`\mathrm{\Sigma }_n\{(s_1,\mathrm{},s_{n+1}):s_i0,s_1+\mathrm{}+s_{n+1}=1\}`$ and $`u(t)u_\rho (t)`$. We now assume that $`\phi `$ is of the form $`\phi _b\phi _f`$, which is automatically the case if $`\phi `$ is factorial with a cluster property, and that $`\phi _b`$ satisfies Haag duality, so that we can apply the results in Sect. 2.
As an irreducible object $`\rho `$ of $`𝒯`$ extends to irreducible endomorphisms of the weak closure on the of $`𝔄_{\alpha ,b}`$ in the state $`\phi `$ (Cor. 2.8), then
###### Proposition 7.4.
$$\tau _n^\rho =d_\phi (u_\rho )\tau _n\rho ^1$$
on $`\rho (𝔄_{\alpha ,b})`$.
###### Proof.
Setting
$$\phi _\rho =\underset{ti}{\text{anal.cont. }}\phi (u(t)),$$
we have (see Sect. 1.1 and ) $`\phi _\rho =d_\phi (u_\rho )\phi \mathrm{\Phi }_\rho `$, thus
$$\phi _\rho |_{\rho (𝔄)}=d_\phi (u_\rho )\phi \rho ^1,$$
which is a super-KMS functional with respect to the evolution Ad$`u(t)\alpha _t=\rho \alpha _t\rho ^1`$ and the superderivation $`\delta ^\rho =\rho \delta \rho ^1`$.
A direct verification shows that $`(\tau _n^\rho )`$ is the JLO cocycle on $`\rho (𝔄_\alpha )`$ associated with $`\phi _\rho |_{\rho (𝔄_\alpha )}`$. $`\mathrm{}`$
Thus the above expression is well defined, if $`(\tau _n)`$ is well defined. In order to have a cocycle on all $`𝔄_{\alpha ,b}`$ we consider the induction $`(\stackrel{~}{\tau }_n^\rho )`$ of $`(\tau _n^\rho )`$ to $`𝔄_\alpha `$.
By the previous discussion and Prop. 7.4 we then have
$$\stackrel{~}{\tau }_n^\rho =d_\phi (u_\rho )\tau _n\overline{\rho }.$$
We summarize our discussion in the following corollary.
###### Corollary 7.5.
Consider the supersymmetric structure as above on $`𝔄=𝔄_b𝔄_f`$ with $`\overline{𝔄}_b`$ the quasi-local C-algebra as in Section 2.
If $`\phi =\phi _b\phi _f`$ is a supersymmetric KMS functional satisfying Haag duality and $`\rho `$ is a translation covariant localized endomorphism of $`𝔄_b`$ mapping $`𝔄_{\alpha ,b}`$ into itself, then
$$d_\phi (\rho )=\stackrel{~}{\tau }_0^\rho (1)\left(\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{(2n)!}{n!}\stackrel{~}{\tau }_{2n}^\rho (1,1,\mathrm{},1)\right)$$
and
$$d_{geo}(\rho )=\sqrt{d_\phi (u_\rho )d_\phi (u_{\overline{\rho }})}=d_{DHR}(\rho ).$$
Remark. The discussion in this section is incomplete in two respects. On one hand if we consider a Quantum Field Theory on a compact non-simply connected low dimensional spacetime, as a chiral conformal net on $`S^1`$, then the quasi-local C-algebra has non-trivial center with a non-trivial action of the superselection structure ; our results then need to be extended to this context with a further study, possibly with a connection with the low dimensional QFT topology. On the other hand, when dealing with quantum field theory on the Minkowski spacetime, as in Section 2, the boundedness requirement for the graded KMS functional $`\phi `$ should be omitted because of Prop. 1.10. Thus the construction of $`(\stackrel{~}{\tau }_n^\rho )`$, in its actual form, is not satisfactory except, perhaps, for QFT on a contractible curved spacetime. This difficulty vanishes if we drop the boundedness requirement for $`\phi `$, but only demand the restriction $`\phi _b\phi |_{𝔄_b}`$ to be bounded. It is then natural to deal only with the restricion of the JLO cocycle $`(\tau _n)`$ to the Bosonic algebra $`𝔄_{\alpha ,b}`$ . An example of such an unbounded super-KMS functionals seems to occur naturally . However a general study unbounded super-KMS functional does not presently exist, in particular it should be check that the JLO formula still gives a well defined and entire cyclic cocycle $`(\tau _n)`$.
## Appendix A Appendix. Some properties of sectors in the C-case.
For the convenience of the reader we describe here a few facts concerning endomorphisms of C-algebras, which are natural counterparts of the corresponding results in the context of factors.
### A.1 The canonical endomorphism.
Let $`𝔄`$ be a unital C-algebra with trivial centre. An endomorphism $`\lambda `$ of $`𝔄`$ is called *canonical* (with finite index) if there exist isometries $`T(\iota ,\lambda )`$ and $`S(\lambda ,\lambda ^2)`$ such that
$$\begin{array}{cc}\hfill \lambda (S)S=& S^2\hfill \\ \hfill S^{}\lambda (T)\backslash \{0\},& T^{}S\backslash \{0\}.\hfill \end{array}$$
(60)
If $`𝔅𝔄`$ is a unital C-subalgebra and $`\lambda `$ is an endomorphism of $`𝔄`$ with $`\lambda (𝔄)𝔅`$ we shall say that $`\lambda `$ is a canonical endomorphism of $`𝔄`$ with respect to $`𝔅`$ (or into $`𝔅`$) if there exist intertwiners $`T(\iota ,\lambda )`$, $`S(\iota |_𝔅,\lambda |_𝔅)`$ such that
$$S^{}\lambda (T)\backslash \{0\},T^{}S\backslash \{0\}.$$
(61)
In this case $`\lambda `$ is clearly canonical, in the sense that equation (60) holds, since $`S𝔅`$. The converse is also true.
###### Proposition A.1.
If $`\lambda `$ is canonical (eq. 60), then it is canonical with respect to a natural C-subalgebra $`𝔅𝔄`$ (eq. 61).
The proof is the same as in the factor case : one defines $`𝔅`$ as the range of $`S^{}\lambda ()S`$ and checks by eq. (60) that $`𝔅`$ is a C-subalgebra and that $`S^{}\lambda ()S`$ is a conditional expectation of $`𝔄`$ onto $`𝔅`$. The above definitions extend to the case of pre-C-algebras, in this case we assume that endomorphisms are bounded and the intertwiners are assumed to live in the -algebras.
Note now that if $`\rho `$ is an endomorphism of $`𝔄`$, a conjugate $`\overline{\rho }`$ of $`\rho `$, if exists, is unique as sector, i.e. up to inner automorphisms of $`𝔄`$, in fact a more general result holds in the context of tensor categories . This implies the following.
###### Proposition A.2.
Let $`𝔅𝔄`$ be unital pre-C-algebras with trivial centre and $`\lambda _1,\lambda _2`$ canonical endomorphisms of $`𝔄`$ into $`𝔅`$ (namely eq. (61) hold). There exists a unitary $`u𝔅`$ such that $`\lambda _2=\text{Ad}u\lambda _1`$.
###### Proof.
The proof could be given similarly to the one given in , Propositions 4.1 and 4.2. However, it is easier to observe that one can consider sectors and conjugate sectors between different C-algebras. Then the canonical endomorphism $`\lambda `$, as a map of $`𝔄`$ into $`𝔅`$, is just the conjugate sector for the embedding $`\iota `$ of $`𝔅`$ into $`𝔄`$ and thus the uniqueness of $`\lambda `$ modulo inner automorphisms of $`𝔅`$ is just a consequence of the uniqueness of the conjugate in a tensor 2-C-category, see . $`\mathrm{}`$
###### Proposition A.3.
If $`\rho `$ and $`\overline{\rho }`$ are conjugate as above, then $`\lambda =\rho \overline{\rho }`$ is the canonical endomorphism of $`𝔄`$ with respect to the subalgebra $`\rho (𝔄)`$.
###### Proof.
Let $`R,\overline{R}`$ be the operators in the conjugate equation for $`\rho ,\overline{\rho }`$; setting $`T=\overline{R}`$ and $`S=\rho (R)`$ equation (61) holds true. $`\mathrm{}`$
### A.2 The quantum double in the C setting.
We consider now a construction in the C context, that corresponds to the one given in in the context of factors, see also .
Let $`𝔄`$ be a unital C-algebra with trivial centre and $`𝒯\mathrm{End}(𝔄)`$ a tensor category of endomorphisms with conjugates and sub-objects. We shall denote by $`𝔄^{op}`$ the opposite C-algebra and by $`\stackrel{~}{𝔄}𝔄𝔄^{op}`$ the tensor product with respect to the maximal C tensor norm.
Let $`I`$ be an index set and choose $`\{\rho _i\}_{iI}`$ a family all irreducible objects of $`𝒯`$, one for each equivalence class, with $`\rho _0=\iota `$ and $`\rho _{\overline{i}}=\overline{\rho }_i`$ and set $`\stackrel{~}{\rho }_i\rho _i\rho _i^{op}`$, where $`\rho ^{op}j\rho j`$ with $`j:𝔄𝔄^{op}`$ the canonical anti-linear isomorphism.
Let $`𝔅_0`$ be the linear space of functions $`X:I\stackrel{~}{𝔄}`$ with finite support. We consider the following product and -operation on $`𝔅_0`$:
$$XY(k)\underset{i,j}{}X(i)\stackrel{~}{\rho }_i(Y(j))C_{ij}^k,$$
(62)
$$X^{}(k)C_{k\overline{k}}^0\stackrel{~}{\rho }_k(X(\overline{k})^{}).$$
(63)
Here $`C_{ij}^k(\stackrel{~}{\rho }_k,\stackrel{~}{\rho }_i\stackrel{~}{\rho }_j)`$ is the canonical intertwiner $`C_{ij}^k=\sqrt{\frac{d(\rho _i)d(\rho _j)}{d(\rho _k)}}_{\mathrm{}}v_{\mathrm{}}j(v_{\mathrm{}})`$ with $`\{v_{\mathrm{}}\}_{\mathrm{}}`$ any orthonormal bases in $`(\rho _k,\rho _i\rho _j)`$. Clearly $`\stackrel{~}{𝔄}`$ can be identified with the subalgebra of $`𝔅_0`$ of functions with support in $`0`$.
Setting $`R_i(k)\delta _{ik}`$, the $`R_i\stackrel{~}{𝔄}`$ satisfy the relations
$$\{\begin{array}{c}R_iX=\stackrel{~}{\rho }_i(X)R_i,X\stackrel{~}{𝔄},\hfill \\ R_i^{}R_i=d(\rho _i)^2,\hfill \\ R_iR_j=_kC_{ij}^kR_k,\hfill \\ R_i^{}=C_{\overline{i}i}^0R_{\overline{i}},\hfill \end{array}$$
(64)
where we have omitted the $``$ in the product. $`𝔅_0`$ is the a -algebra and every $`X𝔅_0`$ as a unique expansion
$$X=\underset{i}{}X(i)R_i$$
where the coefficients are uniquely determined by $`X(i)=\epsilon (XR_i^{})`$. Here $`\epsilon :𝔅_0\stackrel{~}{𝔄}`$ is the conditional expectation given by $`\epsilon (X)=X(0)`$, that can be shown to be faithful as in , App. A.
Extending states from $`\stackrel{~}{𝔄}`$ to $`𝔅_0`$ via $`\epsilon `$ we obtain a faithful family of states, hence $`𝔅_0`$ has a maximal C-norm, the completion under which is a C-algebra $`𝔅`$. Now
$$X^{}X\underset{\psi }{sup}\psi (X^{}X)\underset{\phi _1\phi _2}{sup}\phi _1\phi _2\epsilon (X^{}X)=\epsilon (X^{}X)$$
where $`\psi `$ ranges over the states of $`𝔅_0`$ and $`\phi _1\phi _2`$ over the product states of $`\stackrel{~}{𝔄}`$, namely $`\epsilon `$ is bounded and extends to a conditional expectation $`\epsilon :𝔅\stackrel{~}{𝔄}`$.
To each $`X𝔅`$ we may associate the formal expansion
$$X=\underset{i}{}X(i)R_i$$
where $`X(i)\epsilon (XR_i^{})`$ and one has $`X=0X(i)=0iI`$.
We have $`\stackrel{~}{𝔄}^{}𝔅=`$. Indeed if $`X\stackrel{~}{𝔄}^{}𝔅`$, then for every $`a\stackrel{~}{𝔄}`$
$$\underset{i}{}aX(i)R_i=aX=Xa=\underset{i}{}X(i)R_ia=\underset{i}{}X(i)\stackrel{~}{\rho }_i(a)R_i,$$
thus $`X(i)(\iota ,\stackrel{~}{\rho }_i)`$, thus it follows by the irreducibility of $`\stackrel{~}{\rho }_i`$ that $`X=X(0)𝔄^{}𝔄=`$.
The contact with the construction in is visible by the following proposition.
###### Proposition A.4.
If $`𝒯`$ is rational (namely $`I`$ is a finite set), then
$$\lambda =\underset{iI}{}\rho _i\rho _i^{op}$$
is the restriction to $`\stackrel{~}{𝔄}`$ of the canonical endomorphism of $`𝔅`$ to $`\stackrel{~}{𝔄}`$.
Although the chemical potential can be described via extension of KMS state from $`\stackrel{~}{𝔄}`$ to $`𝔅`$, in Section 3.3 we prefer to deal with a slight variation of the above construction.
Indeed the definition of $`𝔅`$ can be modified by replacing $`\stackrel{~}{𝔄}`$ with $`𝔄`$, where $``$ is any C-algebra with trivial centre where $`𝒯`$ acts faithfully, namely $`𝒯\mathrm{End}()`$. Accordingly one puts in this case $`\stackrel{~}{\rho }_i\rho _i\overline{\rho }_i`$ and $`C_{ij}^k=_{\mathrm{}}\sqrt{\frac{d(\rho _i)d(\rho _j)}{d(\rho _k)}}v_{\mathrm{}}v_{\mathrm{}}^{}`$. In particular we may take $`=𝔄`$ and indeed we specialize to this case in section 3.3. It is rather obvious how to formulate in the above setting the structure and the results obtained for the quantum double case.
Acknowledgments. Among others, we wish to thank in particular S. Doplicher for early motivational comments, and A. Connes, K. Fredenhagen and A. Jaffe for inspiring conversations and invitations respectively at the IHES, Hamburg University and Harvard University, while this work was at different stages. |
warning/0003/cond-mat0003508.html | ar5iv | text | # Test of molecular mode coupling theory: A first resume
## I Introduction
The mode coupling theory (MCT) for supercooled simple liquids proposed by Bengtzelius, Götze and Sjölander interprets the glass transition as a dynamical transition. This picture has been supported by many experiments on several glass formers (see e.g. and Refs. therein) and more recently by detailed analysis of computer simulations for Lennard-Jones systems (see e.g. and references therein). The signature of the dynamical transition, i.e. the asymptotical power laws, have been discovered also via neutron scattering, light scattering, dielectric relaxation and NMR in molecular glassformers (see and references therein), stimulating the extension of MCT to molecular liquids.
Two approaches for such an extension have been proposed recently for rigid molecules. Chong and Hirata introduced a theory based on a site-site description of the molecules. Their approach offers the advantage to be closely related to neutron scattering experiments. The structural information — which is a necessary input of the theory — can be readily obtained from theories of molecular liquids able to predict partial structure factors, like the RISM approximation. The second approach is based on the expansion of the orientational density into a complete set of functions, in analogy to the Fourier expansion of the density related to the translational degrees of freedom. To distinguish the second approach from the site-site theory, it is called molecular mode coupling theory (MMCT). MMCT was derived for a single linear molecule in a simple liquid , for liquids of linear molecules (for some application see ) and for molecules of arbitrary shape . MMCT allows to calculate the glassy dynamics for all orientational degrees of freedom, and it is closely connected to dielectric and NMR experiments. As reorientational motion is also very important in light scattering experiments it may also be helpful in their interpretation. The connection of MMCT to neutron scattering experiments has been discussed recently .
The fundamental MMCT quantities are the time-dependent correlation functions
$$S_{ln,l^{}n^{}}(q,m,t)=\rho _{ln}^{}(q,m,t)\rho _{l^{}n^{}}(q,m)$$
(1)
of the tensorial density modes
$$\rho _{ln}(q,m,t)=i^l(2l+1)^{\frac{1}{2}}\underset{j=1}{\overset{N}{}}e^{i\stackrel{}{q}\stackrel{}{x}_j(t)}𝒟_{mn}^l(\mathrm{\Omega }_j(t)).$$
(2)
Here it is $`𝐪=(0,0,q)`$ and $`l`$ runs over all positive integers including zero, $`m`$ and $`n`$ take integer values between $`l`$ and $`l`$, $`𝒟`$ denotes the Wigner functions . The reader should note that the correlators (1) are diagonal in $`m`$ for $`𝐪=(0,0,q)`$. The MMCT equations of motion for the Laplace transform $`𝐒(q,m,z)=i_0^{\mathrm{}}𝐒(q,m,t)e^{izt}`$, $`(\text{Im }z>0)`$ have been presented in a preceding paper. Here we focus on the unnormalized molecular nonergodicity parameters
$$𝐅(q,m)=\underset{t\mathrm{}}{lim}𝐒(q,m,t)=\underset{z0}{lim}z𝐒(q,m,z)$$
(3)
which we calculated using the following set of equations .
$`𝐅(q,m)=`$ (4)
$`\left[𝐒^1(q,m)+𝐒^1(q,m)𝐊(q,m)𝐒^1(q,m)\right]^1`$ (5)
$`K_{ln,l^{}n^{}}(q,m)=`$ (6)
$`{\displaystyle \underset{\alpha \alpha ^{}}{}}{\displaystyle \underset{\mu \mu ^{}}{}}q_{ln}^{\alpha \mu }(q)q_{l^{}n^{}}^{\alpha ^{}\mu ^{}}(q)\left(\underset{¯}{𝐦}^1(q)\right)_{lmn,l^{}mn^{}}^{\alpha \mu ,\alpha ^{}\mu ^{}}`$ (7)
$`m_{lmn,l^{}m^{}n^{}}^{\alpha \mu ,\alpha ^{}\mu ^{}}(q)={\displaystyle \frac{\rho _0}{(8\pi ^2)^3}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑q_1{\displaystyle \underset{|qq_1|}{\overset{q+q_1}{}}}𝑑q_2`$ (8)
$`{\displaystyle \underset{m_1m_2}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{l_1l_1^{}}{l_2l_2^{}}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{n_1n_1^{}}{n_2n_2^{}}}{}}v_{ln,l_1n_1,l_2n_2}^{\alpha \mu }(qq_1q_2;mm_1m_2)`$ (9)
$`v_{l^{}n^{},l_1^{}n_1^{},l_2^{}n_2^{}}^{\alpha ^{}\mu ^{}}(qq_1q_2;m^{}m_1m_2)F_{l_1n_1,l_1^{}n_1^{}}(q_1,m_1)`$ (10)
$`F_{l_2n_2,l_2^{}n_2^{}}(q_2,m_2)`$ (11)
The memory function matrix $`\underset{¯}{𝐦}`$ in Eq.(11) represent the mode coupling approximation for the correlation function of fluctuating forces. The index $`\alpha `$ labels the translational ($`\alpha =T`$) and rotational ($`\alpha =R`$) currents, each of them consisting of three (spherical) vector components $`\mu \{1,0,1\}`$. The vertex functions $`v`$ are determined only by the matrix of the static molecular structure factors $`𝐒(q,m)`$ and the number density $`\rho _0`$. Their explicit form has been given in Ref. . The coefficients $`q_{ln}^{\alpha \mu }(q)`$ appearing in Eq.(7) are
$$q_{ln}^{\alpha \mu }(q)=\{\begin{array}{cc}0& \alpha =T,\mu =\pm 1\\ q& \alpha =T,\mu =0\\ \frac{1}{\sqrt{2}}\sqrt{l(l+1)n(n+\mu )}& \alpha =R,\mu \pm 1\\ n& \alpha =R,\mu =0\end{array}.$$
(12)
Since $`q_{ln}^{T\pm 1}=0`$ (due to the choice of $`𝐪=(0,0,q)`$) the transversal translational components ($`\alpha =T,\mu =\pm 1`$) of the memory functions enter only indirectly (via the inversion of $`\underset{¯}{𝐦}(q)`$) and thus shall be neglected in the following.
The given set of equations (5)-(11) includes all interactions between translational and rotational degrees of freedom in molecular liquids and thus this set is rather involved. Obviously its numerical solution poses a formidable task.
As model system for our analysis we have chosen SPC/E water . The water molecule possesses a twofold rotational symmetry ($`C_{2v}`$) around the axis given by its dipole moment which has been chosen as the z-axis of the body fixed frame.
The z-axis and the x-axis define the plane which is spanned by the molecule. As discussed in detail in Ref. the $`C_{2v}`$-symmetry can be used to simplify the equations of motion by the restriction that $`n`$ and $`n^{}`$ are even. Preceeding publications on long time MD–simulations for this strong glass former showed that the centre of mass dynamics is in good agreement with the predictions of the asymptotic laws of MCT. Recent work demonstrated that the signature of the dynamic transition can also be observed for the orientational degrees of freedom of the molecule.
A first approach to solve Eqs. (5)-(11) has been given in Ref. . Apart from the necessary truncation of the range of $`l`$ by a cutoff $`l_{co}`$ for which we have chosen $`l_{co}=2`$ we introduced in two more approximations. In the dipole approximation, water was treated as a linear molecule oriented in the direction of its dipole moment. This approximation corresponds to set in Eqs. (5)-(11) the angular indices $`n`$, $`n^{}`$ and the corresponding summation indices $`n_i`$,$`n_i^{}`$ to zero. In Ref. we have also studied the stronger diagonal-dipole approximation, which makes the further assumption that structure factors, nonergodicity parameters and memory functions are diagonal with respect to the angular indices $`l`$ and $`l^{}`$. It is the main purpose of the present communication to present results obtained by solving the ”full” set of MMCT-equations up to $`l_{co}=2`$ without any further approximation. The main difference to the approach in ref. is that we take the nonlinear character of the water molecule serious. In particular, this means that the rotational motion of both protons around the dipolar axis is taken into account. Therefore, we can study the influence of this degree of freedom on the ideal glass transition temperature $`T_c`$ and on the nonergodicity parameters $`F_{l0,l^{}0}(q,m)`$ obtained in ref. . In addition we also obtain the new parameters $`F_{ln,l^{}n^{}}(q,m)`$ with $`n`$ and/or $`n^{}`$ different from zero.
## II Results
We have solved Eqs. (5)-(11) iteratively on a grid of 100 equispaced wave-vectors, up to 110.7 $`nm^1`$. One complete iteration requires 4 days of cpu time on a 533 MHz alpha workstation. The entire calculation to locate the critical temperature and the corresponding nonergodicity parameter, with a tolerance of $`2.5\times 10^4`$ per point, requested more than 250 iterations. On a dedicated 4-nodes parallel machine it required about 250 days. The number of iteration at the critical temperature was $`54`$. We found that at $`T=282K`$ the MMCT equations predict a liquid phase, while at $`T=272K`$ unambiguously a glassy state is predicted. Within the chosen tolerance, we locate $`T_c^{MCT}`$ at $`T_c^{MCT}=279K`$.
Table I summarizes the critical temperatures at which a transition from ergodic to nonergodic behavior is found. While $`T_c`$ is almost equal for both approximations used in ref. and quite close to the result of the MD-simulation we find that the critical temperature is overestimated by almost 50% using the ”full” theory. Thus we confirm our supposition in that the agreement for $`T_c`$ between simulation and both approximation schemes was fortuitous. Such finding could also be expected from the static structure factors used as input. Besides the dominating diagonal structure factors $`S_{00,00}(q,m=0)`$ and $`S_{10,10}(q,m)`$ the most prominent peaks are displayed by $`S_{10,2\pm 2}(q,m)`$ and $`S_{2\pm 2,2\pm 2}(q,m)`$ (See Figures in Ref. ), which were neglected in . The overestimation of the critical temperature is common to the MCT for simple liquids and seems to be a general deficiency of the mode coupling approximation. Thus we see that — although the overestimation of the critical temperature is not desired — it is necessary to include all static structure factors with large amplitude to get a concise description in the MMCT framework.
Figure 1 shows the comparison of the unnormalized critical nonergodicity parameters $`F_{00,00}(q,m=0)`$ for the centre of mass correlations. The oscillations of the MD-result are captured well by all of our MMCT-results. As for the critical temperature we can also observe for the critical nonergodicity parameters that the dipole approximation and the diagonal-dipole approximation lead to nearly the same result. In the vicinity of the maximum of the structure factor the agreement between theory and simulation is improved by removing the additional approximations but for the prepeak as well as for large $`q`$ both approximation schemes perform better than the ”full” theory. In the region of the prepeak and especially for the minimum between prepeak and main peak we observe that the oscillations are less pronounced and the peak positions are slightly shifted. Exactly the same behaviour can be found by a comparison of the static structure factors at the different critical temperatures. Thus the worse performance of the ”full” theory in comparison with the approximation schemes can be attributed at least partly to the overestimation of the critical temperature because of which the static input of the calculations does not reflect accurately the static structure at the ”true” $`T_c`$ of the simulation. Apart from that, it also has to be taken into account that in spite of the computational effort we have made, the fixed point of Eq. (5)-(11) can not be determined with very high precision. Therefore, the result for the critical nonergodicity parameters after 54 iterations may overestimate the exact one by a few percent. The latter reason may also be responsible for the worse performance of the ”full” theory at the minima of $`F_{00,00}(q,m=0)`$ as Nauroth found that the convergence of the iteration is extremely slow in these parts. Finally we want to mention that the disagreement at large $`q`$ between simulation and all theoretical calculations was also found for simple liquids and was considered as a shortcoming of the mode coupling approximation.
As in the case of the centre of mass correlators the $`q`$-dependence of the nonergodicity parameters $`F_{10,10}(q,m)`$ (see fig. 2) is reproduced well by the theory. Here even in the vicinity of the maximum the approximations perform better than the ”full” theory.
Figure 3 shows the comparison for the critical nonergodicity parameters $`F_{20,20}(q,m)`$. Apart from the region of large wave vectors the ”full” theory shows better agreement with the simulation than the two approximations. Thus one observes that the approximations, which neglect terms with $`n0`$ and consequently $`l2`$, have stronger effect on the $`l=2`$-correlators than on those for $`l=1`$ or $`l=0`$. Further one observes that the agreement between theory and simulation is less good for the $`l=2`$–correlators than for those with $`l=0`$ or $`l=1`$. The reader should note that the good agreement at large $`q`$ for the $`l=2`$–correlators must be considered as fortuitous because the overestimation of the nonergodicity parameters (seen for $`l=0,1`$) is compensated by a generally too small amplitude of the $`l=2`$–correlators.
The reason for the worse agreement of the $`l=2`$–correlators is their higher sensibility to the cutoff at $`l_{co}=2`$ which can be understood on a mathematical level by a closer examination of the vertices which have been given in ref. . Let us pick out one special example to illustrate this point. The vertex factor $`v_{00,20,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)`$ is responsible for the coupling of two correlators involving $`l=0`$ and $`l=2`$, respectively. It is of the form
$`v_{00,20,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)=`$ (13)
$`{\displaystyle \underset{l,n}{}}u_{00,ln,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)c_{ln,20}(q_1,0)+(12),`$ (14)
where $`u_{00,ln,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)`$ contains Clebsch-Gordan coefficients of the form $`𝒞(l,2,0;m^{},m^{},0)`$ which enforces $`l=2`$. Therefore the contribution to the memory function matrix caused by the vertex factor in Eq. (13) is exact even for $`l_{co}=2`$. The vertex factor $`v_{20,20,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)`$ instead, which describes the coupling of two correlators involving $`l=2`$ and $`l=2`$, respectively, has the form
$`v_{20,20,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)=`$ (15)
$`{\displaystyle \underset{l,n}{}}u_{20,ln,20}^{\alpha \mu }(q,q_1,q_2;0,0,0)c_{ln,20}(q_1,0)+(12).`$ (16)
The corresponding Clebsch-Gordan coefficient $`𝒞(l,2,2;m^{},m^{},0)`$ allows $`l\{0,1,2,3,4\}`$. Thus for $`l_{co}=2`$ this vertex is not entirely taken into acount, although it is responsible for the coupling of $`l=2`$-correlations. From this example we can see that the $`l=2`$-quantities are more sensitive to the cutoff than those for $`l=0`$.
The off diagonal critical nonergodicity parameters with $`n=n^{}=0`$ shown in figure 4 are reproduced well by the theory. The difference between dipole-approximation and ”full” theory are relatively small.
Figures 5 and 6 show the comparison between the new critical nonergodicity parameters with $`n0`$ and/or $`n^{}0`$ and the simulation results. Since those terms
are neglected by the approximations they can only be calculated using the ”full” set of MMCT-equations. The agreement between theory and simulation is satisfactory for almost all correlators.
In the discussion of the unnormalized critical nonergodicity parameters we have seen that the worse performance of the ”full” theory is partly due to the fact that the static input has to be taken at the ”wrong” temperature. This influences the nonergodicity parameters in two ways: $`i)`$ $`𝐒(q,m)`$ is the initial value of the correlation function $`𝐒(q,m,t)`$ whose limit for $`t\mathrm{}`$ are the unnormalized nonergodicity parameters $`F_{ln,l^{}n^{}}(q,m)`$. $`ii)`$ In the mode coupling approximation the static structure factors determine the vertices which describe the coupling between the tensorial density modes in the system. This influence can partly be eliminated by calculating the normalized critical nonergodicity parameters
$$f_{ln,l^{}n^{}}(q,m)=\frac{F_{ln,l^{}n^{}}(q,m)}{\sqrt{S_{ln,ln}(q,m)S_{l^{}n^{},l^{}n^{}}(q,m)}}.$$
(17)
A selection of the normalized diagonal nonergodicity parameters is shown in Fig.(7)-(9). Due to the normalization the variation with $`q`$ is less pronounced. We observe that the difference in amplitude between the ”full” theory and the approximations is reduced as the initial value $`𝐒(q,m,t=0)`$ of the time-dependent correlation functions is set to 1 at all wavevectors and all temperatures. Nonetheless the dipole and diagonal-dipole approximation still provide the better description of the normalized critical nonergodicity parameter.
In part this may be explained by the better convergence of the iteration for the approximation schemes, but the wrong critical temperature also has its influence as can be seen from the fact that the peak positions for the ”full” MMCT results are still slightly shifted.
## III Summary and Conclusions
From our analysis we can state the following conclusions: i) The mode coupling theory for molecular liquids overestimates the critical temperature $`T_c`$ in the same fashion as its counterpart for simple liquids. ii) The $`q`$-dependence of the critical nonergodicity parameters is well reproduced in the vicinity of the main peaks. Systematic differences exist for large wave vectors. iii) Deviations between theory and simulation are partly due to the overestimation of the critical temperature. Therfore the structural input of the theory does not reflect properly the structure of the liquid at the true glass transition temperature. These deviations still exist for the normalized nonergodicity parameters. iv) Approximation schemes taking into account only part of the correlators like the dipole and the diagonal-dipole approximation can already give a reasonable description. v) The essential correlation functions can be selected on the basis of the static structure factors. I.e. for supercooled water the most important static correlator are those with $`n0`$ and/or $`n^{}0`$ and for $`n=n^{}=0`$ those, which are diagonal in $`l`$ and $`l^{}`$.
We think that in combination with suitable approximations for the static structure factors MMCT will be a useful tool in the understanding of the glass transition in molecular liquids. The solution of the ”full” set of MMCT equations is still a formidable task with present day computers, but approximation schemes can be constructed on the basis of the importance of various static structure factors. A reduction to the most prominent correlators and $`q`$-vectors offers the possibility to construct microscopically based schematic models for the description of molecular liquids.
## IV acknowledgement
The numerical work would not have been possible without the support of INFM-PAIS-C and MURST-PRIN98. F.S. and P.T. acknowledge support also from INFM-PRA. A.L. and R.S. are grateful for financial support from the SFB-262. We also thank L. Fabbian for participating to the early stages of this project and providing us with part of the input data. |
warning/0003/math0003233.html | ar5iv | text | # Untitled Document
On the $`L^2`$-instability and $`L^2`$-controllability of steady flows of an ideal incompressible fluid
A. Shnirelman
School of Mathematical Sciences
Tel Aviv University
69978 Ramat Aviv, Israel
shnirelm@math.tau.ac.il
1. In this work we are studying the flows of an ideal incompressible fluid in a bounded 2-d domain $`M𝐑^2`$, described by the Euler equations
$`{\displaystyle \frac{u}{t}}+(u,)u+p=0;`$ (1)
$`u=0.`$ (2)
Here $`u=u(x,t),xM,t[0,T]`$, and $`u|_M`$ is tangent to $`M`$.
It has been known for a long time, that if the initial velocity field $`u(x,0)`$ is smooth, then there exists unique smooth solution $`u(x,t)`$ of the Euler equations, which is defined for all $`t𝐑`$; see\[M-P\]. The next natural question is, what may be the behavior of this solution, as $`t\mathrm{}`$. This is a problem of indefinite complexity. A restricted problem is the following: suppose that the initial flow field $`u_0(x)`$ is in close to a steady solution $`u_0(x)`$. What may happen with this flow for big $`t`$? Does it stay always close to $`u_0`$, or it can escape far away? Which flows are available, if we start from different initial velocities, close to $`u_0`$?
These are problems of a global, nonlinear perturbation theory of steady solutions of the Euler equations. The first idea is to develop a linear stability theory. The spectrum of a linearized operator is always symmetric w.r.t. both the real and the imaginary axes, for the system is Hamiltonian. Therefore we can never prove an asymptotic stability by the linear method; at best we can prove the absence of a linear instability, which, in its turn, may be a tricky business.
The true, nonlinear stability of some classes of steady flows was first proven by V. Arnold (see \[A1\], \[A2\], \[AK\]). He considered a very strong restriction on perturbation: the perturbation of the vorticity, $`\omega (0)\omega _0=\times u(,0)\times u_0()`$, should be small in $`L^2(M)`$. There are three classes of steady flows which are stable under perturbations small in this sense. The first class contains only one flow with constant vorticity; its stability is obvious. The second and the third classes consist of steady flows, corresponding to a strong local maximum, resp. minimum, of the kinetic energy on the leaf of equivortical fields in the space of all smooth velocity fields in $`M`$ (see \[AK\]).
But this theory breaks down, if we drop the condition on the vorticity perturbation, and consider all (smooth) velocity fields $`u(x,0)`$, which are close to $`u_0(x)`$ in $`L^2`$, without any conditions on the derivatives. Note that this class of perturbations is no less physically significant than the previous one, because it describes perturbations with small energy. In this case, there is apparently no obstacle preventing the flow from going far away from $`u_0`$. So, it is likely that the flow is unstable. (but this does not prove instability, for there may be some other reasons for stability, like, for example in the KAM theory. This is analogous to the situation in the 3-d Euler equations. In the 2-d case, the vorticity is transported by the flow; this means that there exist infinity of integrals of motion, namely the moments of vorticity. These integrals prevent the solution from forming singularities. In the 3-d case, the vorticity field is transformed by the flow as a frozen-in vector field, and we can’t extract additional integrals of motion from the vorticity field. This means that we don’t know any obstacle to the formation of a finite-time singularity from a smooth initial flow. But we have no examples yet of such singularities. May be, this work may give some hints.)
2. In this work we consider a weaker stability problem. Consider the Euler equations with a nonzero right hand side (i.e. external force):
$`{\displaystyle \frac{u}{t}}+(u,)u+p=f;`$ (3)
$`u=0.`$ (4)
Here $`f=f(x,t)`$ is a smooth in $`x`$ vector field, such that $`f=0`$, and $`f(x,t)|_M`$ is parallel to $`M`$. Consider the behavior of $`u(x,t)`$, if $`f`$ is small in the following sense: $`_0^Tf(,t)_{L^2}𝑑t`$ is small, where $`[0,T]`$ is the time interval (assumed to be long), where the flow is considered. For example, if $`f`$ has the form $`f(x,t)=F(x)\delta (t)`$, we return to the initial stability problem.
Definition 1.Suppose that $`u(x_1,x_2)`$ and $`v(x_1,x_2)`$ are two steady flows. We say that the force $`f`$ transfers the flow $`u`$ into the flow $`v`$ during the time interval $`[0,T]`$, if the following is true: if $`w(x,t)`$ is the solution of the nonhomogeneous Euler equations (3), (4) with the initial condition $`w(x,0)=u(x)`$, then $`w(x,T)=v(x)`$.
Consider the simplest basic steady flow, namely a parallel flow. Let $`M`$ be a strip $`0x_21`$ in the $`(x_1,x_2)`$-plane. We restrict ourselves to the flows having period $`L`$ along the $`x_2`$-axis; this period is the same for all flows that are considered below. Suppose that the velocity field $`u_0(x)`$ has the form $`(U(x_2),0)`$, where $`U`$ is a given smooth function (the velocity profile). The original problem was, for which profiles $`U`$ the flow $`u_0`$ is stable. Our main result is the following
Theorem 1. For every nontrivial (i.e. different from constant) velocity profile $`U`$ the flow $`u_0`$ is $`L^2`$-unstable. This means that for every function $`U(x_2)\mathrm{const}`$ there exists $`C>0`$, such that for every $`\epsilon >0`$ the following is true. There exist $`T>0`$ and a smooth force $`f(x,t)`$, defined in $`M\times [0,T]`$, such that $`_0^Tf(,t)_{L^2}𝑑t<\epsilon `$, and $`f`$ transfers the flow $`u_0`$ during the time interval $`[0,T]`$ into a steady flow $`u_1`$, such that $`u_0u_1_{L^2}>C`$.
So, the flow may be considerably changed by arbitrarily small force, provided the time interval is sufficiently long.
Note that if the force $`f`$ satisfies a stronger condition $`_0^T\omega (,t)_{L^2}𝑑t<\epsilon `$, where $`\omega =\times f`$ is the vorticity, then, for every Arnold stable flow $`u_0`$, the resulting perturbation at time $`t`$ will be small, too.
This theorem is implied by a much stronger assertion.
Theorem 2.Suppose that $`U(x_2)`$ and $`V(x_2)`$ are two velocity profiles, such that $`_0^1U(x_2)𝑑x_2=_0^1V(x_2)𝑑x_2`$, and $`_0^1\frac{1}{2}|U(x_2)|^2𝑑x_2=_0^1\frac{1}{2}|V(x_2)|^2𝑑x_2`$; let $`u_0(x_1,x_2)=(U(x_2),0),v_0(x_1,x_2)=(V(x_2),0)`$ be corresponding steady parallel flows (having equal momenta and energies). Then for every $`\epsilon >0`$ there exist $`T>0`$ and a smooth force $`f(x,t)`$, such that $`_0^Tf(,t)_{L^2}<\epsilon `$, and $`f`$ transfers $`u`$ into $`v`$ during the time interval $`[0,T`$\].
This means that the flow of an ideal incompressible fluid is perfectly controllable by arbitrarily small force.
3. Theorems 1, 2 are proven by an explicit construction of the flow.
Note first, that if $`U_1,U_2,\mathrm{},U_N`$ are velocity profiles, and Theorem 2 is true for every pair $`(U_i,U_{i+1})`$ of velocity profiles, then we can pass from $`U_1`$ to $`U_N`$, simply concatenating the flows connecting $`U_i`$ and $`U_{i+1}`$; thus Theorem 2 is true for the pair $`(U_1,U_N)`$. Therefore it is enough to construct the sequence of steady flows with profiles $`U_1,\mathrm{},U_N`$, and the intermediate nonsteady flows connecting every two successive steady ones.
Note also, that it is enough to construct a sequence of piecewise-smooth flows, for it is not difficult to smoothen them, so that the necessary force will have arbitrarily small norm in $`L^1(0,T;L^2(M))`$.
As a first step, we change the flow with the profile $`U=U_1`$ by a piecewise-constant profile $`U_2`$ with sufficiently small steps; this may be done by a force with arbitrarily small norm.
Thus, $`U_2(x_2)`$ is a step function, $`U_2(x_2)=U_2^{(k)}`$ for $`x_2^{(k1)}<x_2<x_2^{(k)},k=1,\mathrm{},K`$. Every next profile $`U_i`$ is also a step-wise function. We are free to subdivide the steps and change a little the values of velocity, if these changes are small enough.
Every flow $`u_k`$ is obtained by the previous one $`u_{k1}`$ by one of two operations, described in the following theorems.
Theorem 3. Let $`U(x_2)`$ be a step function, $`U(x_2)=U^{(k)}`$ for $`x_2^{(k1)}<x_2<x_2^{(k)}`$; let $`V(x_2)`$ be another step function, obtained by transposition of two adjacent segments $`[x_2^{(k1)},x_2^{(k)}]`$ and $`[x_2^{(k)},x_2^{(k+1)}]`$. Let $`u(x_1,x_2),v(x_1,x_2)`$ be parallel flows with velocity profile $`U(x_2),V(x_2)`$. Then for every $`\epsilon >0`$ there exist $`T>0`$ and a piecewise-smooth force $`f(x,t)`$, such that $`_0^Tf(,t)_{L^2}<\epsilon `$, and the force $`f`$ transfers the flow $`u`$ into the flow $`v`$ during the time interval $`[0,T]`$.
To formulate the next theorem, remind the law of an elastic collision of two bodies. Suppose that two point masses $`m_1`$ and $`m_2`$, having velocities $`u_1`$ and $`u_2`$, collide elastically. Then their velocities after collision will be $`v_1=2u_0u_1,v_2=2u_0u_2`$, where $`u_0=(m_1u_1+m_2u_2)/(m_1+m_2)`$ is the velocity of the center of masses. The transformation $`(u_1,u_2)(v_1,v_2)`$ is called a transformation of elastic collision.
Theorem 4. Suppose that the profile $`U(x_2)`$ is like in Theorem 3, and the profile $`V(x_2)`$ is equal to $`U(x_2)`$ outside the segment $`x_2^{(k1)}<x_2<x_2^{(k+1)}`$; on the last segment, $`V(x_2)=v^{(k)}`$, if $`x_2^{(k1)}<x_2<x_2^{(k)}`$, and $`V(x_2)=v^{(k+1)}`$, if $`x_2^{(k})<x_2<x_2^{(k+1)}`$, where $`(v^{(k)},v^{(k+1)})`$ is obtained from $`(u^{(k)},u^{(k+1))}`$ by the transformation of elastic collision, the lengths $`x_2^{(k)}x_2^{(k1)},x_2^{(k+1)}x_2^{(k)}`$ playing the role of masses $`m_1,m_1`$. Let $`u(x_1,x_2),v(x_1,x_2)`$ be parallel flows with profiles $`U(x_2),V(x_2)`$. Then for every $`\epsilon >0`$ there exist $`T>0`$ and a force $`f(x,t)`$, such that $`_0^Tf(,t)_{L^2}<\epsilon `$, and the force $`f`$ transfers the flow $`u`$ into flow $`v`$.
Suppose now, that $`U(x_2)`$ and $`V(x_2)`$ are two velocity profiles, having equal momenta and energies. Then it is not difficult to construct a sequence of step functions $`U_2(x_2),U_3(x_2),\mathrm{},U_N(x_2)`$, so that $`U_2`$ is $`L^2`$-close to $`U_1=U`$, $`U_N`$ is $`L^2`$-close to $`V`$, and every profile $`U_k`$ is obtained from $`U_{k1}`$ by one of two operations, described in Theorems 3 and 4. Using these theorems and the notes above, we construct a piecewise-smooth force $`f(x,t)`$, such that $`_0^Tf(,t)_{L^2}𝑑t<\epsilon `$, and $`f`$ transfers $`U`$ into $`V`$ during the time interval $`[0,T]`$.
4. Theorems 3 and 4 are proven by the variational method.
Let $`𝒟()=𝒟`$ be the group of the volume-preserving diffeomorphisms of the flow domain $`M`$. These diffeomorphisms may be identified with fluid configurations: every configuration is obtained from some fixed one by a permutation of fluid particles, which is assumed to be a smooth, volume preserving diffeomorphism. The flow is a family $`g_t`$ of elements of $`𝒟`$, depending on time $`t,0tT`$. The Lagrangian velocity of the flow is a vector-function $`V(x,t)=\frac{}{t}g_t(x)=\dot{g}_t(x)`$, while the Eulerian velocity is the vector field $`v(x,t)=\dot{g}_t(g_t^1(x))`$. The action of the flow is defined as $`J\{g_t\}_0^T=_0^T\frac{1}{2}\dot{g}_t_{L^2}^2𝑑t`$, and the length $`L\{g_t\}_0^T=_0^T\dot{g}_t_{L^2}𝑑t`$.
The solution $`u(x,t)`$ of the homogeneous Euler equations (1), (2) is an Eulerian velocity field of a geodesic trajectory $`g_t`$ on the group $`𝒟`$: $`u(x,t)=\dot{g}_t(g_t^1(x))`$, such that $`\delta J\{g_t\}_0^T=0`$, provided $`g_0,g_T`$ are fixed. This implies that also $`\delta L\{g_t\}_0^T=0`$. This is the classical Hamiltonian principle (see\[A3\]). The evident idea is to try to construct solutions of the Euler equations by fixing $`g_0,g_T𝒟`$, and looking for the shortest trajectory, connecting these fluid configurations. If the minimum is attainable, then we have constructed some nontrivial solution of the Euler equations.
But this idea does not work well. If $`g_0`$ and $`g_T`$ are $`C^2`$-close, the minimum is assumed at some smooth trajectory. But if $`g_0`$ and $`g_T`$ are far away from each other, which is the only interesting case, then it is possible that the minimum is no more attainable (see\[S\], \[A-K\]). In the 3-d case there are examples of $`g_0,g_T`$, such that for every smooth flow $`g_t`$, connecting $`g_0`$ and $`g_T`$, there exists another smooth flow $`g_t^{}`$, connecting the same fluid configurations, such that $`J\{g_t^{}\}_0^T<J\{g_t\}_0^T`$; so, the minimum is unattainable.
The existence of a minimal geodesic connecting two configurations of a 2-d fluid is neither proven nor disproven, while some physical considerations show that sometimes the minimal smooth flow does not exist .
If there is no smooth solution of the variational problem, we may look for a generalized solution, which is no longer a smooth flow, but belongs to a wider class of object. The appropriate notion of a generalized flow was introduced by Y. Brenier \[B\]. Generalized flow is a probability measure $`\mu `$ in the space $`X=C(0,T;M)`$ of all continuous trajectories in the flow domain $`M`$, satisfying the following two conditions:
1. For every $`t[0,T]`$ and every Borel set $`AM`$,
$$\mu \{x()|x(t)A\}=\mathrm{mes}A;$$
2.
$$J\{\mu \}=_X_0^T|\dot{x}(t)|^2𝑑t\mu \{dx\}<\mathrm{}$$
The meaning of the first condition is that the generalized flow is incompressible; the second condition expresses the finiteness of the mean action (and that $`\mu `$-almost all trajectories belong to $`H^1`$).
Every smooth flow may be regarded as a generalized one; but there is a lot of truly generalized flows.
The generalized variational problem may be posed as follows: given a diffeomorphism $`g𝒟`$; consider all generalized flows which, in addition to the above conditions, satisfy the following:
3. For $`\mu `$-almost all trajectories $`x(t)`$, $`x(T)=g(x(0))`$.
We are looking for a generalized flow, which satisfies all three conditions and minimizes the functional $`J\{\mu \}_0^T=_X_0^T\frac{1}{2}|\dot{x}(t)|^2𝑑t\mu \{dx\}`$.
Y.Brenier has proved, using the simple ideas of weak compactness of a family of measures and semicontinuity of the action functional in $`X`$, that this problem has a solution for every $`g𝒟`$. Simple examples show that this solution may be very far from any smooth, or even measurable, flow.
But in the 2-d case the situation is much better, because there is an additional structure. To see it, consider a smooth incompressible flow $`g_t,0tT,g_0=\mathrm{Id},g_T=g`$. Let $`Q=M\times [0,T]`$ be a cylinder in the $`(x,t)`$-space. Every trajectory $`\lambda _x=\{(g_t(x),t)\},xM`$, is a smooth curve in $`Q`$, connecting the points $`(x,0)`$ and $`(g(x),T)`$. For different points $`x,x^{}`$, the curves $`\lambda _x,\lambda _x^{}`$ do not intersect. So, the lines $`\lambda _x`$ form a braid, containing continuum of threads. Such a braid is called a smooth braid.
Now let us define a generalized braid. Let us fix a volume-preserving diffeomorphism $`g𝒟`$ and a piecewise-smooth incompressible flow $`G_t`$, such that $`G_0=\mathrm{Id},G_T=g`$. The bundle of lines $`(G_t(x),t),xM`$, is called a reference braid and denoted by $`𝐁_0`$.
Definition 1. A generalized flow $`\mu `$ is called a generalized braid, and denoted by $`𝐁`$, if it satisfies conditions 1, 2, 3 above, and the following condition
4. For any $`N`$, let us pick trajectories $`x^1(t),\mathrm{},x^N(t)`$ by random and independently, i.e. with a probability distribution $`\mu \mathrm{}\mu `$ ($`N`$ times). Then for almost all such $`N`$-tuples of trajectories, the lines $`(x^i(t),t)`$ for different $`i`$ do not intersect, and the finite braid, formed by the curves $`(x^1(t),t),\mathrm{},(x^N(t),t)`$, is isotopic to the braid, formed by the curves $`(G_t(x^1(0)),t),\mathrm{},(G_t(x^N(0)),t)`$ (these braids have the same endpoints, so it is possible to define their isotopy).
The braid $`𝐁`$ is called a braid weakly isotopic to a piecewise-smooth braid $`𝐁_0`$.
We are discussing the following variational problem: given a map $`g`$ and a reference braid $`𝐁_0`$; find a generalized braid $`𝐁`$, isotopic to $`𝐁_0`$, which minimizes the functional $`J\{𝐁\}`$.
Theorem 5. The variational problem has a solution for every data $`g,𝐁_0`$.
To prove this theorem, we consider a sequence $`𝐁_i`$ of braids, such that $`J\{𝐁_i\}_0^TJ_0`$, where $`J_0=infJ\{𝐁\}`$ for all braids $`𝐁`$, isotopic to a given braid $`𝐁_0`$. This sequence is weakly compact; its subsequence converges to a generalized flow $`\mu `$, such that $`J\{\mu \}=J_0`$, exactly as for the generalized flows. But the generslized flow $`\mu `$ is, in fact, a braid isotopic to $`𝐁_0`$, which we shall denote by $`\overline{𝐁}`$. This is implied by the fact that the isotopy class of a finite subbraid with fixed endpoints is ”weakly continuous”, and therefore we can pass to the limit and conclude that the weak limit of the braids $`𝐁_i`$, regarded as generalized flows, is a braid, isotopic to $`𝐁_0`$. Let us call $`\overline{𝐁}`$ a minimal braid.
Generally, braids are as nonregular locally as generalized flows. In particular, they have, in general, no definite velocity field: for almost every $`(x,t)`$ there are different trajectories passing through this point with different velocities. But the minimal braid is much more regular. Recall that a measurable flow is defined as a family $`h_{s,t}`$ of measurable maps of $`M`$ into itself, preserving the Lebesgue measure, and such that $`h_{s,t}h_{r,s}=h_{r,t}`$ for all $`r,s,t[0,T]`$.
Theorem 6. Let $`\overline{𝐁}`$ be a minimal braid, isotopic to the reference braid $`𝐁_0`$. Then there exists a measurable flow $`h_{s,t}`$ in $`M`$, such that for $`\mu `$-almost all trajectories $`x(t)`$, $`x(t)=h_{s,t}x(s)`$. Moreover, there exists a vector field $`u(x,t)L^2`$, divergence free and tangent to $`M`$, such that for almost all trajectories $`x(t)`$, $`\dot{x}(t)=u(x(t),t)`$ for almost all $`t`$.
The next fact about the minimal braids is the following
Theorem 7. The velocity field $`u(x,t)`$, corresponding to the minimal braid $`\overline{𝐁}`$, is a weak solution of the Euler equations. This means that for every vector field $`v(x,t)C_0^{\mathrm{}}`$, such that $`v=0`$, and for every scalar function $`\phi (x,t)`$,
$`{\displaystyle \underset{Q}{}}\left[(u,{\displaystyle \frac{v}{t}})+(uu,v)\right]𝑑x𝑑t=0,`$ (5)
$`{\displaystyle \underset{Q}{}}(u,\phi )𝑑x𝑑t=0.`$ (6)
The last fact which we need is the following approximation theorem.
Theorem 8. Suppose that $`\overline{𝐁}`$ is a minimal braid, and $`u(x,t)`$ is its velocity field. Then for every $`\epsilon >0`$ there exists a smooth incompressible flow with velocity field $`w(x,t)`$, which is a solution of the nonhomogeneous Euler equations with the force $`f(x,t)`$, such that $`w(x,t)u(x,t)_{L^2}<\epsilon `$ for all $`t`$, and $`_0^Tf(,t)_{l^2}𝑑t<\epsilon `$.
5. The brais are used to construct the flows, described in Theorems 3 and 4. Consider a flow with a piecewise-constant profile $`U(x_2)`$; then the cylinder $`Q=M\times [0,T]`$ may be divided into slices $`Q_k`$, so that $`U|_{Q_k}=U^{(k)}`$. In every such slice the trajectories are parallel lines with the same slope. Thus, they form a simple, piecewise-smooth braid.
Now let us describe the braids corresponding to the flows described in Theorem 3. Let us divide the domains $`M_k`$, $`M_{k+1}`$, the bases of the cylinders $`Q_k`$, $`Q_{k+1}`$, into small subdomains $`M_{k,j}`$, $`M_{k+1,l}`$. Let us pick one point $`x_{k,j}`$ in every domain $`M_{k,j}`$, and one point $`x_{k+1,l}`$ in every domain $`M_{k+1,l}`$. Let $`\lambda _{k,j},\lambda _{k+1,l}`$ be the trajectories, passing through the points $`(x_{k,j},0)`$ and $`(x_{k+1,l},0)`$. Their endpoints in $`M\times \{T\}`$ are denoted by $`(y_{k,j},T)`$ and $`(y_{k+1,l},T)`$. The trajectories, passing through $`M_{k,j}\times \{0\}`$, and through $`M_{k+1,l}\times \{0\}`$, form subbraids $`𝐁_{k,j}`$ and $`𝐁_{k+1,l}`$.
Now let us define a new braid $`𝐁_0^{}`$. First let us define its threads $`\lambda _{k,j}^{}`$ and $`\lambda _{k+1,l}^{}`$, passing through the points $`(x_{k,j},0)`$ and $`(x_{k+1,l},0)`$. They are straight lines, passing through the points $`(y_{k,j}^{},T)`$ and $`(y_{k+1,l}^{},T)`$, obtained from the points $`(y_{k,j},T)`$ and $`(y_{k+1,l},T)`$ by the shift in the $`x_2`$-direction by, resp., $`(x_2^{(k+1)}x_2^{(k)})`$ and $`(x_2^{(k1)}x_2^{(k)})`$. These lines do not, generally, intersect.
Now let us define a piecewise-smooth braid $`𝐁_0^{}`$. It coincides with $`𝐁`$ outside $`Q_kQ_{k+1}`$. In the last domain the braid $`b_0^{}`$ consists of smooth incompressible subbraids $`𝐁_{k,j}^{}`$ and $`𝐁_{k+1,l}^{}`$ with the bases $`M_{k,j}\times \{0\}`$ and $`M_{k+1,l}\times \{0\}`$. Each subbraid contains one line $`\lambda _{k,j}^{}`$ and $`\lambda _{k+1,l}^{}`$. The interfaces between these subbraids are piecewise-smooth. It is easy to construct such subbraids, while they are not unique.
Now let us use $`𝐁_0^{}`$ as a reference braid, and construct a minimal braid $`𝐁^{}`$, isotopic to $`𝐁_0^{}`$. Using Theorems 5–8, we construct a smooth flow $`w(x,t)`$, supported by a smooth force $`f(x,t)`$, such that $`_0^Tf(,t)_{L^2}𝑑t`$ is arbitrarily small, provided $`T`$ is big enough. This flow is $`L^2`$-close to $`U`$ at $`t=0`$ and to $`V`$ at $`t=T`$; after a small modification of $`w(x,t)`$, requiring an $`L^2`$-small correcting force, we obtain a flow described in Theorem 3.
The proof of Theorem 4 is similar; but in this case the curves $`\lambda _{k,j}^{}`$ and $`\lambda _{k+1,l}^{}`$ are going from points $`(x_{k,j},0)((x_{k+1,l},0))`$ to the points $`(y_{k,j},0)((y_{k+1,l},0))`$ , and some of the curves $`\lambda _{k,j}^{}`$ and $`\lambda _{k+1,l}^{}`$ are linked.
6. Theorems 3 and 4 are true also for circular flows in a disk, with the angular momentum staying in place of momentum in Theorem 4. But for generic 2-d domains the situation is not so clear. We don’t know, whether there is an integral of motion, similar to the angular momentum, in any domain different from the disk. If such integral does not exist, which is the most likely, then the natural conjecture is that for any two flows with equal energies the conclusion of Theorem 4 is true. But this behavior is paradoxical: just imagine a nearly circular flow in a nearly circular domain (e.g. ellipse), which after some long time changes the sign of the angular velocity. This question requires more thinking.
Acknowledges. This work was done during my stay at the Department of Mathematics of Princeton University in the spring semester 1999, with the support of the American Institute of Mathematics. I am very thankful to these institutions for creative atmosphere and generous support.
I am thankful to the organizers of the Journees des Equations Differentielles for invitation to this excellent conference.
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warning/0003/cond-mat0003436.html | ar5iv | text | # On phase transitions in two-dimensional disordered systems
\[
## Abstract
The low-energy limits of models with disorder are frequently described by sigma models. In two dimensions, most sigma models admit either a Wess-Zumino-Witten term or a theta term. When such a term is present the model can have a stable critical point with gapless excitations. We describe how such a critical point appears, in particular in two-dimensional superconductors with disorder. The presence of such terms is required by the underlying (anomalous) symmetries of the original electron model. This indicates that the usual symmetry classes of disordered systems in two dimensions can be further refined. Conversely, our results also indicate that models previously thought to be in different universality classes are in fact the same once the appropriate extra terms are included.
\]
Understanding the effects of disorder is one of the central problems of condensed matter physics. It is natural to focus upon the phase transitions these systems undergo. In this paper, we will describe several ways phase transitions and the corresponding gapless modes arise in disordered models in two dimensions.
The low-energy degrees of freedom for non-interacting fermions with disorder can frequently be described by a sigma model in the replica formulation . A $`G/H`$ sigma model arises when the vacuum manifold of a field is $`H`$, a subgroup of the global symmetry group $`G`$. The low-energy modes of this model then take values in the coset $`G/H`$. Such sigma models in two dimensions do not generically have gapless degrees of freedom. For finite number of replicas, this is a result of the Mermin-Wagner-Coleman theorem. The sigma models are usually gapped because the space $`G/H`$ is curved. Expanding the action in powers of the curvature gives a mass scale: this is the well-known story of asymptotic freedom and dimensional transmutation. At specific values of $`N`$, the curvature and hence the beta function can vanish , but in this paper we will be concerned with the general situation.
Despite the fact that sigma models are usually gapped, there are many non-trivial fixed points in two dimensions. One can add to the sigma model action terms which dramatically change the low-energy physics. One famous example is the topological theta term. For example, the model with $`G=O(3)`$ and $`H=O(2)`$ has a gap unless one adds a theta term to the action. This $`O(3)`$-invariant term is a total derivative, so it depends only on the fields at spacetime infinity and has no effect on perturbation theory around the unstable Gaussian fixed point. Explicitly, the theta term takes the form $`in\theta `$, where $`n`$ is an integer counting how many times a field configuration winds around the $`O(3)/O(2)`$ two-sphere. As has been discussed in the context of the Heisenberg spin-chain, at $`\theta =\pi `$ there is a non-trivial low-energy fixed point . It is unstable when $`\theta `$ is not $`\pi `$, so tuning $`\theta `$ through $`\pi `$ results in a phase transition. This behavior inspired the well-known proposal that the Hall plateau transition is described by $`N0`$ replica limit of the $`U(2N)/U(N)\times U(N)`$ sigma model with a $`\theta `$ term .
The theta angle need not be a tunable parameter. If there is a symmetry of the model (like time-reversal or $`CP`$ invariance) under which the theta term is odd, then $`\theta `$ is fixed at zero or $`\pi `$. Moreover, in certain sigma models such as $`SU(N)/SO(N)`$, the winding number can only be $`0`$ or $`1`$, so $`\theta `$ can be only zero or $`\pi `$. As we will discuss below, in such situations $`\theta `$ is determined by the underlying disordered model. When $`\theta =\pi `$, the $`SU(N)/SO(N)`$ model has a stable low-energy fixed points: it flows to $`SU(N)_1`$ . Thus the fixed value of $`\theta `$ determines whether the model is gapless or not.
There is another way to modify some two-dimensional sigma models to obtain a stable low-energy fixed point. This is to add a Wess-Zumino-Witten (WZW) term to the action . In four-dimensional particle physics, the WZW term for example describes the $`\pi ^02\gamma `$ decay in low-energy effective field theory . Such terms arise when the fields take values in a group. Consider a theory with global symmetry $`H\times H`$, where $`H`$ is any simple Lie group (e.g. $`SU(N),SO(N)`$, or $`Sp(2N)`$). If the vacuum manifold of the theory is chosen such that the modes in the diagonal “vector” subgroup $`H_V`$ become gapped, the low-energy fields, $`h`$, take values in the remaining space, $`H\times H/H_V`$. This space is isomorphic to $`H`$ itself, and is called the axial part of the original symmetry group. For example, for $`SU(N)`$, $`h`$ would be an $`N\times N`$ unitary matrix with determinant one.
Without the WZW term, a sigma model where the fields take values in a simple Lie group is called a principal chiral model. In terms of the coupling, $`g`$ (the inverse stiffness), its action is
$$S_{PCM}=\frac{1}{g}𝑑x𝑑y\text{Tr}\left[_\mu h^\mu h^1\right].$$
(1)
We assume that the fields all fall to a constant at spatial infinity, so that the two-dimensional spacetime can be treated as a sphere. To write the WZW term, one needs to extend the fields $`h(x,y)`$ on the sphere to fields $`h(x,y,z)`$ on a ball which has the original sphere as a boundary. The fields inside the ball are defined so that $`h`$ at the center is the identity matrix, while $`h`$ on the boundary is the original $`h(x,y)`$. It is possible to find a continuous deformation of $`h(x,y)`$ to the identity because $`\pi _2(H)=0`$ for any simple Lie group. Then the WZW term is $`k\mathrm{\Gamma }`$, where
$$\mathrm{\Gamma }=\frac{ϵ_{ijk}}{24\pi }𝑑x𝑑y𝑑z\text{Tr}\left[(h^1_ih)(h^1_jh)(h^1_kh)\right].$$
(2)
The coefficient $`k`$ is known as the level, and for compact groups must be an integer because the different possible extensions of $`h(x,y)`$ to the ball yield a possible ambiguity of $`2\pi `$ times an integer in $`\mathrm{\Gamma }`$. Unlike the $`\theta `$ term, the WZW term does change the equations of motion, but only by terms involving $`h(x,y)`$: the variation of the integrand is a total derivative in $`z`$.
The sigma model with WZW term has a stable fixed point at $`g=16\pi /k`$ , so the model is critical and the quasiparticles are gapless. The corresponding conformal field theory is known as the $`H_k`$ WZW model . The WZW term is invariant under discrete parity transformations (e.g. $`xx`$, $`yy`$) provided $`h`$ is simultaneously transformed via $`hh^1`$. In other words, the WZW term in a parity-invariant theory involves pseudoscalars.
The WZW and theta terms are closely related. For example, if one adds a term $`\lambda (\text{Tr}h)^2`$ to the $`SU(2)_k`$ WZW action, at $`\lambda `$ large the low-energy modes live on the sphere, i.e. the $`SU(2)/O(2)`$ sigma model. One can check explicitly that the WZW term then turns in to the theta term with $`\theta =k\pi `$ . Moreover, in sigma models with only winding number $`0`$ or $`1`$, the topological term can be written in the form of the WZW term (2) .
Given the above considerations, it is important to understand how a WZW term arises in the replica formulation of two-dimensional models with disorder. We illustrate the appearance first by considering a simple theory where such a term appears, related to models studied in . We start with a system of spin-polarized fermions with a triplet p-wave type pairing:
$$H_0=\underset{k}{}ϵ_kc_k^{}c_k+(\mathrm{\Delta }_kc_k^{}c_k^{}+\mathrm{h}.\mathrm{c}.),$$
(3)
where $`\mathrm{\Delta }_k\frac{v_\mathrm{\Delta }}{2}k_x`$. We study disorder weak enough to maintain some notion of the Fermi surface. The low energy excitations of the fermions are then found about two nodes positioned on the $`k_y`$-axis at $`K_\pm =(0,\pm K)`$. We then linearize the theory about these nodes via $`ϵ_{K_\pm +q}=q_yv_F`$ and $`cc_1\mathrm{exp}(iK_+x)+c_2\mathrm{exp}(iK_{}x)`$. Introducing the spinor $`\psi ^{}=(c_1^{},c_2)`$, we can write $`H`$ as
$$H=d^2x\psi ^{}(iv_F\tau _z_y+iv_\mathrm{\Delta }\tau _x_x)\psi .$$
(4)
The Pauli matrices $`\tau _i`$ act in the particle-hole space of the spinors. $`H`$ is precisely the Hamiltonian of a single Dirac fermion in $`2+1`$ dimensions. To obtain the sigma model description of this Hamiltonian, we fix the energy, $`\omega `$, at which we work, and describe the theory in terms of an action of a two-dimensional Euclidean field theory:
$$S_0=d^2x\psi ^{}(iv_F\tau _z_y+iv_\mathrm{\Delta }\tau _x_xi\omega \tau _z)\psi .$$
(5)
With this action we can compute correlators of the form $`1/(H\omega )`$.
We include on-site disorder by adding a term,
$$H_{\mathrm{disorder}}=t(c_1^{}c_1+c_2^{}c_2),$$
(6)
to the Hamiltonian. Here $`t`$ is a random variable with variance $`t(x)t(y)=(2u)^1\delta (xy)`$. To compute correlators with disorder we replicate the theory, $`\psi \psi _k`$. Disorder is then easily averaged over leaving a set of quartic terms, $`S_{\mathrm{disorder}}=\frac{1}{u}(\psi _k^{}\tau ^z\psi _k)(\psi _l^{}\tau ^z\psi _l)`$. Reorganizing the fields via $`\stackrel{~}{\psi }^{}(\psi _R^{},\psi _L^{})=\psi ^{}\tau ^z\mathrm{exp}(i\pi \tau ^x/4)`$ and $`\stackrel{~}{\psi }(\psi _R,\psi _L)=\mathrm{exp}(i\pi \tau ^x/4)\psi `$, gives the action (at $`\omega =0`$) to be:
$$S=d^2xi\stackrel{~}{\psi }_k^{}(v_F_yv_\mathrm{\Delta }\tau ^z_x)\stackrel{~}{\psi }_k\frac{1}{u}(\stackrel{~}{\psi }_k^{}\stackrel{~}{\psi }_k)(\stackrel{~}{\psi }_l^{}\stackrel{~}{\psi }_l).$$
(7)
This theory is invariant under the group $`U(N)_L\times U(N)_R`$ transforming $`\stackrel{~}{\psi }\frac{1}{2}((1+\tau ^z)U_L+(1\tau ^z)U_R)\stackrel{~}{\psi }`$. Adding disorder to the Cooper-pairing term results in another four-fermion term, but preserves this symmetry. As long as the symmetry structure is unchanged, the low-energy sigma model should be the same.
When $`\omega =0`$, this action is equivalent to a massless Dirac fermion in a random magnetic field; the chiral symmetry prevents a mass term from appearing. This model can be solved by bosonization . Free fermions are equivalent to the WZW model $`SU(N)_1\times U(1)`$ . The four-fermion term is then simply expressed as $`(J_L+J_R)^2`$, where $`J_L`$ and $`J_R`$ are the $`U(1)`$ currents, $`J_L=(\psi _L^k\psi _L^k)`$ and $`J_R=(\psi _R^k\psi _R^k)`$. These $`U(1)`$ currents can be expressed in terms of a free boson with $`j_L=i_L\varphi `$ and $`j_R=i_R\varphi `$. The four-fermion coupling then merely determines the boson radius. Therefore, the model reduces to pure $`SU(N)_1`$ together with a decoupled free boson, and hence is critical.
To see directly how the WZW term appears, it is useful to rederive this result from the sigma model approach . We assume there is an energy scale where some fermion bilinear gets an expectation value. Formally speaking, we introduce a Hubbard-Stratonovich matrix field, $`M_{kl}`$, to factor the four-fermion term:
$$S=S_0d^2x\left(\frac{u}{4}\text{Tr}(M^2)+i(\stackrel{~}{\psi }_k^{}M_{kl}\stackrel{~}{\psi }_l)\right),$$
(8)
where $`M`$ is hermitian. Under $`U(N)_L\times U(N)_R`$, $`MUMU^{}`$. The WZW term appears in low-energy effective theories describing saddle points where $`M`$ is off-diagonal, e.g. $`M_{LL}=M_{RR}=0`$, but $`M_{LR}=M_{RL}^{}I`$, the identity. The diagonal subgroup $`U(N)_V`$ leaves this saddle point invariant, and so the low-energy modes $`T=M_{RL}`$ take values in $`U(N)_L\times U(N)_R/U(N)_VU(N)`$. Note that under parity, $`T^{}T`$. Focusing solely upon these modes allows us to write
$$S=S_0d^2x\left(\frac{u}{2}\text{Tr}(TT^{})+i(\stackrel{~}{\psi }_R^{}T\stackrel{~}{\psi }_L+\stackrel{~}{\psi }_L^{}T^{}\stackrel{~}{\psi }_R)\right).$$
(9)
Integrating out the fermions leaves an effective action for the bosonic field $`T`$, which can be expanded in powers of the momentum over the expectation value of $`T`$. At low energy, one can safely neglect four-derivative terms and higher, and one easily obtains the kinetic term (1) for $`T`$. However, a less-obvious term also appears. The case of interest was treated in . To perform a consistent low-energy expansion, one must change field variables. This results in a Jacobian in the path integral , which is precisely the WZW term. Writing $`T=he^{i\varphi }`$, where det($`h)=1`$, we obtain here the WZW term (2) of the $`SU(N)`$ field $`h`$, with level $`k=1`$. The extra $`U(1)`$ field $`\varphi `$ in the sigma model is a theory of a decoupled boson.
Thus in the sigma model approach we recover the critical line described by $`SU(N)_1\times U(1)`$ conformal field theory in the limit $`N0`$. This critical line is a part of the critical space of the Gade hopping model; the full space is larger because the $`SU(N)`$ coupling does not flow when $`N=0`$ . Moreover, this same CFT also applies to a time-reversal symmetric version of the hopping model, which is described by a $`U(N)/O(N)`$ sigma model. When $`\theta =\pi `$, this sigma model flows to the same $`SU(N)_1\times U(1)`$ conformal field theory . We can see how the two are related by perturbing the $`U(N)`$ WZW model action by a term which gives a mass to the modes outside a $`U(N)/O(N)`$ subspace (the subspace corresponds to symmetric $`U(N)`$ matrices). As with the $`N=2`$ case mentioned above, the low energy limit of the perturbed $`U(N)`$ WZW model reduces to the $`U(N)/O(N)`$ sigma model with $`\theta =k\pi `$. Thus the Gade model with and without time reversal invariance provides a concrete realization of the equivalence of theta and WZW terms in disordered systems.
In no sense does the WZW term arise in the above models as a result of fine tuning: it must appear when there is a chiral anomaly. In the model treated above, the fermions have a $`U(N)_L\times U(N)_R`$ symmetry. As is well known, chiral symmetries involving fermions are frequently anomalous, so that the Noether currents do not all remain conserved in the quantum theory. For massless fermions in $`1+1`$ dimensions, this was shown in detail in . The WZW term is the effect of the anomaly on the low-energy theory. Even though the fermions effectively become massive when $`T`$ gets an expectation value, their presence still has an effect on the low-energy theory, even if this mass is arbitrarily large. This violation of decoupling happens because the chiral anomaly must be present in the low-energy theory. In other words, the anomaly coefficient does not renormalize. This follows from an argument known as ’t Hooft anomaly matching . One imagines weakly gauging the anomalous symmetry (in our case the axial $`SU(N)`$). It is not possible to gauge an anomalous symmetry in a renormalizable theory, but one can add otherwise non-interacting massless chiral fermions to cancel the anomaly. Adding these spectator fermions ensures that the appropriate Ward identities are obeyed and the symmetry can be gauged. In the low-energy effective theory, the Ward identities must still be obeyed and the theory must remain anomaly-free. Because the massless spectator fermions are still present in the low-energy theory, there must be a term in the low-energy action which cancels the anomaly from the spectators. This is the WZW term.
We now turn to a case of considerable interest, that of $`d_{x^2y^2}`$ superconductors in the presence of disorder . We now study fermions with spin symmetry, so that $`c,c^{}`$ carry a node index, a spin index, and a replica index. We label the nodes $`(1,\pm )`$ for the pair near the wavevectors $`k=\pm (\pi /2,\pi /2)`$ and $`(2,\pm )`$ for the pair near $`k=\pm (\pi /2,\pi /2)`$. It is convenient to group particles and holes together. Thus define four doublets, $`\psi _{a\alpha }^j=(c_{a\alpha }^j,i(\sigma ^ys^xc^j)_{a\alpha })`$, $`a=\pm ,\alpha =/`$, for each pair of nodes $`j=1,2`$. The Pauli matrices $`\sigma ^a`$ and $`s^a`$ act upon the spin/$`\pm `$ node indices respectively. If we rotate the system in the $`ab`$ plane by $`45^o`$, the linearized Hamiltonian for a superconductor with time-reversal and spin symmetry is
$$H=i\overline{\psi }^1(v_Fs^z\tau ^z_x+v_\mathrm{\Delta }s^z\tau ^x_y)\psi ^1+(x,1y,2),$$
(10)
where the Pauli matrix $`\tau `$ acts in particle/hole space and $`\overline{\psi }^j(s^x\sigma ^y\tau ^y\psi ^j)^T`$. After adding disorder, the replicated Hamiltonian is invariant under the group $`Sp(2N)_L\times Sp(2N)_R`$ . It sends $`\psi ^iU\psi ^i`$, where $`U=\frac{1}{2}\left[U_L(1+\tau ^y)+U_R(1\tau ^y)\right],`$ with $`U_L`$ and $`U_R`$ each elements of $`Sp(2N)`$ acting only upon spin and replica indices. An element $`Q`$ of $`Sp(2N)`$ is an invertible $`2N\times 2N`$ real matrix obeying $`Q^T\sigma ^yQ=\sigma ^y`$. With Gaussian on-site disorder and $`\omega =0`$, the low-energy modes of the corresponding action take values in $`Sp(2N)_L\times Sp(2N)_R/Sp(2N)_V`$. The low-energy action thus includes the term (1) with $`h`$ in $`Sp(2N)`$ .
This model also has a WZW term. We first consider the case where disorder does not couple the pairs of nodes. Once disorder has been averaged over and the quartic terms have been factorized, the fermions interact with hermitian matrix fields, $`M^{(1)}`$ and $`M^{(2)}`$, via terms of the form
$$\overline{\psi }_\alpha ^1M_{\alpha \beta }^{(1)}\psi _\beta ^1+\overline{\psi }_\alpha ^2M_{\alpha \beta }^{(2)}\psi _\beta ^2,$$
(11)
as before. At a saddle point with non-zero values of the off-diagonal elements, the argument of shows that we will again obtain level 1 WZW theories for $`M^{(1)}`$ and for $`M^{(2)}`$, in this case $`Sp(2N)_1`$.
We make this explicit by showing how the spin degrees of freedom expand the previous $`U(N)`$ symmetry to that of $`Sp(2N)`$. $`Sp(2N)`$ has a $`U(N)`$ subgroup consisting of orthogonal real matrices of the form
$$\left(\begin{array}{cc}A& B\\ B& A\end{array}\right),$$
where we put the up spins in the first $`N`$ components of $`\psi `$, and the down spins in the second half. Thus the $`U(N)`$ subgroup acts on the combinations, $`\psi _{}\pm i\psi _{}`$, with the matrices $`AiB`$. In group-theoretic language, the $`2N`$ dimensional representation of $`Sp(2N)`$ decomposes into the $`N+\overline{N}`$ representation of $`U(N)`$. We can examine the WZW model describing the low-energy excitations which lie in the $`U(N)`$ subspace. It must be exactly the same as that treated above, with the exception that there are now two species of “spinless” fermions $`\psi _{}\pm i\psi _{}`$. This results in two WZW terms, one with $`h`$ and one with $`h^{}`$. As the WZW term is invariant under complex conjugation, the two contributions are identical, and so add. This theory is a $`U(N)_2`$ WZW theory (as was also derived using bosonization ; the corresponding result in the supersymmetric formulation was derived in ). The level of the WZW term in the full $`Sp(2N)`$ theory is determined by the embedding: when an $`H_l`$ theory is embedded in $`G_k`$, the levels obey $`kr=l`$, where $`r`$ is a group-theory factor called the index of the embedding of $`H`$ into $`G`$. For the embedding of $`U(N)`$ into $`Sp(2N)`$, $`r=2`$. Thus we have shown that the low-energy theory for (11) is given by an $`Sp(2N)_1`$ WZW model for each pair of nodes.
When a WZW term is present, the quasi-particles of the spin/time-reversal invariant superconductor need not be localized. Specifically, the coupling constant of the $`Sp(2N)_k`$ theory, inversely related to the spin conductance, now flows to a fixed value (provided its bare value is sufficiently small), as opposed to becoming arbitrarily large. Thus with a WZW term such models are no longer spin insulators. The $`Sp(2N)_k`$ WZW models seem to be well behaved as $`N0`$. The dimensions of the fundamental operators are given by $`m(2m)/(2k+2)`$, where $`m`$ is a positive integer. To compute the density of states at a finite energy, $`\omega `$, we add a term of the form $`\omega \mathrm{Tr}(h+h^{})`$ to the action. The density of the states is then given by $`\rho (\omega )\mathrm{Tr}(h+h^{})`$. This field $`\text{tr }h`$ has dimension $`\mathrm{\Delta }=1/(2k+2)`$, and a simple scaling argument yields $`\rho (\omega )\omega ^{\mathrm{\Delta }/(2\mathrm{\Delta })}=\omega ^{1/(4k+3)}`$. This result agrees with a previous computation when $`k=1`$.
When the two pairs of nodes are not coupled by the disorder, the low-energy theory is therefore described by two $`Sp(2N)_1`$ WZW models, which have symmetry currents, $`J_{L1,2},J_{R1,2}`$. Once the pairs of nodes are coupled, there are three possibilities. The first is that the two $`Sp(2N)_L\times Sp(2N)_R`$ symmetries are preserved, so the low-energy theory remains two $`Sp(2N)_1`$ WZW models. The second is that only simultaneous chiral transformations on the two pairs remain a symmetry. In the decoupled theory, this symmetry is generated by $`J_{L1}+J_{L2}`$ and $`J_{R1}+J_{R2}`$. It has level $`2`$ (in the language used above, the index of the embedding is $`2`$). Once the nodes are coupled, the generators may be deformed away from $`J_{L1}+J_{L2}`$ and $`J_{R1}+J_{R2}`$, but the level cannot change. This occurs for the same reasons that anomalies remain in the low-energy theory – the anomaly coefficient cannot renormalize. Thus the low-energy behavior of this model is described by an $`Sp(2N)_2`$ WZW model. If one modifies the Hamiltonian so that the two pairs of nodes are no longer identical, then the result is the same as long as the chiral symmetry is not explicitly broken: the chiral anomaly does not change under perturbation.
The third possibility is that the anomalies cancel between the two theories. This is in fact what happens in $`d_{x^2y^2}`$ superconductors . Consider two coupled WZW models with fields $`h_1`$ and $`h_2`$ and action
$$S_{PCM}(h_1)+\mathrm{\Gamma }(h_1)+S_{PCM}(h_2)+\mathrm{\Gamma }(h_2)+\lambda \text{Tr}(h_1h_2+h.c.).$$
The global symmetries of the decoupled model are $`h_iU_{Li}h_iU_{Ri}`$. In the coupled model this global symmetry is broken to $`h_1U_{L1}h_1U_{L2}^{}`$, $`h_2U_{L2}h_2U_{L1}^{}`$, appropriate to the underlying lattice symmetry of the d-wave superconductor. These symmetries are anomaly free; the axial anomalies of the two nodes cancel upon coupling. This can be seen roughly in that the coupling $`\lambda `$ is relevant; in the strong coupling limit $`h_1=h_2^{}`$ (i.e. $`\mathrm{\Gamma }(h_1)=\mathrm{\Gamma }(h_2)`$) and so the WZW terms cancel. This omits terms that might arise in integrating out massive degrees of freedom. This possibility cannot be ignored, because we know this is precisely how a WZW term arises in . However, this argument can be made more precise in the case when $`h_1`$ and $`h_2`$ are elements of $`O(N)`$ instead of $`Sp(2N)`$. Indeed, then one can fermionize one of the two WZW models, writing $`(h^2)_{ab}=\overline{\psi }_b\psi _a`$, with the resulting action
$$S_{PCM}(h_1)+\mathrm{\Gamma }(h_1)+\overline{\psi }\gamma ^\mu _\mu \psi +\lambda (\overline{\psi }h_1\psi +h.c.).$$
For $`\lambda `$ large, we can integrate out these fermions according to as before. In doing so, a WZW term for $`h_1`$ is induced which precisely cancels the original WZW term for $`h_1`$. If we had instead coupled the theories via a term $`\text{Tr}(h_1h_2^{}+h.c.)`$, the two terms would have added, and we would have recovered the level two theory discussed above. This cancellation occurs even if the two theories have differing Fermi velocities. Including Fermi velocities directly in the action (easily done through dimensional analysis) shows that the terms $`S_{PCM}`$ are affected while the WZW terms $`\mathrm{\Gamma }`$ are not. Thus the cancellation between the two WZW terms can proceed as above.
A d-wave superconductor with spin rotational but broken time reversal invariance is realized when the gap wave function is $`\mathrm{\Delta }d_{x^2y^2}+id_{xy}`$. This breaks the $`Sp(2N)\times Sp(2N)`$ symmetry explicitly, and the appropriate sigma model in the replica formulation for this theory is $`Sp(2N)/U(N)`$ with a theta term . Because the time-reversal invariance is broken, the coefficient of $`\theta `$ is not constrained by symmetry, but at $`\theta =\pi `$ one expects a critical point in the same fashion as . It is conjectured in that at $`\theta =\pi `$, this model flows to the $`Sp(2N)_1`$ conformal field theory. By utilizing the supergroup formulation of the corresponding network model (spin quantum Hall effect), it was shown that the critical point point of the disordered model is equivalent to classical percolation . Thus the replica limit $`N0`$ of $`Sp(2N)_1`$ should be equivalent to the critical point in the spin quantum Hall effect. The exponents computed above do indeed agree with those computed in , for example $`\rho (E)E^{1/7}`$. We again see a feature that appeared in the $`U(N)`$ class of problems: related models seemingly in two different universality classes are in the same one, once the WZW terms and $`\theta `$ terms are specified.
We have shown how stable critical points can arise in a number of replica sigma models. This shows that a number of disordered symmetry classes in two dimensions allow further refinement. In particular, in models with a chiral symmetry, one must compute the anomaly in the original underlying short-distance replica (or supergroup) model to determine the WZW term in the low-energy model. Likewise, in sigma models with $`𝐙_2`$ winding numbers, one must determine whether $`\theta =0`$ or $`\pi `$.
We would like to thank K. Intriligator, A. Ludwig, C. Nayak, T. Senthil, B. Simons and M. Zirnbauer for extremely helpful conversations. This work was supported by a DOE OJI Award, a Sloan Foundation Fellowship, and by NSF grant DMR-9802813. |
warning/0003/nucl-th0003038.html | ar5iv | text | # I Introduction
## I Introduction
Collisions at relativistic heavy ion colliders like the Relativistic Heavy Ion Collider RHIC/Brookhaven and the Large Hadron Collider LHC/CERN (operating in its heavy ion mode) are mainly devoted to the search of the Quark Gluon Plasma. However, peripheral heavy ion collisions also open up a broad area of studies as advocated by Baur and collaborators . Examples are the possible discovery of an intermediate-mass Higgs boson or beyond standard model physics using peripheral ion collisions, which have been discussed at length in the literature. More promissing than these may be the study of hadronic physics, which will appear quite similarly to the two-photon hadronic physics at $`e^+e^{}`$ machines with the advantage of a huge photon luminosity peaked at small energies . Due to this large photon luminosity it will become possible to discover resonances that couple weakly to the photons .
Double-pomeron exchange will also occurs in peripheral heavy ion collisions and their contribution is similar to the two-photons one as discussed by Baur and Klein . A detailed calculation performed by Müller and Schramm of Higgs boson production have shown that the diffractive contribution is much smaller than the electromagnetic one . We can easily understand this result remembering that the coupling between the Higgs boson and the pomerons is intermediated by quarks, and according to the pomeron model of Donnachie and Landshoff when in the vertex pomeron-quark-quark any of the quark legs goes far “off-shell” the coupling with the pomeron decreases. Therefore, we do not need to worry about the pomeron-pomeron contribution in peripheral heavy ion collisions when heavy (or far “off-shell”) quarks are present. However, this is not what happens in the case of light resonances , where double diffraction were claimed to be as important as photon initiated processes. In particular, Engel et al. have shown that at the LHC the diffractive production of hadrons may be a background for the photonic one.
In Ref. it was remarked that the effect of removing “central collisions” should also be performed in the double-pomeron calculation, implying in a considerable reduction of the background calculated in Ref.. This claim is the same presented by Baur and Cahn and Jackson in the case of early calculations of peripheral heavy ion collisions. Roughly speaking we must enforce that the minimum impact parameter ($`b_{min}`$) should be larger than ($`R_1+R_2`$), where $`R_i`$ is the nuclear radius of the ion “$`i`$”, in order to have both ions coming out intact after the interaction.
In this work we will compute the production of resonances, pion-pairs and a hadron cluster with invariant mass $`M_X`$ through photon-photon and pomeron-pomeron fusion in peripheral heavy ion collisions at the energies of RHIC and LHC. We will take into account the effect of the impact parameter as discussed in the previous paragraph for photons as well as for pomerons. We also compare this approach to cut the central collisions with the use of an absorption factor in the Glauber approximation. The inclusion of pion-pairs production is important because they certainly will be studied at these colliders, and they also represent a background for glueball (and other hadrons) detection. The pomeron physics within the ion will be described by the Donnachie and Landshoff model . We will focus on the values of the cross sections that shall be measured in the already quoted ion colliders, and point out when pomeron-pomeron processes can be considered competitive or not with photon-photon collisions. The arrangement of our paper is the following: Section 2 contains a discussion of the photons and pomerons distributions in the nuclei. In Sect. 3 we introduce the cross section for the elementary processes. Finally, Sect. 4 contains the results and conclusions.
## II Photons and pomerons distribution functions
### A Photons in the nuclei
The photon distribution in the nucleus can be described using the equivalent-photon or Weizsäcker-Williams approximation in the impact parameter space. Denoting by $`F(x)dx`$ the number of photons carrying a fraction between $`x`$ and $`x+dx`$ of the total momentum of a nucleus of charge $`Ze`$, we can define the two-photon luminosity through
$`{\displaystyle \frac{dL}{d\tau }}={\displaystyle _\tau ^1}{\displaystyle \frac{dx}{x}}F(x)F(\tau /x),`$ (1)
where $`\tau =\widehat{s}/s`$, $`\widehat{s}`$ is the square of the center of mass (c.m.s.) system energy of the two photons and $`s`$ of the ion-ion system. The total cross section $`ZZZZ\gamma \gamma ZZX`$, where $`X`$ is the particle produced within the rapidity gap, is
$`\sigma (s)={\displaystyle 𝑑\tau \frac{dL}{d\tau }\widehat{\sigma }(\widehat{s})},`$ (2)
where $`\widehat{\sigma }(\widehat{s})`$ is the cross-section of the subprocess $`\gamma \gamma X`$.
There remains only to determine $`F(x)`$. In the literature there are several approaches for doing so, and we choose the conservative and more realistic photon distribution of Ref.. Cahn and Jackson , using a prescription proposed by Baur , obtained a photon distribution which is not factorizable. However, they were able to give a fit for the differential luminosity which is quite useful in practical calculations:
$$\frac{dL}{d\tau }=\left(\frac{Z^2\alpha }{\pi }\right)^2\frac{16}{3\tau }\xi (z),$$
(3)
where $`z=2MR\sqrt{\tau }`$, $`M`$ is the nucleus mass, $`R`$ its radius and $`\xi (z)`$ is given by
$$\xi (z)=\underset{i=1}{\overset{3}{}}A_ie^{b_iz},$$
(4)
which is a fit resulting from the numerical integration of the photon distribution, accurate to $`2\%`$ or better for $`0.05<z<5.0`$, and where $`A_1=1.909`$, $`A_2=12.35`$, $`A_3=46.28`$, $`b_1=2.566`$, $`b_2=4.948`$, and $`b_3=15.21`$. For $`z<0.05`$ we use the expression (see Ref. )
$$\frac{dL}{d\tau }=\left(\frac{Z^2\alpha }{\pi }\right)^2\frac{16}{3\tau }\left(\mathrm{ln}(\frac{1.234}{z})\right)^3.$$
(5)
The condition for realistic peripheral collisions ($`b_{min}>R_1+R_2`$) is present in the photon distributions showed above, and the applications of Sect. 4 are straightforward once we determine the cross sections for the elementary processes.
### B Pomerons in the nuclei
In the case where the intermediary particles exchanged in the nucleus-nucleus collisions are pomerons instead of photons, we can follow closely the work of Müller and Schramm and make a generalization of the equivalent photon approximation method to this new situation. So the cross section for particle production via two pomerons exchange can be written as
$`\sigma _{AA}^{PP}={\displaystyle 𝑑x_1𝑑x_2f_P(x_1)f_P(x_2)\sigma _{PP}(s_{PP})},`$ (6)
where $`f_P(x)`$ is the distribution function that describe the probability for finding a pomeron in the nucleus with energy fraction $`x`$ and $`\sigma _{PP}(s_{PP})`$ is the subprocess cross section with energy squared $`s_{PP}`$. In the case of inclusive particle production we use the form given by Donnachie and Landshoff
$`f_P(x)={\displaystyle \frac{1}{4\pi ^2x}}{\displaystyle _{\mathrm{}}^{(xM)^2}}𝑑t|\beta _{AP}(t)|^2|D_P(t;s^{})|^2,`$ (7)
where $`D_P(t;s^{})`$ is the pomeron propagator
$`D_P(t;s)={\displaystyle \frac{(s/m^2)^{\alpha _P(t)1}}{\mathrm{sin}(\frac{1}{2}\pi \alpha _P(t))}}\mathrm{exp}\left({\displaystyle \frac{1}{2}}i\pi \alpha _P(t)\right),`$ (8)
with $`s`$ the total squared c.m. energy. The Regge trajectory obeyed by the pomeron is $`\alpha _P(t)=1+\epsilon +\alpha _P^{}t`$, where $`\epsilon =0.085`$, $`\alpha _P^{}=0.25`$ GeV<sup>-2</sup> and $`t`$ is a small exchanged four-momentum square, $`t=k^2<<1`$, so the pomeron behaves like a spin-one boson. The term in the denominator of the pomeron propagator, $`[\mathrm{sin}(\frac{1}{2}\pi \alpha _P(t))]^1`$, is the signature factor that express the different properties of the pomeron under C and P conjugation. At very high c.m. energy this factor falls very rapidly with $`𝐤^2`$, whose exponential slope is given by $`\alpha _P^{}\mathrm{ln}(s/m^2)`$, $`m`$ is the proton mass, and it is possible to neglect this $`𝐤^2`$ dependence,
$`\mathrm{sin}{\displaystyle \frac{1}{2}}\pi (1+\epsilon \alpha _P^{}𝐤^2)\mathrm{cos}({\displaystyle \frac{1}{2}}\pi \epsilon )1.`$ (9)
If we define the pomeron range parameter $`r_0`$ as
$`r_0^2=\alpha _P^{}\mathrm{ln}(s/m^2),`$ (10)
the pomeron propagator can be written as
$`|D_P(t=𝐤^2;s)|=(s/m^2)^\epsilon e^{r_0^2𝐤^2}.`$ (11)
Since we are interested in the spatial distribution of the virtual quanta in the nuclear rest frame we are using $`t=𝐤^2`$.
The nucleus-pomeron coupling has the form
$`\beta _{AP}(t)=3A\beta _0F_A(t),`$ (12)
where $`\beta _0=1.8`$ GeV<sup>-1</sup> is the pomeron-quark coupling, $`A`$ is the atomic number of the colliding nucleus, and $`F_A(t)`$ is the nuclear form factor for which is usually assumed a Gaussian expression (see, e.g., Drees et al. in )
$`F_A(t)=e^{t/2Q_0^2},`$ (13)
where $`Q_0=60`$ MeV.
Performing the $`t`$ integration of the distribution function in Eq.(7) we obtain
$`f_P(x)`$ $`=`$ $`{\displaystyle \frac{(3A\beta _0)^2}{(2\pi )^2x}}\left({\displaystyle \frac{s^{}}{m^2}}\right)^{2\epsilon }{\displaystyle _{\mathrm{}}^{(xM)^2}}𝑑te^{t/Q_0^2}`$ (14)
$`=`$ $`{\displaystyle \frac{(3A\beta _0Q_0)^2}{(2\pi )^2x}}\left({\displaystyle \frac{s^{}}{m^2}}\right)^{2\epsilon }\mathrm{exp}\left[\left({\displaystyle \frac{xM}{Q_0}}\right)^2\right].`$ (15)
The total cross section for a inclusive particle production is obtained with the above distribution and also with the expression for the subprocess $`PPX`$ as prescribed in Eq.(6). However, in Eq.(6) the cases where the two nuclei overlap are not excluded. To enforce the realistic condition of a peripheral collision it is necessary to perform the calculation taking into account the impact parameter dependence, $`b`$. It is straightforward to verify that in the collision of two identical nuclei the total cross section of Eq.(6) is modified to
$`{\displaystyle \frac{d^2\sigma _{AA}^{PPX}}{d^2b}}={\displaystyle \frac{Q^2}{2\pi }}e^{Q^2b^2/2}\sigma _{AA}^{PP},`$ (16)
where $`(Q^{})^2=(Q_0)^2+2r_0^2`$. The total cross section for inclusive processes is obtained after integration of Eq.(16) with the condition $`b_{min}>2R`$ in the case of identical ions.
For exclusive particle production the determination of the pomeron distribution function in the nuclei is slightly modified, because in this case it is necessary some specific assumption about the pomeron internal structure . Following Ref. the distribution function of pomerons is
$`f_P(x)={\displaystyle \frac{(3A\beta _0)^2}{(2\pi )^2x}}{\displaystyle _{\mathrm{}}^{(xM)^2}}𝑑t(tx^2M^2)F_A(t)^2|D(t)|^2,`$ (17)
and the cross section for a resonance production as a function of the impact parameter is
$`{\displaystyle \frac{d^2\sigma _{AA}^{PPR}}{d^2b}}`$ $`=`$ $`2\pi \left({\displaystyle \frac{3A\beta _0}{2\pi ^2}}\right)^4{\displaystyle \frac{dx_1}{x_1}\frac{dx_2}{x_2}Q_1^4Q_2^4\stackrel{~}{Q}^2e^{x_1^2M^2/Q_1^2}e^{x_2^2M^2/Q_2^2}}`$ (18)
$`\times `$ $`\left({\displaystyle \frac{x_1x_2s^2}{m^4}}\right)^{2\epsilon }\sigma _{AA}^{PPR}(x_1x_2s)b^2\stackrel{~}{Q}^2e^{b^2\stackrel{~}{Q}^2/2},`$ (19)
with $`\sigma _{AA}^{PPR}(x_1x_2s)`$ indicating the subprocess cross section (double pomeron fusion producing a resonance), and where
$`\stackrel{~}{Q}^2={\displaystyle \frac{1}{2}}(Q_1^2+Q_2^2),`$ (20)
with $`Q_i^2Q_0^2+2r_0^2`$ for idential ions. In the calculations we are going to perform we noticed that the aproximation $`Q_i^2Q_0^2`$ is quite reasonable, because for the energies that we shall consider the pomeron range parameter (Eq.(10)) is smaller than the width of the Gaussian form factor and consequently $`\stackrel{~}{Q}^2Q_0^2`$. Therefore, we obtain the final expression
$`{\displaystyle \frac{d^2\sigma _{AA}^{PPR}}{d^2b}}`$ $`=`$ $`2\pi \left({\displaystyle \frac{3A\beta _0Q_0^2}{2\pi ^2}}\right)^4{\displaystyle \frac{dx_1}{x_1}\frac{dx_2}{x_2}e^{x_1^2M^2/Q_0^2}e^{x_2^2M^2/Q_0^2}}`$ (21)
$`\times `$ $`\left({\displaystyle \frac{x_1x_2s^2}{m^4}}\right)^{2\epsilon }\sigma _{AA}^{PPR}(x_1x_2s)b^2Q_0^4e^{b^2Q_0^2/2}.`$ (22)
As discussed previously, to enforce the condition of peripheral collisions we integrate Eq.(22) with the condition $`b_{min}>2R`$.
Another way of to exclude events due to inelastic central collisions is through the introduction of an absortion factor computed in the Glauber aproximation . This factor modifies the cross section in the following form
$`{\displaystyle \frac{d\sigma _{AA}^{gl}}{d^2b}}`$ $`=`$ $`{\displaystyle \frac{d\sigma _{AA}^{PPR}}{d^2b}}\mathrm{exp}\left[A^2b\sigma _0{\displaystyle \frac{dQ^2}{(2\pi )^2}F_A^2(Q^2)e^{iQb}}\right]`$ (23)
$`=`$ $`{\displaystyle \frac{d\sigma _{AA}^{PPR}}{d^2b}}\mathrm{exp}\left[A^2b\sigma _0{\displaystyle \frac{Q_0^2}{4\pi }}e^{Q_0^2b^2/4}\right],`$ (24)
where $`\sigma _0`$ is the nucleon-nucleon total cross section, whose value for the different energy domains that we shall consider is obtained directly from the fit of Ref.
$`\sigma _0=Xs^ϵ+Y_1s^{\eta _1}+Y_2s^{\eta _2},`$ (25)
with $`X=18.256`$, $`Y_1=60.19`$, $`Y_2=33.43`$, $`ϵ=0.34`$, $`\eta _1=0.34`$, $`\eta _2=0.55`$, $`F_A(Q^2)=e^{Q^2/2Q_0^2}`$ and we exemplified Eq.(24) for the case of resonance production, i.e., $`\sigma _{AA}^{PPR}`$ is the total cross section for the resonance production to be discussed in the next section. The integration in Eq.(24) is over all impact parameter space and in the last section we discuss the differences between the two approaches showed above for removing central collisions.
## III Subprocesses initiated by photons and pomerons
### A Resonances
The main motivation to study resonance production in peripheral heavy ion collisions is that the high photon luminosity will allow us to observe resonances that couple very weakly to photons. The simplicity of this calculation also enable us to test the methods for removing central collisions, as well as to check up to which degree the double pomeron exchange is or not a background for the two photon physics.
To estimate the production of single spin-zero resonances, we note that these states can be formed by photon-photon fusion with a coupling strength that is measured by their photonic width
$`\widehat{\sigma }_{\gamma \gamma R}={\displaystyle \frac{8\pi ^2}{M_R\widehat{s}}}\mathrm{\Gamma }_{R\gamma \gamma }\delta \left(\tau {\displaystyle \frac{M_R^2}{\widehat{s}}}\right),`$ (26)
where $`M_R`$ is the ressonance mass and $`\mathrm{\Gamma }_{R\gamma \gamma }`$ its decay width in two photons. Using this expression into Eq.(2) we obtain the total cross section for the production of pseudo-scalar mesons.
To compute the cross section of the subprocess $`PPR`$ we can use the pomeron model of Donnachie and Landshoff . In this model it is assumed that the pomeron couples to the quarks like a isoscalar photon . This means that the cross sections of $`PPX`$ subprocesses can be obtained from suitable modifications on the cross-section for $`\gamma \gamma X`$. Another aspect to be considered is that the pomeron-quark-quark vertex is not point-like, and when either or both of the two quark legs in this vertex goes far off shell the coupling is known to decrease. So the quark-pomeron coupling $`\beta _0`$ must be replaced by
$`\stackrel{~}{\beta }_0(q^2)={\displaystyle \frac{\beta _0\mu _0^2}{\mu _0^2+Q^2}},`$ (27)
where $`\mu _0^2=1.2`$ GeV<sup>2</sup> is a mass scale characteristic of the pomeron, in the case of resonance production $`Q=M_R/2`$ measures how far one of the quark legs is off mass shell and $`M_R`$ is the resonance mass. Therefore, the process $`PPR`$ is totally similar to the one initiated by photons unless from an appropriate change of factors. The cross section we are looking for is obtained changing the fine-structure constant $`\alpha `$ by $`9\stackrel{~}{\beta }/16\pi ^2`$, where $`\stackrel{~}{\beta }`$ is giving by Eq.(27) and $`9=3^2`$ is a color factor, leading to
$`\sigma _{PP}^R={\displaystyle \frac{9}{2}}{\displaystyle \frac{\stackrel{~}{\beta }^4}{\alpha ^2}}{\displaystyle \frac{\mathrm{\Gamma }(R\gamma \gamma )}{M_R}}\delta (x_1x_2sM_R^2).`$ (28)
Using this expression in Eq.(22) the total cross section is equal to
$`\sigma _{AA}^{PPR}`$ $`=`$ $`{\displaystyle \frac{9\pi }{8}}{\displaystyle \frac{(\stackrel{~}{\beta }Q_0)^4}{\alpha ^2}}\left({\displaystyle \frac{3A\beta _0Q_0}{2\pi }}\right)^4{\displaystyle \frac{\mathrm{\Gamma }(R\gamma \gamma )}{M_R}}\left({\displaystyle \frac{M_R^2s}{m^4}}\right)^{2ϵ}{\displaystyle \frac{Q_0^4}{M_R^2}}`$ (29)
$`\times `$ $`{\displaystyle \frac{dx}{x}\mathrm{exp}\left[\left(\frac{M_R^2M}{sQ_0x}\right)^2\frac{(xM)^2}{Q_0^2}\right]_{b_{min}}^{\mathrm{}}𝑑b2\pi b^3e^{Q_0^2b^2/2}},`$ (30)
where $`b_{min}=2R`$.
### B Pion pair production
The continuous production of pion pairs ($`\pi ^+\pi ^{}`$) is also an interesting signal to be observed in peripheral heavy ion collisions, mostly because they are a background for glueball and other resonances decays. Here we discuss the subprocess cross sections for two photons, $`ZZ\gamma \gamma ZZ\pi ^+\pi ^{}`$, and two pomerons exchange $`ZZPPZZ\pi ^+\pi ^{}`$.
The cross section for pion pair production by two photons can be calculated approximately by using a low energy theorem derived from partially-conserved-axial-vector-current hypothesis and current algebra and is equal to
$`\sigma (\gamma \gamma \pi ^+\pi ^{}){\displaystyle \frac{2\pi \alpha ^2}{s}}\left(1{\displaystyle \frac{4m_\pi ^2}{s}}\right)^{(1/2)}\left[{\displaystyle \frac{m_V^4}{\left(\frac{1}{2}s+m_V^2\right)\left(\frac{1}{4}s+m_V^2\right)}}\right]^2,`$ (32)
where $`m_\pi `$ is the pion mass an $`s`$ its squared energy, $`m_V`$ is a free parameter, whose value that provides the best fit to the experimental data is $`m_V1.4`$ GeV. This expression shows a nice agreement with the experimental data . For large values of $`s`$ it deviates from the Brodsky and Lepage formula . However, since most of the photon distribution is concentrated in the small $`x`$ region, i.e., the photons carry a small fraction of the momentum of the incoming ion, the difference is negligible.
Using Eqs. (32) and (2) we obtain
$`\sigma (s)={\displaystyle \frac{2\pi \alpha ^2}{s}}{\displaystyle _{\tau _{min}}^1}{\displaystyle \frac{d\tau }{\tau }}\left(1{\displaystyle \frac{4m_\pi ^2}{s\tau }}\right)^{(1/2)}\left[{\displaystyle \frac{m_V^4}{\left(\frac{1}{2}s\tau +m_V^2\right)\left(\frac{1}{4}s\tau +m_V^2\right)}}\right]^2{\displaystyle \frac{dL}{d\tau }}.`$ (33)
In the case of double pomeron exchange producing a pion pair we use once again the Donnachie and Landshoff model for the pomeron, obtaining the cross section for $`PP\pi ^+\pi ^{}`$ from the photonic one changing $`\alpha ^29\stackrel{~}{\beta }_0^4/16\pi ^2`$ in $`\sigma (\gamma \gamma \pi ^+\pi ^{})`$, and the resulting expression replaces $`\sigma _{AA}^{PPR}(x_1x_2s)`$ in Eq.(22). The total cross section appears after we perform the integration in the parameter space representation of the following equation
$`{\displaystyle \frac{d^2\sigma _{AA}^{PP\pi ^+\pi ^{}}}{db^2}}`$ $`=`$ $`\left({\displaystyle \frac{\pi ^2}{8}}\right){\displaystyle \frac{9}{4}}{\displaystyle \frac{(\stackrel{~}{\beta }_0Q_0)^4}{s}}\left({\displaystyle \frac{3A\beta _0Q_0}{2\pi ^2}}\right)^4{\displaystyle \frac{dx_1}{x_1^2}\frac{dx_2}{x_2^2}e^{(x_1M)^2/Q_0^2}e^{(x_2M)^2/Q_0^2}}`$ (34)
$`\times `$ $`\left({\displaystyle \frac{x_1x_2s^2}{m^4}}\right)^{2\epsilon }\left(1{\displaystyle \frac{4m_\pi ^2}{x_1x_2s}}\right)^{1/2}\left[{\displaystyle \frac{m_V^4}{\left(\frac{x_1x_2s}{2}+m_V^2\right)\left(\frac{x_1x_2s}{4}+m_V^2\right)}}\right]^2`$ (35)
$`\times `$ $`{\displaystyle _{b_{min}}^{\mathrm{}}}𝑑b2\pi Q_0^4b^3e^{b^2Q_0^2/2}.`$ (36)
### C Multiple particle production
The elementary cross section for multiple-particle production via two photons fusion can be described by the parametrization
$`\sigma _{\gamma \gamma hadrons}=C_1\left({\displaystyle \frac{s}{s_0}}\right)^ϵ+C_2\left({\displaystyle \frac{s}{s_0}}\right)^\eta ,`$ (37)
where $`C_1=173`$ nbarn, $`C_2=519`$ nbarn, $`s_0=1`$ GeV<sup>2</sup>, $`ϵ=0.079`$ and $`\eta =0.4678`$. The total cross section comes out from Eq.(2).
Within the Donnachie and Landshoff model it is straightforward to see that with the above parametrization the differential cross section to produce a cluster of particles with mass $`M_X`$ through double pomeron exchange is
$`{\displaystyle \frac{d\sigma }{dM_X}}`$ $`=`$ $`{\displaystyle \frac{(3A\beta _0\stackrel{~}{\beta }_0\mu _0)^4}{(2\pi )^4R_N^4(16\pi ^2\alpha ^2)}}{\displaystyle \frac{1}{2M_X}}{\displaystyle \frac{ds^{}}{s^{}}\left[C_1\left(\frac{s^{}}{s_0}\right)^ϵ+C_2\left(\frac{s^{}}{s_0}\right)^\eta \right]}`$ (38)
$`\times `$ $`\mathrm{exp}\left[\left({\displaystyle \frac{s^{}MR_N}{s}}\right)^2\left({\displaystyle \frac{M_X^2MR_N}{s^{}}}\right)^2\right]{\displaystyle _{b_{min}}^{\mathrm{}}}𝑑bb{\displaystyle \frac{e^{b^2/2R_N^2}}{R_N^2}},`$ (39)
to obtain this expression we used the pomeron distribution function in the nucleus for inclusive process (Eq.(7)).
To be possible a comparison with the work of Engel et al. , we also make use of the Ter-Martirosyan model for diffractive multiparticle production. In this model the subprocess $`PPhadrons`$ is characterized by the cross section
$`\sigma _{PP}^{tot}(\mathrm{ln}(M_X^2/m^2),t_1,t_2)8\pi r(t_1)r(t_2),`$ (40)
which is a function of the triple-pomeron vertex $`r(t)`$, where t is the exchanged momentum. Using the value of $`r(0)`$ from Ref. , $`\sigma _{𝒫𝒫}^{tot}=8\pi r^2(0)140`$ $`\mu `$ barn. Note that we have clear differences between the approaches described above. Eq.(37) is a parametrization valid for a wide range of momenta, and with this one we naively apply the model of Ref. to compute the total cross section for multiparticle production. On the other hand Eq.(40) is obtained in another specific model and it is not expected to be valid for the same range of energies as Eq.(37). This difference is going to be discussed in the last section.
Streng applied the model of Ref. for proton-proton collisions where the initial protons are scattered almost elastically, emerging with a very large fraction of the initial energy,
$`|x_1|,|x_2|c,c0.9.`$ (41)
The double pomeron exchange produce a particle cluster within a large rapidity gap and with a mass of the order
$`M_X^2s(1|x_1|)(1|x_2|),`$ (42)
where $`s`$ is the reaction energy squared. As the scattering is almost elastic, i.e., the emerging beam has approximately the same energy as the incident one, the following kinematical boundaries can be introduced
$`M_0M_X(1c)\sqrt{s},`$ (43)
$`{\displaystyle \frac{M_X^2}{(1c)}}s_1(1c)s,`$ (44)
where $`M_0=2`$ GeV and $`c=0.9`$. These limits have been translated for the case of heavy ions by Engel et al. and we will proceed as them. If we consider Eq.(40), dress it with the pomeron distribution functions within the nuclei, and subtract the central collisions considering the absortion factor computed in the Glauber aproximation we reproduce the results of Ref. .
## IV Results and Conclusions
Peripheral collisions at relativistic heavy ion colliders provide an arena for interesting studies of hadronic physics. Resonances coupling weakly to photons can be studied due to the large photon luminosity, the continuous production of pion pairs will be observed not only as a reaction of interest as well as a possible background for some resonance decays. A hadron cluster produced within a large rapidity gap will give information about photon-photon and double pomeron exchange. In this work we estimate the cross section for these processes. One of the main points is to verify if the double pomeron exchange is or not a background for the purely electromagnetic process. We discussed double pomeron exchange according to the Donnachie and Landshoff model and calculated the cross sections in the impact parameter space. The condition for a realistic peripheral collision is imposed integrating the cross section with $`b_{min}>2R`$ in the case of two identical ions with radius $`R`$.
We considered the production of pseudoscalars resonances in the collision of <sup>238</sup>U for energies available at RHIC ($`\sqrt{s}=200`$ GeV/nucleon), and collisions of <sup>206</sup>Pb at energies available in LHC ($`\sqrt{s}=6.300`$ GeV/nucleon). Our results are shown in Table 1. Contrarily to the result of Ref. the double pomeron exchange is not important when the cut in the impact parameter is introduced. For a realistic peripheral collision in the case of resonance production the pomeron-pomeron process is at least two orders of magnitude below the photon-photon one. Note in Table 1 that the rate of diffractive resonance production decreases with the increase of the meson mass. The main reason for this behavior lies in the fast decrease of the pomeron-quark coupling as shown in Eq.(27).
Note that the results of this table assume $`100\%`$ of efficiency in tagging the peripheral collision, even if we consider a small efficiency we recall that the cross section for light resonances imply in approximately billions of events/yr which easily survive the cuts for the background separation proposed by Nystrand and Klein (see the last paper of Ref. ). One of the most important cuts to separate inelastic nuclear reactions, associated with grazing collisions, is the small multiplicity of the final state, and this is exactly what we may expect in the final state of the particles discussed in Table 1.
The decays of $`\pi ^0`$, $`\eta `$, etc… will be dominated by two (or three) body decays in the central region of rapidity, and easily separated from the larger multiplicity common to inelastic collisions. It is interesting that in the case of $`\pi ^0`$ and $`\eta `$ production we may focus on the $`2\gamma `$ decay, and even if it is possible to separate the background from inelastic nuclear reactions, we still have the background of the photon-photon scattering through the QED box diagram producing the same final state. The box diagram will be dominated by light quarks, electron and muon, and for these we can use the asymptotic expression of $`\gamma \gamma `$ scattering ($`\sigma (s)20/s`$). Integrating this expression in a bin of energy (proportional to the resonance partial width into two-photons) centered at the mass of the resonance, we obtain a cross-section smaller than the resonant one with subsequent decay into two-photons. We do not considered the interference between the box and resonant diagram because on resonance the two processes are out of phase. It is opportune to mention that the decay products will fill the central region of rapidity, which is also one of the conditions proposed in Ref. to isolate the peripheral collisions.
As discussed in Sect. 2 we have two different ways to enforce a realistic peripheral heavy ion collision. One is a geometrical cut in the impact parameter space where $`b_{min}>2R`$ is imposed, the other is through the introduction of the absorption factor in the Glauber approximation as given by Eq.(24). In Table 2 we compare the ratios between the total cross sections for diffractive resonance production computed with Eq.(24) and the one with the cut on the impact parameter (given by Eq.(LABEL:sigmaex)), in the collision of <sup>238</sup>U for energies available at RHIC ($`\sqrt{s}=200`$ GeV/nucleon), and collisions of <sup>206</sup>Pb at energies available in LHC ($`\sqrt{s}=6.300`$ GeV/nucleon).
The values of Table 2 show that the geometrical cut is less restrictive than the one given by the Glauber absorption factor. However, which one is more realistic also depends on the energy and on the ion that we are considering. In Table 3 we present the cross section for $`\pi ^0`$ production for different ions and at different energies. From Eq.(24) we notice that small variations in $`\sigma _0`$ (the nucleon-nucleon total cross section) are also promptly transmitted to the total cross section, and modify the ratios between the different methods to exclude inelastic collisions. Table 3 shows that the difference between the methods also become less important for light ions, but the most striking fact in this table is that for light ions double pomeron exchange starts becoming a background for photon-photon processes! According to Table 3 for <sup>28</sup>Si the diffractive $`\pi ^0`$ production is a factor of 2 down the electromagnetic one (assuming the geometrical cut). This is not surprising because we know that for proton-proton the double pomeron exchange process should be larger than the electromagnetic one for producing a light resonance.
In Table 4 we show the pion pair cross section for different ions. The values were obtained using the geometrical cut, and even if with this procedure the diffractive cross section is a little bit overestimated for heavy ions we verify that photon-photon dominates. For light ions the diffractive process is already of the order of $`10\%`$ of the electromagnetic one.
The simulations discussed in the last paper of Ref. have shown that the $`\gamma \gamma `$ interactions produce final states with small summed transverse momentum $`(|\overline{p}_T|)`$. Therefore, a cut of $`|\overline{p}_T|40\mathrm{\hspace{0.17em}\hspace{0.17em}100}MeV/c`$ can reduce considerably the background of non-peripheral collisions. In Table 4 we present the cross section for pion pair production through double photon interaction with $`|\overline{p}_T|100MeV/c`$. With this cut the cross section was reduced almost by a factor of $`4`$. The electromagnetic process with the restriction on $`p_T`$ is still larger than those of double pomeron exchange without this cut, and the introduction of this cut in the diffractive process produces a similar reduction.
The results for a hadron cluster production with invariant mass $`M_X`$ is depicted in Fig.1. In the figure it is shown the cross section for four different ions (Pb, Au, Ag, Ca) at energies that will be available at RHIC and LHC. The results were obtained integrating the cross sections with the condition $`b_{min}>2R`$. At LHC the photon-photon process will dominate the cross sections for heavy ions, whereas for light ions and small invariant mass they become of the same order. For heavy ions the diffractive process is indeed negligible. Note that our result for photons is similar to the one of Engel et al. , but the diffractive cross sections is slightly smaller than the one of Ref. . We credit this deviation to the differences in our approachs to calculate the subprocess cross section, mainly in the use of Eq.(37) with the changes prescribed by the Donnachie and Landshoff instead of Eq.(40) given by the Ter-Martirosyan model. They also use a value for $`\sigma _0`$ that is smaller than the one we considered here, which gives a smaller cut of the central collisions. We believe that the use of Eq.(37) and the model of Ref. is more appropriate for the full range of momenta. Actually, diffractive models are plagued by uncertainties and the measurement of the double pomeron exchange in heavy ion colliders will provide useful information to distinguish between different models.
For multiple particle production we will not have the criteria of low multiplicity to help us to select the truly peripheral collisions, as well as it is far from clear if the cut in transverse momentum will be very effective to select the $`\gamma \gamma `$ events. However, we can separate the peripheral events on the basis of a clustering in the central region of rapidity, although an extensive and detailed simulation of the background processes will be necessary in order to set the precise interval of rapidity needed to cut the inelastic nuclear collisions.
As verified by Drees, Ellis and Zeppenfeld , Eq.(13) is a reasonable approximation for the form factor obtained from a Fermi or Woods-Saxon density distribution. However, their result shows that for heavy final states the photon-photon luminosity is slightly underestimated, and we can expect the same for the Pomeron one. A simple form factor expression consistent with the Fermi distribution has been recently obtained in Ref. , and its use would yield a few percent larger cross section in the case of a very heavy hadron cluster production.
In the case of peripheral heavy ion collisions at RHIC we surely cannot neglect the diffractive contribution, for light ions and a hadron cluster with low invariant mass it surely dominates photon-photon collisions. Notice that these results may change if we use the Glauber absorption factor to compute the cross section (depending on the energy, the ion and invariant mass), but the actual fact is that double pomeron exchange cannot be neglected at RHIC.
In conclusion, we estimated the production of resonances, pion pairs and a cluster of hadrons with invariant mass $`M_X`$ in peripheral heavy ion collisions at energies that will be available at RHIC and LHC. The condition for a realistic peripheral collision was studied with the use of a geometrical cut, where the minimum impact parameter was forced to be larger than twice the identical nuclei radius. The introduction of an absorption factor in the Glauber approximation to eliminate central collisions was also studied. We find out that the most restrictive method to account for inelastic collisions depended on the energy, the ion, as well as on the value of $`\sigma _0`$ (the nucleon-nucleon total cross section). The geometrical cut is not allways the most restrictive way to enforce peripheral collisions, an this is a topic that should be answered by the future experiments. In both cases we noticed that at energies of the LHC operating in the heavy ion mode and for very heavy ions the double pomeron exchange is not a background for the two photon process. The situation changes considerably for light ions and mostly for the energies available at RHIC, where double pomeron exchange cannot be neglected.
## Acknowledgments
One of us (C.G.R.) thanks Paulo S. R. da Silva for useful discussions, and C. A. Bertulani for a helpful remark. This research was supported in part by the Conselho Nacional de Desenvolvimento Cientifico e Tecnologico (CNPq) (AAN), Fundação de Amparo a Pesquisa do Estado de São Paulo (FAPESP) (CGR and AAN), and by Programa de Apoio a Núcleos de Excelência (PRONEX). |
warning/0003/hep-ph0003105.html | ar5iv | text | # Exact and Approximate Dynamics of the Quantum Mechanical 𝑂(𝑁) Model
## I Introduction
Initial value problems in quantum field theory are of great interest in areas such as heavy ion collisions, dynamics of phase transitions, and early Universe physics. However, the solution of the corresponding functional Schrödinger equation is essentially impossible and one is forced to resort to approximate methods such as mean field approaches of the Hartree type or the large $`N`$ expansion. The application of variational techniques such as Hartree is limited in scope since the errors are uncontrolled. While $`1/N`$ methods promise better error control since they are based on a systematic expansion, at next-to-leading order these methods can become extremely complicated and expensive to implement. The motivation for our work in this paper is to implement the $`1/N`$ expansion at the first nontrivial order in a quantum mechanical example. Not only does this simplify the analysis but it also opens the possibility of comparing the approximate results with numerical simulations of the time dependent Schrödinger equation, a luxury not available in the field theoretic case. However, it should be kept in mind that quantum mechanics and quantum field theory are very different. For example, in the quantum mechanics applications discussed below, the $`O(1/N)`$ corrections do not correspond to inter-particle collisions (as they do in field theory) since we are restricting ourselves to one-particle quantum mechanics. Nevertheless, as discussed in more detail below, quantum mechanical examples provide excellent test-beds for key issues such as positivity violation and late-time accuracy of the approximations.
The $`O(N)`$ model has been extensively employed in time independent applications in statistical physics and quantum field theory and several recent applications have studied time-dependent phenomena. The dynamics of the chiral phase transition following the expansion of a quark-gluon plasma produced during a relativistic heavy ion collision has been modeled by an $`O(4)`$ $`\sigma `$-model at leading order in $`1/N`$. The nonequilibrium dynamics of an $`O(N)`$-symmetric $`\lambda \varphi ^4`$ theory, again treated at leading order, has been investigated in detail . Even at leading order, the $`1/N`$ expansion captures the phase transition, but does not contain enough of the dynamics to allow for rethermalization, since direct scattering first occurs at next order. The $`O(N)`$ model has been used in inflationary models of the early Universe with the scalar field often starting at the top of a hill in the potential and “rolling” down, giving rise to a quantum roll problem. It has also been applied to study primordial perturbations arising from defect models of structure formation .
The general method for obtaining the dynamical $`1/N`$ approximation via path integral techniques in quantum field theory was discussed earlier in Ref. and applied later to a quantum mechanical system of $`N+1`$ coupled oscillators (a one-dimensional version of scalar electrodynamics). Two different sets of approximate actions were considered, which differed by terms of order $`1/N^2`$, both of them being energy conserving. The first method in Ref. was a perturbative expansion of the generating functional in powers of $`1/N`$. The second method was to first Legendre transform the action to order $`1/N`$, and then find the equations of motion. When these two methods diverged from each other, they also diverged from an exact solution for the case $`N=1`$. However, due to computational restrictions it was not possible to study numerically the accuracy of the approximation as a function of $`N`$. Remedying that deficiency is the main motivation of the present study, since for the quantum roll problem numerical solutions can be obtained for arbitrary $`N`$.
One of the subtle issues in expansions involving moment based truncation schemes such as $`1/N`$, which is present both in quantum field theory and in quantum mechanics, relates to the imposition of constraints arising from the positivity of the underlying probability density function or functional. The importance of these constraints is well known in areas such as turbulence and beam dynamics . In this paper, we show that possible violations of these constraints must be tamed in $`1/N`$ expansions if the approximation is at all expected to succeed at moderate values of $`N`$. This is possible by using certain resummation schemes which we will discuss elsewhere.
In this paper we show that the naive next-to-leading order $`1/N`$ expansion violates unitarity (or more generally, positivity) leading to an instability at least for $`N`$ less than some value $`N_T`$. We have numerical evidence for a sharp threshold at $`NN_T`$ beyond which we have not been able to detect an instability. This behavior appears to be related to the nature of the effective potential at next-to-leading order: At this order, the effective potential has the property of not being defined everywhere for values of $`N<N_c`$, where $`N_c`$ depends on the values of the parameters specifying the potential, and $`N_TN_c`$. However, for $`N>N_c`$, the effective potential is defined everywhere.
A comparison of the $`1/N`$ expansion and Hartree is of interest since both agree at infinite N. At finite $`N`$, the next-to-leading order large N and Hartree approximations differ and provide alternative routes to improving the leading order result which, for the quantum roll problem, consists of harmonic oscillations in $`r^2`$ where $`r`$ is the radial degree of freedom. At finite values of $`N`$, the inclusion of nonlinearities leads to amplitude modulation effects on top of the harmonic motion. The ability to capture this modulation is a good test for the next-to-leading order large N and Hartree approximations. Our numerical results provide evidence that neither of these methods are satisfactory at late times (relative to the oscillation time), though they work reasonably well at short to intermediate times.
Our results suggest that it is important to find ways to improve the naive $`1/N`$ expansion at next-to-leading order. Work using resummation schemes is in progress and short discussions of relevant issues are included in this paper.
The paper is organized as follows. In Section II we present the $`O(N)`$ model as it pertains to quantum mechanics and in Section III we derive equations of motion for the large-$`N`$ approximation to order $`1/N`$. We derive the corresponding equations for the time dependent Hartree approximation (TDHA) in Section IV. In section V we show how the same TDHA equations can be obtained from an equal-time Green’s function approach which is computationally more attractive. The energies for the various approximations are calculated in Section VI. Section VII describes the two initial conditions which preserve the $`O(N)`$ symmetry, namely a quantum roll, and the time evolution of an offset Gaussian centered at an $`O(N)`$ symmetric point. In Section VIII we determine the effective potential to both order $`1/N`$ and for the Hartree approximation. Numerical results and comparisons with the approximations are discussed in Section IX and our conclusions are discussed in Section X.
## II The $`O(N)`$ model
The Lagrangian for the $`O(N)`$ model in quantum mechanics is given by:
$$L(x,\dot{x})=\frac{1}{2}\underset{i=1}{\overset{N}{}}\dot{x}_i^2V(r),$$
(1)
where $`V(x)`$ is a potential of the form
$$V(r)=\frac{g}{8N}\left(r^2r_0^2\right)^2,r^2=\underset{i=1}{\overset{N}{}}x_i^2.$$
(2)
The time-dependent Schrödinger equation for this problem is given by:
$$i\frac{\psi (x,t)}{t}=\left\{\frac{1}{2}\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+V(r)\right\}\psi (x,t).$$
(3)
For arbitrary initial conditions, given present computational constraints, these equations can be numerically integrated only for small $`N4`$. The initial conditions for the quantum roll problem allow a numerical solution for all $`N`$, and in this case we can attempt to study fully the behavior of the large-$`N`$ expansion. (For the shifted Gaussian initial conditions, however, this is not possible, and we used numerical solutions obtained for $`N=1`$ and $`2`$ to benchmark the large-$`N`$ approximations and the TDHA solutions at short times.)
The symmetry of the quantum roll problem is such that only the radial part of the wave function is of interest. Assuming a solution of the form
$$\psi (r,t)=r^{(1N)/2}\varphi (r,t),$$
(4)
the time dependent Schrödinger equation for $`\varphi (r,t)`$ reduces to:
$$i\frac{\varphi (r,t)}{t}=\left\{\frac{1}{2}\frac{^2}{r^2}+U(r)\right\}\varphi (r,t)$$
(5)
with an effective one dimensional potential $`U(r)`$ given by
$$U(r)=\frac{(N1)(N3)}{8r^2}+\frac{g}{8N}\left(r^2r_0^2\right)^2.$$
(6)
It is further useful to make the rescaling:
$$r^2=Ny^2,r_0^2=Ny_0^2.$$
(7)
The potential (6) then becomes:
$$u(y,N)=\frac{U(y)}{N}=\frac{(N1)(N3)}{8N^2y^2}+\frac{g}{8}(y^2y_0^2)^2,$$
(8)
corresponding to the new Schrödinger equation,
$$i\frac{\varphi (y,\stackrel{~}{t})}{\stackrel{~}{t}}=\left\{\frac{1}{2N^2}\frac{^2}{y^2}+u(y,N)\right\}\varphi (y,\stackrel{~}{t})$$
(9)
where $`\stackrel{~}{t}=Nt`$.
The method of choice to investigate the long-time behavior of the exact solution is the split-operator method, which has been presented in detail in Ref. . The wave function is expanded as a Fourier series in the radial component, and the solution is obtained as the repeated application of a time-evolution operator in symmetrically split form. As a result, the use of a Fast-Fourier Transform algorithm is required. For the purpose of the present implementation, 256 radial grid points, a value of 20 for the radial grid boundary, and a time step size of 0.01, provides a conservation of the wave function unitarity to better than 9 significant figures. The accuracy of the method has been established by comparing results with a second method, where we first solve for the eigenvalues and eigenfunctions, and then use the expansion
$$\varphi (r,t)=\underset{n}{}C_ne^{iE_nt}\varphi _n(r),$$
(10)
where $`C_n`$ was determined from the initial conditions. This method was restricted to moderate values of $`N`$. Results from the two methods agreed in the cases where they were used together.
## III The Large-$`N`$ approximation
The large $`N`$ approximation has been worked out for the $`O(N)`$ model in $`1+3`$ dimensions in Ref. . The Lagrangian (1) with the potential function (2) is obtained from that paper by specializing to $`0+1`$ dimensions, and replacing $`\varphi _a(t)x_i(t)`$, $`vr_0`$, and $`\lambda g`$.
To implement the large $`N`$ expansion, it is useful to rewrite the Lagrangian in terms of the composite field $`\chi `$ by adding a constraint term to (1), given by:
$$\frac{N}{2g}\left[\chi \frac{g}{2N}(r^2r_0^2)\right]^2,$$
(11)
which yields an equivalent Lagrangian,
$$L^{}(x,\dot{x},\chi )=\underset{i}{}\frac{1}{2}\left(\dot{x}_i^2\chi x_i^2\right)+\frac{r_0^2}{2}\chi +\frac{N}{2g}\chi ^2.$$
(12)
The generating function $`Z[j,J]`$ is given by the path integral over the classical fields $`x_i(t)`$:
$`Z[j,J]=e^{iW[j,J]}`$ $`=`$ $`{\displaystyle d\chi \underset{i}{}\mathrm{d}x_i\mathrm{exp}\left\{iS[x,\chi ;j,J]\right\}},`$
$`S[x,\chi ;j,J]`$ $`=`$ $`{\displaystyle _𝒞}dt\left\{L^{}+{\displaystyle \underset{i}{}}j_ix_i+J\chi \right\}.`$
The effective action, to order $`1/N`$, is obtained by integrating the path integral for the generating functional for the Lagrangian (12), over the $`x_i`$ variables, and approximating the integral over $`\chi `$ by the method of steepest descent (keeping terms up to order $`1/N`$). A Legendre transform of the resulting generating functional then yields the effective action, which we find to be:
$`\mathrm{\Gamma }[q,\chi ]=`$ (15)
$`{\displaystyle _𝒞}\mathrm{d}t\{{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}[\dot{q}_i^2(t)\chi (t)q_i^2(t)]+{\displaystyle \frac{i}{2}}{\displaystyle \underset{i}{}}\mathrm{ln}[G_{ii}^1(t,t)]`$
$`+{\displaystyle \frac{r_0^2}{2}}\chi (t)+{\displaystyle \frac{N}{2g}}\chi ^2(t)+{\displaystyle \frac{i}{2}}\mathrm{ln}[D^1(t,t)]\},`$
where the integral is over the close time path $`𝒞`$, discussed in Ref. and $`q(t)=x_i(t)`$. Here $`G_{ij}^1(t,t^{})`$ and $`D^1(t,t^{})`$ are the lowest order in $`1/N`$ inverse propagators for $`x_i`$ and $`\chi `$, given by
$`G_{ij}^1(t,t^{})`$ $`=`$ $`\left\{{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}t^2}}+\chi (t)\right\}\delta _𝒞(t,t^{})\delta _{ij}G^1(t,t^{})\delta _{ij},`$
$`D^1(t,t^{})`$ $`=`$ $`{\displaystyle \frac{N}{g}}\delta _𝒞(t,t^{})\mathrm{\Pi }(t,t^{}),`$
where
$`\mathrm{\Pi }(t,t^{})`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{i,j}{}}G_{ij}(t,t^{})G_{ji}(t^{},t)`$ (17)
$`+{\displaystyle \underset{i,j}{}}q_i(t)G_{ij}(t,t^{})q_j(t^{}).`$
Here $`\delta _𝒞(t,t^{})`$ is the closed time path delta function.
The equations of motion for the classical fields $`q_i(t)`$, to order $`1/N`$, are
$`\left\{{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}t^2}}+\chi (t)\right\}q_i(t)`$ (18)
$`+i{\displaystyle \underset{j}{}}{\displaystyle _𝒞}dt^{}G_{ij}(t,t^{})D(t,t^{})q_j(t^{})=0,`$ (19)
with the gap equation for $`\chi (t)`$ given by
$$\chi (t)=\frac{g}{2N}r_0^2+\frac{g}{2N}\underset{i}{}\left[q_i^2(t)+\frac{1}{i}𝒢_{ii}^{(2)}(t,t)\right].$$
(20)
The next-to-leading order $`x_i`$ propagator $`𝒢_{ij}^{(2)}(t,t^{})`$ and self energy $`\mathrm{\Sigma }_{ij}(t,t^{})`$ to order $`1/N`$ turn out to be
$`𝒢_{ij}^{(2)}(t,t^{})=G_{ij}(t,t^{})`$ (21)
$`{\displaystyle \underset{k,l}{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2G_{ik}(t,t_1)\mathrm{\Sigma }_{kl}(t_1,t_2)G_{lj}(t_2,t^{}),`$ (22)
$`\mathrm{\Sigma }_{kl}(t,t^{})=iG_{kl}(t,t^{})D(t,t^{})q_k(t)D(t,t^{})q_l(t^{}).`$ (23)
These equations agree with (2.18–2.22) of Ref. . We mention here that the actual equation for $`𝒢`$ which follows from the effective action differs from Eq. (21) in that the final $`G`$ in the integral equation is replaced by the full $`𝒢`$. This leads to a partial resummation of the $`1/N`$ corrections which which guarantees positivity of $`x^2(t)`$ (this restricted result does not imply that the full positivity problem for the density matrix has been solved). However, it does not improve the long-time accuracy of the results .
In order to solve for $`D(t,t^{})`$, we first write
$$\frac{N}{g}D(t,t^{})=\delta _𝒞(t,t^{})+\frac{N}{g}\mathrm{\Delta }D(t,t^{}).$$
(24)
Then $`\mathrm{\Delta }D(t,t^{})`$ satisfies the integral equation,
$$\frac{N}{g}\mathrm{\Delta }D(t,t^{})=\frac{g}{N}\mathrm{\Pi }(t,t^{})_𝒞dt^{\prime \prime }\mathrm{\Pi }(t,t^{\prime \prime })\mathrm{\Delta }D(t^{\prime \prime },t^{}),$$
(25)
in agreement with (2.13–2.16) of Ref. .
We are now in a position to solve these coupled equations for the motion of $`q_i(t)`$ and $`\chi (t)`$ for given initial conditions. For the initial conditions discussed in Sec. VII, we find:
$$G_{ij}(t,t^{})/i=\delta _{ij}f(t)f^{}(t^{}),$$
(26)
where $`f(t)`$ and $`f^{}(t)`$ satisfy the homogeneous equation,
$$\left\{\frac{\mathrm{d}^2}{\mathrm{d}t^2}+\chi (t)\right\}\left(\begin{array}{c}f(t)\\ f^{}(t)\end{array}\right)=0,$$
(27)
with initial conditions:
$$f(0)=\sqrt{G},\dot{f}(0)=1/(2\sqrt{G}).$$
(28)
However, $`\mathrm{\Delta }D(t,t^{})`$ cannot be factored into products of functions like $`G_{ij}(t,t^{})`$.
We solve (19) and (20) simultaneously with (21) and (25), using the Chebyshev expansion technique of Appendices A and B of Ref. .
## IV The time dependent Hartree approximation
It is useful to compare our results for the large-$`N`$ approximation to the time dependent Hartree approximation (TDHA) suitably formulated for the $`O(N)`$ problem. The static Hartree approximation is based on the idea of varying the parameters of a Gaussian wave function so as to minimize the energy (the generalization to the time-dependent case is given below). For the $`O(N)`$ problem this amounts to placing an $`N`$-dimensional Gaussian some (radial) distance away from the origin and then carrying out the minimization procedure. In contrast, the leading-order large $`N`$ wave function is a Gaussian which is locked at the origin. At infinite $`N`$ the TDHA becomes exact and equivalent to the leading order large-$`N`$ approximation, a well-known result. At finite $`N`$, the TDHA and the next-to-leading order large $`N`$ approximation can be thought of as two competing schemes to improve on the leading-order result.
There are several ways of implementing the Hartree approximation: The most common is by using the time-dependent variational principle of Dirac . This has the advantage of giving a classical Hamiltonian description for the dynamics of the variational parameters, which can be hidden in other formulations.
The idea behind this approach is that the variation of
$`\mathrm{\Gamma }[\psi ,\psi ^{}]`$ $`=`$ $`{\displaystyle dt\psi (t)|i\frac{}{t}H|\psi (t)}`$ (29)
is stationary for the exact solution of the Schrödinger equation. We consider Gaussian trial wave functions of the form:
$`\psi (x,t)=𝒩\mathrm{exp}[ip_i(t)z_i(t)`$ (30)
$`z_i(t)({\displaystyle \frac{G_{ij}^1(t)}{4}}\mathrm{\Pi }_{ij}(t))z_j(t)],`$ (31)
where $`𝒩`$ is the normalization constant, and we have set $`z_i(t)=x_iq_i(t)`$. Here $`q_i(t)`$, $`p_i(t)`$, $`G_{ij}(t)`$ and $`\mathrm{\Pi }_{ij}(t)`$ are time-dependent variational parameters, to be determined by minimizing the Dirac action. We note that $`\mathrm{\Pi }_{ij}(t)`$, which is used only in this section, is conjugate to $`G_{ij}(t)`$ and is not to be confused with the self energy $`\mathrm{\Pi }(t,t^{})`$ defined in Eq. (17).
The $`n`$-point functions can be calculated from the generating functional using the formula,
$$z_iz_j\mathrm{}z_n=\frac{^nZ[j]}{j_ij_j\mathrm{}j_n}|_{j=0},$$
(32)
where
$`Z[j]`$ $`=`$ $`𝒩^2{\displaystyle \underset{s}{}\mathrm{d}x_s}`$ (34)
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{2}}z_i(t)G_{ij}^1z_j(t)+j_iz_i(t)\right]`$
$`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{j_iG_{ij}j_j}{2}}\right].`$ (35)
The expectation value of the time derivative is given by
$`i{\displaystyle \frac{}{t}}=p_i\dot{q}_iG_{ij}\dot{\mathrm{\Pi }}_{ij}`$ (36)
and the expectation value of the kinetic energy is
$$\frac{1}{2}\frac{^2}{x_i^2}=\frac{p_ip_i}{2}+\frac{1}{8}G_{ii}^1+\mathrm{\hspace{0.17em}2}\mathrm{\Pi }_{ij}G_{jk}\mathrm{\Pi }_{ki}.$$
(37)
For the expectation value of $`V`$ we first expand the potential in a Taylor series about $`z_i=0`$,
$`V(q,z)=V(q)+V_i(q)z_i+{\displaystyle \frac{1}{2}}V_{ij}(q)z_iz_j+\mathrm{}`$
where
$`V_i(q)`$ $`=`$ $`{\displaystyle \frac{g}{2N}}q_i\left(q_sq_sr_0^2\right),`$ (38)
$`V_{ij}(q)`$ $`=`$ $`{\displaystyle \frac{g}{2N}}\left[\delta _{ij}\left(q_sq_sr_0^2\right)+2q_iq_j\right],`$ (39)
$`V_{ijk}(q)`$ $`=`$ $`{\displaystyle \frac{g}{N}}\left(\delta _{ij}q_k+\delta _{ik}q_j+\delta _{jk}q_i\right),`$ (40)
$`V_{ijkl}(q)`$ $`=`$ $`{\displaystyle \frac{g}{N}}\left(\delta _{ij}\delta _{kl}+\delta _{il}\delta _{jk}+\delta _{ik}\delta _{jl}\right).`$ (41)
Thus
$`V(q,z)={\displaystyle \frac{g}{8N}}[(q_jq_jr_0^2)^2+(z_iz_i)^2+4(z_iz_i)(z_jq_j)`$
$`+4(z_iq_i)^2+2(z_iz_i)(q_jq_jr_0^2)+4(z_iq_i)(q_jq_jr_0^2)]`$
Taking the expectation value, we obtain:
$`V`$ $`=`$ $`{\displaystyle \frac{g}{8N}}[(q_jq_jr_0^2)^2+2G_{ii}(q_jq_jr_0^2)`$ (43)
$`+4G_{ij}q_iq_j+G_{ii}G_{jj}+2G_{ij}G_{ji}]`$
The Hartree equations of motion are Hamilton’s equations for the variational parameters:
$`\dot{q}_i`$ $`=`$ $`p_i,`$ (44)
$`\dot{p}_i`$ $`=`$ $`V_i{\displaystyle \frac{1}{2}}V_{ijk}G_{jk},`$ (45)
$`\dot{G}_{ij}`$ $`=`$ $`2\left(G_{ik}\mathrm{\Pi }_{kj}+G_{jk}\mathrm{\Pi }_{ki}\right),`$ (46)
$`\dot{\mathrm{\Pi }}_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{8}}G_{ik}^1G_{kj}^12\mathrm{\Pi }_{ik}\mathrm{\Pi }_{kj}{\displaystyle \frac{1}{2}}V_{ij}{\displaystyle \frac{1}{4}}V_{ijkl}G_{kl}.`$ (47)
Solutions of this set of equations determine the time-dependent Hartree approximation to the true solution of the Schrödinger equation.
## V Method of equal-time Green’s functions
Solutions of the TDHA equation (44) require computing the matrix inverse of $`G_{ij}`$. This can be difficult to carry out in practice for large $`N`$. Fortunately, the method of equal-time Green’s functions provides a way to avoid this technical difficulty . We begin by considering the time-evolution of the one point functions:
$`\dot{q}_i`$ $`=`$ $`p_i,`$ (48)
$`\dot{p}_i`$ $`=`$ $`V_i{\displaystyle \frac{1}{2}}V_{ijk}G_{jk},`$ (49)
$`=`$ $`{\displaystyle \frac{g}{2N}}\left\{q_i(q_kq_k+G_{kk}r_0^2)+q_k(G_{ik}+G_{ki})\right\}`$ (50)
where $`V_i`$ and $`V_{ijk}`$ are given by (38), as well as the evolution of the two-point functions:
$`G_{ij}(t)`$ $`=`$ $`z_iz_j,K_{ij}=\dot{z}_i\dot{z}_j,`$ (51)
$`F_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[z_i\dot{z}_j+\dot{z}_jz_i].`$ (52)
Here we have again set $`z_i(t)=x_iq_i(t)`$. All of the expectation values are taken with respect to the Gaussian trial wave function, Eq. (30). To obtain the equations of motion for the two-point functions, we use the exact equation of motion and the factorization resulting from the Gaussian approximation. This yields
$`\dot{G}_{ij}`$ $`=`$ $`F_{ij}+F_{ji},`$ (53)
$`\dot{F}_{ij}`$ $`=`$ $`K_{ij}z_i{\displaystyle \frac{V}{x_j}},`$ (54)
$`\dot{K}_{ij}`$ $`=`$ $`\left[{\displaystyle \frac{V}{x_i}}\dot{z}_j\dot{z}_i{\displaystyle \frac{V}{x_j}}\right].`$ (55)
Here we have used the Lagrange equations of motion
$$\ddot{x}_i+\frac{V}{x_i}=0,$$
(56)
where
$`{\displaystyle \frac{V}{x_i}}`$ $`=`$ $`V_{ij}z_j+{\displaystyle \frac{1}{6}}V_{ijkl}z_jz_kz_l`$ (58)
$`+\text{terms with even powers of }z_i,`$
and the fact that for our Gaussian wave packet, $`z_i=0.`$ The canonical commutation relations give:
$$z_i\dot{z}_j\dot{z}_jz_i=x_i\dot{x}_j\dot{x}_jx_i=i\delta _{ij},$$
(59)
and we have
$`z_iz_jz_kz_l=G_{ij}G_{kl}+G_{il}G_{jk}+G_{ik}G_{jl},`$ (60)
$`z_iz_jz_k\dot{z}_l=G_{ij}F_{kl}+F_{il}G_{jk}+G_{ik}F_{jl}`$ (61)
$`+i\left(G_{ij}\delta _{kl}+G_{jk}\delta _{il}+G_{ik}\delta _{jl}\right)/2,`$ (62)
$`\dot{z}_iz_jz_kz_l=F_{ji}G_{kl}+F_{li}G_{jk}+F_{ki}G_{jl}`$ (63)
$`i\left(G_{kl}\delta _{ji}+G_{jk}\delta _{li}+G_{jl}\delta _{ki}\right)/2.`$ (64)
Finally, from Eqs. (54) and (55) we get:
$`\dot{F}_{ij}`$ $`=`$ $`K_{ij}V_{jk}G_{ik}`$ (66)
$`V_{jklm}\left(G_{ik}G_{lm}+G_{im}G_{kl}+G_{il}G_{km}\right)/6,`$
$`\dot{K}_{ij}`$ $`=`$ $`V_{ik}F_{kj}V_{jk}F_{ki}`$ (69)
$`V_{iklm}\left(G_{kl}F_{mj}+G_{lm}F_{kj}+G_{km}F_{lj}\right)/6`$
$`V_{jklm}\left(F_{ki}G_{lm}+F_{mi}G_{kl}+F_{li}G_{km}\right)/6.`$
For Gaussian initial conditions, the equal-time Green’s function method is assured to give the same result as the Hartree method, if $`F_{ij}(0)`$ and $`K_{ij}(0)`$ satisfy the requirements:
$`F_{ij}(0)=0,K_{ij}(0)=G_{ij}^1(0)/4.`$
Choosing $`K_{ij}`$ independently of $`G_{ij}`$, corresponds to a mixed initial density matrix, rather than a pure state. If we further choose $`G_{ij}(0)`$ to be diagonal and equal to the same number $`G_0`$,
$`G_{ij}(0)=\delta _{ij}G_0,`$
then $`K_{ij}(0)`$ is given by:
$`K_{ij}(0)=\delta _{ij}/(4G_0).`$
For the initial conditions pertinent to the quantum roll, $`q_i(t)=0`$ for all $`t`$. Then $`G_{ij}(t)`$, $`F_{ij}(t)`$, and $`K_{ij}(t)`$ are all proportional to the unit matrix, and have no off-diagonal terms. For the offset initial condition, we choose large-$`N`$ symmetric initial conditions so that $`q_i(0)=q_0`$, and $`p_i(0)=0`$, and $`G_{ij}(0)=G_0\delta _{ij}`$. In that case, all the $`q`$’s and $`p`$’s are identical,
$`q_i(t)=q(t),p_i(t)=p(t),`$
and the matrices $`G_{ij}(t)`$, $`F_{ij}(t)`$ and $`K_{ij}(t)`$ become off-diagonal in a simple way so that all the diagonal elements are equal and all the off-diagonal elements are equal. That is, we can write
$`G_{ij}(t)`$ $`=`$ $`G(t)\delta _{ij}+\overline{G}(t)(1\delta _{ij}),`$
$`F_{ij}(t)`$ $`=`$ $`F(t)\delta _{ij}+\overline{F}(t)(1\delta _{ij}),`$
$`K_{ij}(t)`$ $`=`$ $`K(t)\delta _{ij}+\overline{K}(t)(1\delta _{ij}).`$
For this case, Eqs. (48), (53), (66), and (69) simplify to the following set of coupled equations:
$`\dot{q}`$ $`=`$ $`p,`$ (70)
$`\dot{p}`$ $`=`$ $`{\displaystyle \frac{g}{2N}}q\left\{Nq^2r_0^2+(N+2)G+2(N1)\overline{G}\right\},`$ (71)
$`\dot{G}`$ $`=`$ $`2F,\dot{\overline{G}}=2\dot{\overline{F}},`$ (72)
$`\dot{F}`$ $`=`$ $`K{\displaystyle \frac{g}{2N}}\{G[(N+2)(q^2+G)r_0^2]`$ (74)
$`+2(N1)\overline{G}(q^2+\overline{G})\},`$
$`\dot{\overline{F}}`$ $`=`$ $`\overline{K}{\displaystyle \frac{g}{2N}}\{\overline{G}[(3N2)q^2+(N+4)G`$ (76)
$`+2(N2)\overline{G}r_0^2]+2Gq^2\},`$
$`\dot{K}`$ $`=`$ $`{\displaystyle \frac{g}{N}}\{F[(N+2)(q^2+G)r_0^2]+`$ (78)
$`+\overline{F}\mathrm{\hspace{0.17em}2}(N1)(q^2+\overline{G})\},`$
$`\dot{\overline{K}}`$ $`=`$ $`{\displaystyle \frac{g}{N}}\{\overline{F}[(3N2)q^2+(N+2)G`$ (80)
$`+2(N2)\overline{G}r_0^2]+F\mathrm{\hspace{0.17em}2}(q^2+\overline{G})\}.`$
If we let $`r_0^2=Ny_0^2`$, and then take the limit $`N\mathrm{}`$, we recover the leading order in the large-$`N`$ result, as discussed in Ref. .
The equal time Green’s function method is easier to implement numerically than the Hamiltonian system described by Eq. (44), since no matrix inversion is involved. However, if one wants to find the wave function or the energy, instead of just obtaining the Green’s functions, matrix inversion is once again required.
## VI Energy
It is important to note that even though the Hartree and large $`N`$ approximations are truncations of the true dynamics, they are nevertheless energy conserving. In the large-$`N`$ approximation, to order $`1/N`$, the expectation value of the Hamiltonian is given by
$`E`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\left[\dot{x_i}^2(t)+\widehat{\chi }(t)x_i^2(t)\right]`$ (82)
$`{\displaystyle \frac{r_0^2}{2}}\widehat{\chi }(t){\displaystyle \frac{N}{2g}}\widehat{\chi }^2(t).`$
We write these expectation values in terms of the CTP Green’s functions. By definition, the disconnected two-point Green’s functions are introduced as
$`D_{\mathrm{dis}}(t,t^{})`$ $`=`$ $`i\mathrm{T}_𝒞[\widehat{\chi }(t)\widehat{\chi }(t^{})]`$ (83)
$`=`$ $`i\chi (t)\chi (t^{})+𝒟(t,t^{}),`$ (84)
$`G_{ij,\mathrm{dis}}(t,t^{})`$ $`=`$ $`i\mathrm{T}_𝒞[x_i(t)x_j(t^{})]`$ (85)
$`=`$ $`iq_i(t)q_j(t^{})+𝒢_{ij}(t,t^{}),`$ (86)
where $`𝒟`$ and $`𝒢_{ab}`$ denote the connected two-point Green’s functions
$`𝒟(t,t^{})`$ $`=`$ $`\left[{\displaystyle \frac{\delta ^2W[J,j]}{\delta J(t)\delta J(t^{})}}\right]_{J,j=0},`$ (87)
$`𝒢_{ij}(t,t^{})`$ $`=`$ $`\left[{\displaystyle \frac{\delta ^2W[J,j]}{\delta j_i(t)\delta j_j(t^{})}}\right]_{J,j=0}.`$ (88)
To obtain the energy from (82), we require the expectation values
$`\widehat{\chi }(t)`$ $`=`$ $`\chi (t),`$
$`\widehat{\chi }^2(t)`$ $`=`$ $`\chi ^2(t)+D(t,t)/i,`$
$`x_i^2(t)`$ $`=`$ $`q_i^2(t)+𝒢_{ii}(t,t)/i,`$
$`\dot{x}_i^2(t)`$ $`=`$ $`\dot{q}_i^2(t)+{\displaystyle \frac{^2𝒢_{ii}(t,t^{})/i}{tt^{}}}|_{t=t^{}},`$
$`\widehat{\chi }(t)x_i^2(t)`$ $`=`$ $`\chi (t)\left[q_i^2(t)+𝒢_{ii}(t,t)/i\right]K_{ii}(t,t,t),`$
where $`K_{ij}(t_1,t_2,t_3)`$ is the 3-point Green’s function defined as
$`K_{ij}(t_1,t_2,t_3)={\displaystyle _𝒞}dtG_{ik}(t_1,t)G_{kj}(t_2,t)D(t,t_3).`$
The energy for the next-to-leading order large-N approximation is then given by:
$`E`$ $`=`$ $`{\displaystyle \frac{r_0^2}{2}}\chi (t){\displaystyle \frac{N}{2g}}\left\{\chi ^2(t)+\mathrm{\Delta }D(t,t)/i\right\}`$ (91)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\left\{\dot{q}_i^2(t)+{\displaystyle \frac{^2𝒢_{ii}(t,t^{})/i}{tt^{}}}|_{t=t^{}}\right\}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\left\{\chi (t)\left[q_i^2(t)+𝒢_{ii}(t,t)/i\right]K_{ii}(t,t,t)\right\}.`$
Using the equations of motion, we can show directly that (91) is conserved.
It is easy to evaluate the energy at $`t=0`$ for the quantum roll problem using the initial conditions (28). The result for the leading and next-to-leading order large-N approximation is:
$`{\displaystyle \frac{E}{N}}=ϵ_0+{\displaystyle \frac{1}{N}}ϵ_1+\mathrm{}`$
where,
$`ϵ_0`$ $`=`$ $`{\displaystyle \frac{1}{8G}}+{\displaystyle \frac{1}{8}}gy_0^4{\displaystyle \frac{1}{4}}gGy_0^2+{\displaystyle \frac{1}{8}}gG^2,`$ (92)
$`ϵ_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}gG^2.`$ (93)
Our initial wave function was chosen to be Gaussian, so that the parameters of the Hartree approximation agree exactly with the energy and parameters of the exact wave function at $`t=0`$. However, the leading order in the large-$`N`$ approximation for the same value of $`G`$ will disagree with the exact energy by $`ϵ_1`$. This discrepency dissapears when we include the $`1/N`$ corrections.
Since the Hartree approximation leads to a canonical Hamiltonian dynamical system, the corresponding energy in that approximation is also a constant of the motion. It is given by:
$`E`$ $`=`$ $`{\displaystyle \frac{1}{2}}p_i^2+{\displaystyle \frac{1}{8}}G_{ii}^1+2\mathrm{\Pi }_{ij}G_{jk}\mathrm{\Pi }_{ki}`$ (96)
$`+V(q)+{\displaystyle \frac{1}{2}}V_{ij}G_{ij}`$
$`+{\displaystyle \frac{1}{4!}}V_{ijkl}\left(G_{ij}G_{kl}+G_{il}G_{jk}+G_{ik}G_{jl}\right).`$
We used this expression to check the accuracy of our numerical solutions. As with the next-to-leading order $`1/N`$ expression, for the quantum roll initial condition, Eq. (96) agrees with the exact result.
## VII Initial Conditions
### A Quantum Roll
We wish to study initial conditions which are consistent with $`O(N)`$ symmetry. This implies immediately that all the $`x_i(t)`$ have to be identical, with $`x_i(0)=0`$, and $`G_{ij}(t)`$ must be diagonal. The quantum roll problem is defined by a Gaussian initial wave function that is centered on the origin:
$$\psi _0(r)=\frac{1}{(2\pi G)^{N/4}}\mathrm{exp}\left\{\frac{r^2}{4G}\right\}.$$
(97)
In this section, $`GG(0)`$.
One of the difficulties in studying the systematics of the $`1/N`$ expansion is the fact that, at next to leading order, every different value of $`N`$ (with all other parameters held constant) defines a different initial value problem. In this sense one cannot naively compare individual solutions, exact or approximate, at different values of $`N`$. In effect one has to tune the parameters of the problem at each $`N`$ in order to maintain certain invariance properties which allow different $`N`$ evolutions to be compared to each other. This parameter tuning process is described below.
Since the infinite $`N`$ limit has very precise properties, several technical issues arise when one wants to approach this limit starting at $`N=1`$ in a uniform manner. To study the large $`N`$ limit it is convenient to make a rescaling to the $`y`$ variables, given in (7). At very large $`N`$, the potential energy $`u(y,N)`$ is (8):
$`u(y,N)`$ $`=`$ $`{\displaystyle \frac{(N1)(N3)}{8N^2y^2}}+{\displaystyle \frac{g}{8}}(y^2y_0^2)^2,`$ (98)
$``$ $`{\displaystyle \frac{1}{8y^2}}+{\displaystyle \frac{g}{8}}(y^2y_0^2)^2.`$ (99)
In this limit, $`u(N,y)`$ has a minimum which is independent of $`N`$, and the large $`N`$ limit consists of harmonic oscillations about this minimum \[the reason for this is that the large $`N`$ limit also corresponds to an effectively large mass limit in the Schrödinger equation (9)\].
One way to uniformly study the motion of a wave packet as a function of $`N`$ is to choose initial conditions so that there is a uniform overlap of the initial wave function with the set of eigenfunctions of the Hamiltonian in the $`N\mathrm{}`$ limit. We can obtain this constant overlap if we allow the coupling constant $`g`$ to be a slowly varying function of $`N`$. This can be done in several ways that differ by terms of order $`1/N^2`$. The method presented below leads to uniform results even at $`N=1`$ as we change the parameters with $`N`$. Our method is to keep the distance between the centers of the initial wave function and the position of the minimum of the potential a constant as $`N`$ is varied.
Using (97), we define $`\stackrel{~}{r}`$ by:
$$\stackrel{~}{r}=r=\sqrt{2G}\left[\frac{\mathrm{\Gamma }((N+1)/2)}{\mathrm{\Gamma }(N/2)}\right],$$
(100)
and $`G_{\mathrm{}}`$ by the variance:
$`{\displaystyle \frac{G_{\mathrm{}}}{2}}=r^2r^2=G\left\{N2\left[{\displaystyle \frac{\mathrm{\Gamma }((N+1)/2)}{\mathrm{\Gamma }(N/2)}}\right]^2\right\}.`$
Solving the above equations for $`G`$, we have
$$G(N)=\frac{G_{\mathrm{}}}{2N4\left[\mathrm{\Gamma }((N+1)/2)/\mathrm{\Gamma }(N/2)\right]^2}.$$
(101)
Substitution of this expression into (100) yields
$$\stackrel{~}{r}(N)=\sqrt{\frac{G_{\mathrm{}}}{N[\mathrm{\Gamma }(N/2)/\mathrm{\Gamma }((N+1)/2)]^22}}.$$
(102)
In the limit when $`N`$ goes to infinity, we have
$`G(N)G_{\mathrm{}},\stackrel{~}{r}(N)\stackrel{~}{r}_{\mathrm{}}=\sqrt{(N1)G_{\mathrm{}}},`$
which defines $`\stackrel{~}{r}_{\mathrm{}}`$, and agrees with the asymptotic form of the rescaled version of the initial wave function (97):
$`\varphi _0(r)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi G)^{N/4}}}\mathrm{exp}\left\{{\displaystyle \frac{r^2}{4G}}+{\displaystyle \frac{N1}{2}}\mathrm{ln}r\right\},`$
$``$ $`{\displaystyle \frac{1}{(2\pi G)^{N/4}}}\mathrm{exp}\left\{{\displaystyle \frac{(r\stackrel{~}{r}_{\mathrm{}})^2}{2G_{\mathrm{}}}}+O(1/\sqrt{N})\right\}.`$
In order to ensure that the initial wave function has a finite overlap with the energy eigenfunctions of the Schrödinger equation at large $`N`$, we will keep the value of $`G_{\mathrm{}}`$ (and not $`G`$) fixed in our simulations.
Another quantity that should be kept constant is the basic oscillation frequency. In order to do this, we first find the Gaussian oscillations about the minimum of the one dimensional potential, defined by Eq. (6). We expand $`U(r)`$ as
$$U(r)=U(\overline{r})+\frac{1}{2}\overline{m}^2(r\overline{r})^2+\mathrm{},$$
(103)
where $`\overline{r}`$ is given by the solution of the equation,
$$\frac{(N1)(N3)}{4\overline{r}^4}=\frac{g}{2N}(\overline{r}^2r_0^2),$$
(104)
and $`\overline{m}^2`$ by
$`\overline{m}^2`$ $`=`$ $`{\displaystyle \frac{3(N1)(N3)}{4\overline{r}^4}}+{\displaystyle \frac{g}{2N}}(3\overline{r}^2r_0^2),`$ (105)
$`=`$ $`{\displaystyle \frac{g}{N}}(3\overline{r}^22r_0^2).`$ (106)
The frequency of oscillation is determined by $`\overline{m}`$, and this is the quantity to be kept fixed as $`N`$ is changed.
The last technical issue is to keep the distance between the center of the initial wave function, $`\stackrel{~}{r}`$, and the minimum of the potential, $`\overline{r}`$, a constant as we vary $`N`$. That is, we keep
$`\delta r=\overline{r}\stackrel{~}{r},`$
constant for all $`N`$.
With this strategy of keeping $`G_{\mathrm{}}`$, $`\overline{m}`$ and $`\delta r`$ fixed, we can now determine how the coupling constant must vary with $`N`$. We first define $`m^2`$ to be the second derivative of $`V(r)`$ evaluated at $`r=\overline{r}`$,
$$m^2=\frac{\mathrm{d}^2V(r)}{\mathrm{d}r^2}|_{r=\overline{r}}=\frac{g}{2N}(3\overline{r}^2r_0^2).$$
(107)
Then, (105) becomes
$$m^2(N)=\overline{m}^2\frac{3(N1)(N3)}{4\overline{r}^4}.$$
(108)
Solving (107) for $`r_0`$, substituting into (104) and solving for $`g`$ gives:
$$g(N)=\frac{N}{\overline{r}^2}\left\{\overline{m}^2\frac{(N1)(N3)}{\overline{r}^4}\right\}.$$
(109)
with $`\overline{r}=\stackrel{~}{r}+\delta r`$. The value of $`r_0^2`$ is then determined by (104):
$$r_0^2(N)=\overline{r}^2\frac{2N}{g(N)}\frac{(N1)(N3)}{4\overline{r}^4}.$$
(110)
Thus, for fixed values of $`G_{\mathrm{}}`$, $`\overline{m}`$ and $`\delta r`$, Eqs. (101), (102), (109), and (110) determine values for $`G(N)`$, $`\stackrel{~}{r}(N)`$, $`g(N)`$, and $`r_0(N)`$ all values of $`N`$.
In the limit, $`N\mathrm{}`$, we find that
$`g(\mathrm{})`$ $`=`$ $`{\displaystyle \frac{1}{G_{\mathrm{}}}}\left(\overline{m}^2{\displaystyle \frac{1}{G_{\mathrm{}}^2}}\right)`$ (111)
$`m^2(\mathrm{})`$ $`=`$ $`\overline{m}^2{\displaystyle \frac{3}{4G_{\mathrm{}}^2}}`$ (112)
To summarize, in order to establish appropriate initial conditions for the quantum roll problem, we have kept the variance $`G_{\mathrm{}}`$ constant instead of $`G`$, and have allowed the parameters describing the potential function $`g`$ and $`r_0`$ to change with $`N`$ in order to compare solutions that have close to the same oscillation frequencies. In our numerical runs, we chose the values
$$G_{\mathrm{}}=1;\overline{m}^2=2;\delta r=2;g(\mathrm{})=1$$
(113)
Fig. 1 displays the variation of the potential parameters with $`N`$.
### B Shifted Gaussian Initial Conditions
The second $`O(N)`$ invariant initial condition we investigated had a wave function localized in a wave packet near the center of the valley of the classical potential at $`r=r_0`$. For $`N=1`$ this would be the standard double-well tunnelling problem; for higher values of $`N`$, tunneling is avoided by going around the barrier. Therefore this initial condition is qualitatively different from the roll problem and provides a different arena for testing approximations. However, since this initial condition violates the $`O(N)`$ symmetry of the potential, numerical solution is at present possible only for very small $`N`$.
We take the initial wave function to be a shifted Gaussian of the form:
$$\psi _0(x)=\frac{1}{(2\pi G)^{N/4}}\mathrm{exp}\left\{\underset{i}{}\frac{(x_ir_0/\sqrt{N})^2}{4G}\right\}.$$
(114)
The energy $`E`$ of this state can be determined from Eq. (96) by the substitutions,
$`G_{ij}\delta _{ij}G;q_ir_0/\sqrt{N};p_i0;\mathrm{\Pi }_{ij}0,`$
from which we find:
$$E=\frac{N}{8G}+\frac{g}{8N}\left\{N(N+2)G^22NGr_0^2+r_0^4\right\}.$$
(115)
On the other hand, the height of the classical potential barrier is given by:
$$E_b=\frac{g}{8N}r_0^4.$$
(116)
For $`N=1`$, the necessary requirement for tunneling is that $`E<E_b`$.
In the general case (arbitrary $`N`$), we have:
$$M^2=\frac{^2V(r)}{r^2}|_{r_0}=\frac{g}{N}r_0^2\mathrm{or}r_0^2=\frac{N}{g}M^2.$$
(117)
If the initial state is close to the ground state of a harmonic potential that approximates the potential at the bottom of the well, then the width $`G`$ of the wave function is, approximately,
$$G=\frac{1}{2\sqrt{M^2}},$$
(118)
which can be combined with Eq. (115) to give the desired energy of the initial state in terms of the values of $`N`$ and $`g`$.
We are interested in initial conditions where the energy per oscillator does not increase as a function of $`N`$. To implement this we fix $`M^2=1`$, which corresponds to $`G=1/2`$ for the initial width. The barrier height is then given by
$$E_b=\frac{N}{8g},$$
(119)
and the total energy by
$$E=\frac{N+1}{4}+\frac{N+2}{32}g$$
(120)
We explored three cases: $`E=0.5E_b`$, $`E=E_b`$, and $`E=2E_b`$. For each of these cases, Eqs. (119) and (120) determine $`g`$ for each $`N`$. In all cases we took $`x_i(0)=r_0/\sqrt{N}`$, $`\dot{x}_i=0`$, $`G_{ij}(0)=G\delta _{ij}`$, and $`\dot{G}_{ij}(0)=0`$. As a consequence, all of the oscillators $`x_i(t)`$ move identically.
## VIII Effective potential
It is well-known that the static effective potential is not always a useful guide to the true dynamics of the system (See, e.g., Ref. ). Nevertheless, one may seek to gain qualitative insight into some aspects of quantum dynamics this way, though care is certainly indicated (See, e.g., Ref. for the Gaussian effective potential). Indeed, there appears to be an interesting connection with the properties of the effective potential at next-to-leading order and with the corresponding dynamical evolution (discussed in the next section).
The effective potential in the large $`N`$ approximation has been previously obtained by Root to order $`1/N`$; however, we recalculate it here using our equations. When $`x_i`$ and $`\chi `$ are independent of time, we can ignore the closed time path ordering and use Fourier transforms, passing the poles by using the Feynman contour. Then, from the action given in (15), we find
$`V_{\text{eff}}^{[1]}(r,\chi )={\displaystyle \frac{N\chi }{g}}\left(\mu ^2{\displaystyle \frac{\chi }{2}}\right)+{\displaystyle \frac{1}{2}}\chi r^2`$ (122)
$`+{\displaystyle \frac{N}{2}}{\displaystyle \frac{dk}{2\pi i}\mathrm{ln}[\stackrel{~}{G}^1(k)]}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{dk}{2\pi i}\mathrm{ln}[\stackrel{~}{D}^1(k)]},`$
where $`\chi `$ satisfies the requirement
$`{\displaystyle \frac{}{\chi }}V_{\text{eff}}(r,\chi )=0.`$ (123)
In this section to make contact with Root , we have $`\mu ^2=gr_0^2/(2N)<0`$. In order to examine the large $`N`$ limit, we again rescale (7) to the $`y`$ variables. Then for the Green’s functions, we find:
$`\stackrel{~}{G}^1(k)`$ $`=`$ $`(k^2\chi ),`$
$`\stackrel{~}{D}^1(k)`$ $`=`$ $`{\displaystyle \frac{N}{g}}Ny^2\stackrel{~}{G}(k)+{\displaystyle \frac{iN}{2}}{\displaystyle \frac{dp}{2\pi }\stackrel{~}{G}(p)\stackrel{~}{G}(kp)}`$
$`=`$ $`{\displaystyle \frac{N}{g}}\left\{1g{\displaystyle \frac{y^2}{k^2\chi }}{\displaystyle \frac{g}{2\sqrt{\chi }}}{\displaystyle \frac{1}{k^24\chi }}\right\}`$
$`=`$ $`{\displaystyle \frac{N}{g}}{\displaystyle \frac{(k^2m_+^2)(k^2m_{}^2)}{(k^24\chi )(k^2\chi )}},`$
where $`m_\pm ^2=b\pm \sqrt{b^2c}`$, with
$`b`$ $`=`$ $`{\displaystyle \frac{5}{2}}\chi +{\displaystyle \frac{g}{2}}\left(y^2+{\displaystyle \frac{1}{2\sqrt{\chi }}}\right)`$
$`c`$ $`=`$ $`4\chi ^2+g\left(4y^2\chi +{\displaystyle \frac{1}{2}}\sqrt{\chi }\right).`$
For the Feynman contour, we have
$$\frac{dk}{2\pi i}\mathrm{ln}(k^2\chi )=\sqrt{\chi }+\text{constant terms}.$$
(124)
Thus the effective potential (122) becomes
$`{\displaystyle \frac{V_{\text{eff}}^{[1]}(y,\chi )}{N}}`$ $`=`$ $`{\displaystyle \frac{\chi }{2}}\left(y^2y_0^2\right){\displaystyle \frac{\chi ^2}{2g}}+{\displaystyle \frac{\sqrt{\chi }}{2}}`$ (126)
$`+{\displaystyle \frac{1}{2N}}\left(m_++m_{}3\sqrt{\chi }\right).`$
The gap equation which determines $`\chi `$ follows from (123)
$`\chi `$ $`=`$ $`{\displaystyle \frac{g}{2}}\left(y^2y_0^2\right)+{\displaystyle \frac{g(N3)}{4N\sqrt{\chi }}}+{\displaystyle \frac{g}{2N}}{\displaystyle \frac{(m_++m_{})}{\chi }}.`$ (127)
To leading order in the large $`N`$ expansion, Eqs. (126, 127) reduce to the parametric set
$`{\displaystyle \frac{V_{\text{eff}}^{[0]}(\chi )}{N}}`$ $`=`$ $`{\displaystyle \frac{\chi ^2}{2g}}+{\displaystyle \frac{\sqrt{\chi }}{4}},`$ (128)
$`y^2(\chi )`$ $`=`$ $`y_0^2+{\displaystyle \frac{2}{g}}\chi {\displaystyle \frac{1}{2\sqrt{\chi }}}.`$ (129)
Equations (126) and (127) agree with Root, however he used the leading order expression for $`\chi `$ in (129), rather than the full $`\chi `$ of (127).
There exist two real solutions of Eq. (127) for $`\chi `$ with $`y`$ greater than some minimum value $`y_{\text{min}}`$. The next-to-leading order large $`N`$ effective potential, from Eq. (126) is therefore double valued for $`y>y_{\text{min}}`$, and does not exist for smaller values of $`y`$. The physical solution branch corresponds to the one that matches on to the leading order result; the other branch is an unphysical solution. Since it follows from a Legendre transformation, the effective potential (at any order in $`1/N`$) has to be a convex function. The nonexistence of the effective potential at $`y<y_{\text{min}}`$ implies that no quantum state can be associated with the next-to-leading order large $`N`$ approximation in this range.
In Fig. 2, we plot the physical branch of the effective potential as a function of $`y`$, for values of $`N`$ from $`1`$ to $`100`$, for the case $`g=1`$ and $`y_0=2`$. For comparison, we also show in this figure the leading order potential function from Eq. (129), which is single-valued and finite for all $`y`$. (In contrast to the next-to-leading order case we can always associate a Gaussian wave function with the leading order approximation.)
In the case of the Hartree approximation, one can define an “effective potential” as the expectation value of the Hamiltonian using the variational wave function (30) for static configurations . Setting $`p_i(t)=0`$ and $`\mathrm{\Pi }_{ij}(t)=0`$, and putting $`_iq_i^2=r^2`$ and $`G_{ij}=\delta _{ij}G`$ in Eqs. (37) and (43), we find:
$`{\displaystyle \frac{V_{\text{eff}}^{[H]}(y,G)}{N}}`$ $`=`$ $`{\displaystyle \frac{1}{8G}}+{\displaystyle \frac{g}{8}}\left(y^2y_0^2\right)^2`$ (131)
$`+g{\displaystyle \frac{N+2}{4N}}\left(y^2+{\displaystyle \frac{1}{2}}G\right)G{\displaystyle \frac{g}{4}}y_0^2G,`$
The value of $`G`$ is fixed by the requirement that
$$\frac{V_{\text{eff}}^{[H]}(y,G)}{G}=0,$$
which gives the gap equation for the Hartree approximation:
$$\chi =\frac{g}{2}(y^2y_0^2)+\frac{g}{N}\left(y^2+\frac{1}{2\chi }\right)+\frac{g}{4\sqrt{\chi }},$$
(132)
where we have set $`G=1/2\sqrt{\chi }`$. Parametric equations for the Hartree effective potential are then given by:
$`{\displaystyle \frac{V_{\text{eff}}^{[H]}(\chi )}{N}}`$ $`=`$ $`{\displaystyle \frac{1}{2g}}\left({\displaystyle \frac{N}{N+2}}\right)^2\chi ^2{\displaystyle \frac{N}{(N+2)^2}}y_0^2\chi `$ (135)
$`+{\displaystyle \frac{N+4}{4(N+2)}}\sqrt{\chi }+{\displaystyle \frac{g}{2(N+2)^2}}y_0^4`$
$`+{\displaystyle \frac{g}{4(N+2)}}{\displaystyle \frac{y_0^2}{\sqrt{\chi }}}{\displaystyle \frac{g}{16N}}{\displaystyle \frac{1}{\chi }}`$
$`y^2(\chi )`$ $`=`$ $`{\displaystyle \frac{N}{N+2}}\left(y_0^2+{\displaystyle \frac{2}{g}}\chi {\displaystyle \frac{1}{2\sqrt{\chi }}}\right){\displaystyle \frac{1}{(N+2)\chi }}.`$ (136)
Note that in the limit $`N\mathrm{}`$, Eq. (135) reduces to Eq. (129), the leading order large $`N`$ result. Note also that the Hartree effective potential is not derived from a Legendre transform and hence is not subject to a convexity constraint.
In Fig. 3, we plot the Hartree effective potential from Eq. (135) as a function of $`y`$, for our chosen parameters $`g_{\mathrm{}}=1`$ and $`y_0=2`$, for different values of $`N`$. In contrast to the smooth behavior exhibited by the large $`N`$ effective potentials, the Hartree effective potential shows a “first-order transition” in the placement of the minimum of the potential as a function of $`N`$.
The minimum of the effective potential corresponds to a determination of the ground state energy. In Fig. 4, we show values of the minimum energies of the large $`N`$ and Hartree effective potentials as a function of $`N`$. For $`N2`$, the Hartree minimum is generally greater than that for the next-to-leading order large-$`N`$ approximation. Since the Hartree approximation is a variational ansatz it gives an upper bound to the minimum energy. The fact that the next-to-leading order large $`N`$ results are lower than this bound is encouraging, although no guarantee of absolute accuracy.
It is interesting to ask the question how the point $`y_{\text{min}}`$, below which the next-to-leading order large $`N`$ effective potential does not exist, changes as a function of $`N`$. We know that at “infinite $`N`$,” (leading order), $`y_{\text{min}}=0`$ but it is important to know how this limit is reached. For instance, is there a finite value of $`N`$ beyond which $`y_{\text{min}}=0`$? In Fig. 5, we plot $`y_{\text{min}}`$ as a function of $`N`$, for the Hartree and next-to-leading order large $`N`$ approximations. As already stated, the Hartree approximation displays a first order phase transition between the broken and unbroken symmetry solutions at $`N=6.2`$, whereas for the next-to-leading order large $`N`$ approximation a different type of behavior is found: for $`N18.6`$, $`y_{\text{min}}`$ is finite, but for $`N18.6`$, it hits the origin. Thus for $`N18.6`$, we can associate a quantum state (though not known explicitly) with the next-to-leading order approximation.
The critical value of $`N`$ is fixed by the value of $`\chi `$ at the gap equation at the inflection point. If we write the gap equation (127) as
$$f(\chi ,y^2,N)=f_0(\chi ,y^2)+\frac{1}{N}f_1(\chi ,y^2)=0,$$
(137)
where
$`f_0(\chi ,y^2)`$ $`=`$ $`{\displaystyle \frac{g}{2}}(y^2y_0^2)+{\displaystyle \frac{g}{4}}{\displaystyle \frac{1}{\sqrt{\chi }}}\chi ,`$
$`f_1(\chi ,y^2)`$ $`=`$ $`{\displaystyle \frac{3g}{4}}{\displaystyle \frac{1}{\sqrt{\chi }}}+{\displaystyle \frac{g}{2}}{\displaystyle \frac{(m_++m_{})}{\chi }}.`$
Then the critical point is determined by
$$N_\text{c}=\frac{f_1(\overline{\chi },0)}{f_0(\overline{\chi },0)}=\frac{f_1(\overline{\chi },0)/\overline{\chi }}{f_0(\overline{\chi },0)/\overline{\chi }},$$
(138)
where $`\overline{\chi }`$ is given by the solution of this system of equations. For the parameters of Eq. (113), we find numerically that $`N_\text{c}=18.60`$ in excellent agreement with the results shown in Fig. 5.
## IX Numerical Results
### A Quantum Roll
We begin with a discussion of our results for the quantum roll problem. We first examine the short time behavior, $`0<t<3`$, to see if the next-to-leading order large $`N`$ approximation gives an improvement over the leading order solutions. In Fig. 6, we plot the values of $`r^2/N`$ from the numerical solution, the leading and next-to-leading order large-$`N`$ approximations, and the Hartree approximation, for $`N=20`$. The next-to-leading order large $`N`$ approximation is clearly better than the leading order solution, and also better than the Hartree results. Similar behavior is seen for other values of $`N`$ (we also ran $`N=50,80,`$ and 100).
The long time behavior of these approximations is typically of much more interest. We examined behavior over times $`0<t<100`$ to see how long the approximations remained viable. Fig. 7 displays $`r^2/N`$ for the numerical solution, the leading and next-to-leading order large-$`N`$ approximations, and the Hartree approximation for $`N=20`$ and 100. The next-to-leading order large-$`N`$ approximation for $`N=20`$ blows up at $`t84`$. This instability is connected to a violation of unitarity in the particular implementation of this approximation and will be discussed in greater detail below. In general, at these moderate values of $`N`$, the approximations track the numerical solutions reasonably well though they do get out of phase as time progresses. As $`N`$ is increased, the phase errors are considerably reduced as is apparent in the results for $`N=100`$.
The energy of the next-to-leading order large $`N`$ and Hartree approximations is the same as the exact one, but the energy of the leading order large $`N`$ approximation differs from it by terms of order $`1/N`$. (This is because we need to keep the initial values of the parameters the same.) To make a comparison between the approximations this difference has to be compensated for; we do this by rescaling time by a constant multiplicative factor so as to match the last oscillation maxima. This effect is of order $`1/N`$. For $`N=100`$, we find that for $`0<t<100`$, the next-to-leading order $`1/N`$ approximation is always more accurate than the leading order; however, when comparing the next-to-leading order with the Hartree, although less accurate for $`t<50`$, the Hartree approximation starts becoming more accurate at $`t50`$ (however, the errors are very similar in magnitude, of the order of a few per cent).
We now return to a discussion of the blow ups first encountered in the $`N=20`$ case discussed above. The failure of truncation schemes (of which $`1/N`$ is an example) to maintain a rigid connection with the existence of a probability distribution function is a well-known problem in nonequilibrium statistical mechanics. It is often the case that, when this connection is lost (failure of reality or positivity conditions), instability soon follows. In our case, the violation of unitarity is manifest in that positive expressions such as $`r^2/N`$ can turn out to be negative. Note that since both the Hartree and leading order $`1/N`$ approximations are variational in nature, they can never violate unitarity. In Fig. 8, we show the blow up or failure time for the next-to-leading order large $`N`$ approximation as a function of $`N`$; the failure time being defined as the time at which $`r^2/N`$ becomes negative. It is interesting to note that for values of $`N`$ near the critical value $`N_\text{c}=18.60`$ where the effective potential extends down to the origin, the failure time starts to increase rapidly (consistent with our interpretation of the connection of the static effective potential to an associated quantum state). For values of $`N`$ greater than 21, we could not find failure for $`t<150`$. Thus, it may well be that above a certain value of $`N`$, the unitarity violation is pushed out to times of no practical significance, or may even disappear altogether.
We discussed previously how to choose initial conditions so as to reach the large $`N`$ limit in a controlled manner. Starting with these initial conditions, we show in Fig. 9 results from numerical solution of the Schrödinger equation for the time evolution of $`r^2/N`$. In the strict large $`N`$ limit one expects pure harmonic oscillations about the minimum of the infinite-$`N`$ effective potential. For any finite value of $`N`$, however, at sufficiently large times the pure harmonic motion is overcome by nonlinear effects, and interesting behavior is found, as shown in Fig. 10, indicating the presence of a very non-Gaussian wave function. The ability to capture this long time behavior is an important test of $`1/N`$ methods.
In order to test whether Hartree or the next-to-leading order approximation incorporate nonlinearities correctly so as to capture the late time modulation behavior, we ran a comparison against the numerical results for $`N=21`$, the results being displayed in Fig. (11). It is clear that both approximations do not give satisfactory results. This provides additional motivation for the development of alternative $`1/N`$ expansions which would incorporate selective resummations in order to reduce the coefficient of the error term at late times.
### B Shifted Gaussian Initial Conditions
We now discuss the time evolution of a quantum state having an initial wave function given by Eq. (114). For this problem, because of the lack of symmetry, exact solutions were only obtained for $`N2`$. For $`N=1`$, depending on whether the energy is above or below the barrier height, one observes either slow tunneling with rapid oscillations in one well, or slower oscillation in the complete range. At these low values of $`N`$, the large-$`N`$ expansion breaks down quickly, as in the quantum roll, but even here at $`N=1`$, the $`1/N`$ corrections improve the short time accuracy of $`q(t)`$.
A more relevant comparison is to consider larger values of $`N`$ at which the approximations have a better chance of capturing the exact behavior. In the next figure, we compare the Hartree with the leading and next-to-leading order large-$`N`$ approximation for $`q(t)`$, at $`N=50`$. Fig. 12 displays the results for a run with $`E>E_b`$ using the equal time Green’s function approximation (see Sec. V) method for obtaining the Hartree results. Unlike the situation for the roll initial condition where the Hartree and next-to-leading order large $`N`$ results are not dramatically different, here the qualitative behavior is quite dissimilar (whereas the Hartree and leading order results are in fact very close).
## X Conclusions
Testing the $`1/N`$ approximation in quantum mechanics has already enabled us to arrive at some useful conclusions. In order to interpret our results, it is important to keep in mind that $`1/N`$ approximations are a form of resummed perturbation theory and are therefore only valid at weak coupling. Thus for couplings of order unity, it is unrealistic to expect the approximation to give good results for small values of $`N`$. Our results have shown that at sufficiently large $`N`$ the next-to-leading order approximation is a clear improvement over the leading order approximation, however, at late times this approximation (as well as Hartree) fails to capture the nonlinear effects that lead to nontrivial amplitude modulation of the radial oscillations in the quantum roll problem.
We have noted the presence of a finite-time breakdown in the evolution given by the next-to-leading order approximation. This result is related to the fact that the large $`N`$ expansion for the expectation values does not necessarily correspond to a positive semi-definite density matrix when truncated at any finite order in $`1/N`$. (At lowest order, the $`1/N`$ approximation is equivalent to a Gaussian variational ansatz for the density matrix and does not have this problem.) This last aspect is already clear even in static situations such as the lack of a real effective potential for all $`N`$ in the next-to-leading order approximation. This type of finite-time breakdown induced by unitarity/positivity violation has also been noted in simulations of quantum systems where the coupled equal times Green’s function approach was truncated at fourth order or where high order cumulant expansion methods were used .
Two aspects of this breakdown deserve further mention: First, the time at which breakdown occurs appears to be strongly connected with the behavior of the effective potential. For values of $`N`$ not very much bigger than the critical value $`N_c`$ (beyond which the effective potential exits over the entire range of $`y`$), the breakdown time increases extremely steeply and may even be pushed to times long enough to be no longer an obstacle to practical calculations ( this still needs to be demonstrated). Second, it is important to point out that avoiding the breakdown via a partial resummation does not automatically guarantee better late time accuracy (or convergence) since such a scheme is also only next-to-leading order accurate. However, it will help in the sense that one may carry out simulations at smaller values of $`N`$, thus making it easier to compare against the late-time numerical solutions of the corresponding Schrödinger equation.
One possible way of correcting the problem of a manifestly positive operator such as $`r^2`$ becoming negative is to solve for the full Green’s function $`𝒢_{ij}(t,t^{})`$:
$`𝒢_{ij}(t,t^{})=G_{ij}(t,t^{})`$ (139)
$`{\displaystyle \underset{k,l}{}}{\displaystyle _𝒞}dt_1{\displaystyle _𝒞}dt_2G_{ik}(t,t_1)\mathrm{\Sigma }_{kl}(t_1,t_2)𝒢_{lj}(t_2,t^{}),`$ (140)
rather than the next-to-leading order one as in (21). This equation is the exact equation one obtains by varying the effective action and it contains terms of all orders in $`1/N`$ (thus, strictly speaking, one is no longer truncating at some fixed order).
However, just making this correction does not increase the time period during which the approximation is accurate. In order to extend the accuracy of the $`1/N`$ approximation to late times, it appears necessary to use a more robust approximation based on the Schwinger-Dyson equations. Several approximations of this sort are possible, which may both cure the positivity problem as well as lead to accurate results at late times. These will be discussed separately .
## XI Acknowledgements
The authors acknowledge helpful conversations with Luis Bettencourt, Yuval Kluger, and Emil Mottola. S.H. acknowledges stimulating discussions with Larry Yaffe. The work of B.M. and J.F.D. at UNH is supported in part by the U.S. Department of Energy under grant DE-FG02-88ER40410. B.M. and J.F.D thank the Theory Group (T-8) at LANL and the Institute for Nuclear Theory at the University of Washington for hospitality during the course of this work. F.C. would like to thank the Physics Department at Yale University for their hospitality where some of this research was carried out. |
warning/0003/quant-ph0003025.html | ar5iv | text | # Λ’s, V’s, and optimal cloning with stimulated emission
## I Introduction
An ideal quantum cloning machine is a device that produces an arbitrary number of perfect copies of a given (unknown) quantum system. Such a device would allow the exact determination of the quantum state of a system. It has been shown that such a device would violate the linearity of quantum mechanics and also relativistic locality because it would make superluminal communication possible .
Non-perfect copying, though, can be realized in quantum mechanics. Since the seminal paper of Bužek and Hillery , quantum cloning has been extensively studied theoretically. Bruss et al. derived bounds on the possible fidelity of quantum cloners, Gisin and Massar, and Bužek and Hillery discovered optimal universal cloning transformations, and finally Werner, and Keyl and Werner discussed optimal universal cloning in great generality.
While optimal cloning was previously discussed in terms of quantum networks, in a recent paper some of the authors have shown that optimal universal cloning can be comparatively easily realized via stimulated emission . In this scheme the general qubit to be cloned is represented by the polarization state of a photon. When cloning is realized via stimulated emission, the fidelity of the clones is limited by the unavoidable presence of spontaneous emission. It was shown that the bounds on the fidelity given by the above-mentioned fundamental principles can nevertheless be saturated.
In Sec. II of the present paper we present a scheme for optimal universal cloning based on stimulated emission in three-level systems of the Lambda type. Our main analytic tool is the formal equivalence between systems of Lambda atoms and coupled harmonic oscillators, which is established in Sec. II B. In Sec. II C, this equivalence is used to analyze the structure of the transformations realized in detail and to prove their optimality. More specifically, we explicitly demonstrate their equivalence to the optimal cloning transformations for qubits discovered before . In particular, it will become clear that the atomic states play the double role of photon source and ancilla, and that the universal NOT operation is realized in the ancilla states. In the same way, we show that the down-conversion cloner presented in is obtained from the present schemes as a limiting case. In Sec. III we demonstrate that optimal cloning can also be achieved with pairwise entangled V atoms, using an interesting equivalence between the two systems. In Sec. IV we discuss the physical differences that exist between our stimulated emission cloners and the qubit cloners considered previously, and we give an explicit proof that the bounds derived for qubit cloning indeed apply to our situation as well. Section V gives our conclusions.
## II Cloning via stimulated emission in Lambda-atoms
The general principles of universal cloning via stimulated emission are the following. Consider an inverted medium that can spontaneously emit photons of any polarization with the same probability. If a photon (or several) of a given polarization interacts with such a medium, it stimulates the emission of photons of the same polarization. In the final photonic state there will be a majority of photons polarized parallel to the incoming photon, while some photons will be in the orthogonal polarization due to spontaneous emission. In this way the photons in the final state can be considered as clones of the original incoming photon, where the fidelity of the clones is given by the relative frequency of photons of the correct polarization in the final state.
The inverted medium that we will use as a cloning device consists of an ensemble of Lambda-atoms. These are three-level systems that have two degenerate ground states $`|g_1`$ and $`|g_2`$ and an excited level $`|e`$. The ground states are coupled to the excited state by two modes of the electromagnetic field, $`a_1`$ and $`a_2`$, respectively. These two modes define the Hilbert space of our qubit to be cloned, i.e. we want to clone general superposition states $`(\alpha a_1^{}+\beta a_2^{})|0,0=\alpha |1,0+\beta |0,1`$. We can think of $`a_1`$ and $`a_2`$ as being orthogonal polarizations of one photon with a specific frequency, but we do not have to restrict ourselves to such a specific example, in fact we can think about other systems and other degrees of freedom, as long as they are described by the same formalism, e.g. $`a_1`$ and $`a_2`$ could also refer to the center-of-mass motion (phonons) in an ion trap. In the interaction picture, after the usual dipole and rotating wave approximations, the interaction Hamiltonian between field and atoms has the following form:
$$_i=\gamma (a_1\underset{k=1}{\overset{N}{}}|e^kg_1^k|+a_2\underset{k=1}{\overset{N}{}}|e^kg_2^k|)+h.c.=\gamma (a_1\underset{k=1}{\overset{N}{}}\sigma _{+,1}^k+a_2\underset{k=1}{\overset{N}{}}\sigma _{+,2}^k)+h.c.$$
(1)
The index $`k`$ refers to the $`k`$-th atom. Note that in (1) the atoms couple to only one single spatial mode of the electromagnetic field. In particular this means that spontaneous emission into all other modes is neglected. Situations where this is a good approximation can now be achieved in cavity QED . We also assume that the coupling constant $`\gamma `$ is the same for all atoms, which in a cavity QED setting means that they have to be in equivalent positions relative to the cavity mode. Trapping of atoms inside a cavity has recently been achieved . Finally note that our Hamiltonian has no spatial dependence, which means that the effect of the field on the motion of the atoms is neglected, their spatial wavefunction is assumed to be unchanged .
The Hamiltonian (1) is invariant under simultaneous unitary transformations of the vectors $`(a_1,a_2)`$ and $`(|g_1,|g_2)`$ with the same matrix $`U`$. If one furthermore chooses an initial state of the atoms that has the same invariance, then the system behaves equivalently for all incoming photon polarizations, i.e. universal cloning is achieved. This can be seen in the following way. Consider an incident photon in a general superposition state $`(\alpha a_1^{}+\beta a_2^{})|0,0`$. Together with the orthogonal one-photon state this defines a new basis in polarization space, which is connected to the original one by a unitary transformation. If the atomic states are now rewritten in the basis that is connected to the original one by the same unitary transformation, then under the above assumptions the interaction Hamiltonian and initial state of the atoms look exactly the same as in the original basis. The initial state where all atoms are excited to $`|e`$ has the required invariance: it is completely unaffected by the above-mentioned transformations.
We can therefore, without loss of generality, restrict ourselves to the cloning of photons in mode $`a_1`$. We consider an initial state
$$|\mathrm{\Psi }_{in}=_{k=1}^N|e^k\frac{(a_1^{})^m}{\sqrt{m!}}|0,0,$$
(2)
i.e. we are starting with $`m`$ photons of a given polarization, and we want to produce a certain (larger) number $`n`$ of clones.
### A The simplest case
For illustrative purposes let us first consider the simplest case of one Lambda-atom and one photon polarized in direction $`1`$:
$$|\mathrm{\Psi }_{in}=|ea_1^{}|0,0=|e|1,0=:|_0$$
(3)
To study the time development, we expand the evolution operator $`e^{it}`$ into a Taylor series and determine the action of powers of $``$ on the state $`|\mathrm{\Psi }_{in}`$.
$`|\mathrm{\Psi }_{in}`$ $`=`$ $`\gamma (|g_1a_1^{}|1,0+|g_2a_2^{}|1,0)=\gamma \sqrt{3}{\displaystyle \frac{(\sqrt{2}|g_1|2,0+|g_2|1,1)}{\sqrt{3}}}=:\gamma \sqrt{3}|_1`$ (4)
$`^2|\mathrm{\Psi }_{in}`$ $`=`$ $`\gamma ^2(|ea_1\sqrt{2}|2,0+|ea_2|1,1)=3\gamma ^2|e|1,0=3\gamma ^2|_0`$ (5)
$`\mathrm{}`$ (6)
The result is
$`e^{it}|\mathrm{\Psi }_{in}`$ $`=`$ $`\mathrm{cos}(\gamma \sqrt{3}t)|e|1,0i\mathrm{sin}(\gamma \sqrt{3}t)(\sqrt{{\displaystyle \frac{2}{3}}}|g_1|2,0+\sqrt{{\displaystyle \frac{1}{3}}}|g_2|1,1)`$ (8)
$`=`$ $`\mathrm{cos}(\gamma \sqrt{3}t)|_0i\mathrm{sin}(\gamma \sqrt{3}t)|_1`$ (9)
$`|_0`$ and $`|_1`$ denote the states of the system atom-photons that lie in the subspace with $`1`$ and $`2`$ photons respectively. $`|_0`$ is in the subspace where no cloning has taken place and $`|_1`$ in the one where one additional photon has been emitted, so that the two photons can now be viewed as clones with a certain fidelity. This way of labeling the states will turn out to be convenient below. The probability that the system acts as a cloner is $`p(1)=\mathrm{sin}^2(\gamma \sqrt{3}t)`$. The fidelity $`F_1`$ of the cloning procedure can be defined as the relative frequency of photons in the correct polarization mode in the final state $`|_1`$ (cf. Sec. IV). One finds
$$F_1=\frac{2}{3}1+\frac{1}{3}\frac{1}{2}=\frac{5}{6},$$
(10)
which is exactly the optimal fidelity for a 1-to-2 cloner . Actually, the state
$$|_1=\sqrt{\frac{2}{3}}|2,0|g_1+\sqrt{\frac{1}{3}}|1,1|g_2$$
(11)
is exactly equivalent to the three-qubit state
$$\sqrt{\frac{2}{3}}|11|+\sqrt{\frac{1}{3}}\left(\frac{1}{\sqrt{2}}(|01+|10)\right)|$$
(12)
produced by the Bužek-Hillery cloner, see , Eq. (3.29b). The equivalence is established, if the photonic states in Eq. (11) are identified with the corresponding symmetrized two-qubit states (both photons in mode 1 means both qubits in state $`|1`$, one photon in each mode means one qubit in state $`|1`$, one in state $`|0`$) in Eq. (12), while the atomic states $`|g_1`$ and $`|g_2`$ are identified with the states $`|`$ and $`|`$ of the ancillary qubit. This is another way of proving the optimality of Eq. (11). Note that in our case the universality follows directly from the symmetry of initial state and Hamiltonian, as explained above. In the following we show that a similar equivalence holds between our cloning scheme and the Gisin-Massar cloners in the completely general case (arbitrary numbers of photons and atoms).
### B Equivalence to coupled harmonic oscillators
We now turn to the discussion of the general case, i.e. we consider the initial state (2). We are going to show the equivalence of our system defined by (1) and (2) to a system of coupled harmonic oscillators. First note that both the initial state (2) and the Hamiltonian (1) are invariant under permutations of the atoms, which implies that the state vector of the system will always be completely symmetric. Furthermore the Hamiltonian (1) can be rewritten as
$$=\gamma \left(a_1J_{+,1}+a_2J_{+,2}\right)+h.c.$$
(13)
in terms of “total angular momentum” operators
$$J_{+,r}=\underset{k=1}{\overset{N}{}}\sigma _{+,r}^k=\underset{k=1}{\overset{N}{}}|e^kg_r^k|(r=1,2),$$
(14)
By the above considerations one is led to use a Schwinger type representation for the angular momentum operators:
$$J_{+,r}=b_rc^{}(r=1,2),$$
(15)
where $`c^{}`$ is a harmonic oscillator operator creating “$`e`$” type excitations, while $`b_1`$ destroys “$`g_1`$” excitations. Note that $`J_{+,1}`$ and $`J_{+,2}`$ share the operator $`c^{}`$ because both ground levels $`g_1`$ and $`g_2`$ are connected to the same upper level $`e`$ by the Hamiltonian (1), and correspondingly for the Hermitian conjugates. In terms of these operators, (1) acquires the form
$$_{osc}=\gamma (a_1b_1+a_2b_2)c^{}+h.c.,$$
(16)
while the initial state (2) is now given by
$$|\psi _i=\frac{(a_1^{})^m}{\sqrt{m!}}\frac{(c^{})^N}{\sqrt{N!}}|0=|m_{a1},0_{a2},0_{b1},0_{b2},N_c|m,0,0,0,N.$$
(17)
Actually, for reasons that will become apparent below, it is slightly more convenient for our purposes to use the following Hamiltonian instead of (16):
$$=\gamma (a_1b_2a_2b_1)c^{}+h.c.,$$
(18)
which can be obtained from (16) by a simple unitary transformation in mode $`b`$, corresponding to a simple redefinition of the atomic states in (1). This is the Hamiltonian that is going to be used in the rest of this paper. The invariance properties of (18) are linked to those of (1) or equivalently (16) discussed above: (18) is invariant under simultaneous identical SU(2) transformations in modes $`a`$ and $`b`$ (because the determinant of such a transformation is equal to unity), while a phase transformation in either mode can be absorbed into $`\gamma `$. This ensures the universality of the cloning procedure.
We are now dealing with five harmonic oscillator modes defined by the operators $`c,b_1,b_2,a_1`$, and $`a_2`$. Action of (18) on (17) generates Fock basis states of the general form
$$|(m+j)_{a1},i_{a2},i_{b1},j_{b2},(Nij)_c=|m+j,i_{photons}|i,j,Nij_{atoms}.$$
(19)
Remember that $`a_1`$ is now coupled to $`b_2`$ etc. Expressed in terms of individual atoms, $`|i,j,Nij_{atoms}`$ is the completely symmetrized state with $`i`$ atoms in level $`g_1`$, $`j`$ atoms in level $`g_2`$, and $`Nij`$ atoms in level $`e`$. The correctness of (15) can be checked by explicit application of left hand side and right hand side to such a general state, written in terms of the individual atoms and in terms of harmonic oscillator eigenstates respectively (see Appendix A).
Note that the use of the Schwinger representation is only convenient because the initial state of the atomic system in (2) is completely symmetric under permutation of the atoms.
Studying the Hamiltonian in the form (18) instead of (1) is helpful in several respects. The number of atoms $`N`$ that is explicit in the Hamiltonian (1) now appears only as a part of the initial conditions of our system, which makes it easy to treat the general case of N atoms in one go. We will do this in the next subsection.
Furthermore, the connection to cloning by parametric down-conversion (PDC) as proposed in is now obvious. The Hamiltonian (18) can also be seen as a Hamiltonian for down-conversion with a quantized pump-mode described by the operator $`c`$, while $`a_r`$ and $`b_r`$ are the signal and idler modes respectively, where $`r`$ labels the polarization degree of freedom. There is only one difference between (18) and the Hamiltonian used in (see Eq. 6 of that reference): in the operator $`c`$ of (18) is replaced by a c-number. In the context of down-conversion, this corresponds to the limit of a classical pump field. Thus the PDC scheme, which was shown to achieve optimal universal cloning in , is obtained as a limiting case from the schemes discussed here.
In passing we note that the above dynamical equivalence generalizes to atoms with more than $`2`$ ground-states $`|g_n`$ that are coupled each to a different degree of freedom of photons $`a_n`$. By similar arguments a system of $`N`$ identical atoms with $`r`$ ground states $`\{|g_1,\mathrm{},|g_r\}`$ governed by a Hamiltonian
$$^r=\gamma \underset{k=1}{\overset{N}{}}\underset{n=1}{\overset{r}{}}|e^kg_n^k|a_n+h.c.$$
(20)
is equivalent to a system of $`r+1`$ coupled harmonic oscillators with lowering operators $`c`$ and $`b_1,\mathrm{},b_r`$ governed by the interaction Hamiltonian
$$_{osc}^r=\gamma \underset{n=1}{\overset{r}{}}cb_n^{}a_n^{}+h.c.$$
(21)
### C Cloning of $`m`$ photons with $`N`$ Lambda-atoms: Proof of optimality
We are now going to show that the system defined by (17) and (18) indeed realizes optimal cloning for arbitrary $`N`$ and $`m`$. The idea of the proof is the following. After evolution in time the system that started with a certain photon number $`m`$ will be in a superposition of states with different total photon numbers, where total means counting photons in mode $`a_1`$ and $`a_2`$, i.e. both “good” and “bad” copies. We will show that the general form of the state vector after a time interval $`t`$ is
$$|\mathrm{\Psi }(t)=e^{it}|\mathrm{\Psi }_{in}=\underset{l=0}{\overset{N}{}}f_l(t)|_l,$$
(22)
where $`l`$ denotes the number of additional photons that have been emitted and
$$|_l:=\left(\genfrac{}{}{0pt}{}{m+l+1}{l}\right)^{\frac{1}{2}}\underset{i=0}{\overset{l}{}}(1)^i\sqrt{\left(\genfrac{}{}{0pt}{}{m+li}{m}\right)}|(m+li)_{a1},i_{a2},i_{b1},(li)_{b2},(Nl)_c.$$
(23)
Note that the number of photons can never become smaller than $`m`$ since all the atoms start out in the excited state. $`|_l`$ is a normalized state of the system with $`m+l`$ photons in total. To see that $`|_l`$ is properly normalized note that $`_{i=0}^l\left(\genfrac{}{}{0pt}{}{m+i}{m}\right)=\left(\genfrac{}{}{0pt}{}{m+l+1}{l}\right)`$.
The states $`|_l`$ are formally identical to the states obtained in , which have been shown to realize optimal universal cloning and the optimal universal NOT simultaneously. The ideal universal NOT is an operation that produces the orthogonal complement of an arbitrary qubit. Like perfect cloning, it is prohibited by quantum mechanics. The transformation in links universal cloning and universal NOT (anti-cloning): the ancilla qubits of the cloning transformation are the anti-clones. In our case, the clones are the photons in the $`a`$-modes and the anti-clones are the atoms in the $`b`$-modes (atomic ground states $`g_1`$ and $`g_2`$). From the Hamiltonian (18) and (23) it is clear that for every “good” emitted photon-clone (in mode $`a_1`$) there is an excitation in mode $`b_2`$ which corresponds to an anti-clone (atomic ground state $`|g_2`$). The only difference to the states in is the presence of the fifth harmonic oscillator mode $`c`$, describing the “e” type excitations, which counts the total number of clones that have been produced (equal to the number of atoms having gone to one of the ground states) and doesn’t affect any of the conclusions.
A distinguishing feature of our cloner is that the output state (23) is a superposition of states with different total numbers of clones. Cloning with a certain fixed number of produced copies can be realized by measuring the number of atoms in the excited state $`|e`$ (corresponding to mode $`c`$) and post-selection.
To see that the $`|_l`$ are indeed the output of an optimal cloner, let us calculate the fidelity of the cloning, given by the mean relative frequency of photons in the correct mode ($`a_1`$ in our case). In the state $`|m+li,i_{photons}`$ the relative frequency of correct photons is $`(m+li)/(m+l)`$. Therefore
$`F_l`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{m+l+1}{l}}\right)^1{\displaystyle \underset{i=0}{\overset{l}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{m+li}{m}}\right){\displaystyle \frac{m+li}{m+l}}={\displaystyle \frac{m(m+2)+l(m+1)}{(m+l)(m+2)}}`$ (24)
which corresponds to the fidelity of an optimal universal $`mm+l`$ cloner . Note again that the universality in our case follows from the symmetry of the Hamiltonian and the initial atomic state.
To prove that the system is indeed always in a superposition of the states $`|_l`$ as in Eq. (23) we use induction: The initial state of the system is $`|\mathrm{\Psi }_{in}=|_0`$. Now we will show that if $`|\mathrm{\Phi }`$ is a superposition of states $`|_l`$ then $`|\mathrm{\Phi }`$ is so, too. Then, since $`|\mathrm{\Psi }(t)=e^{it}|\mathrm{\Psi }_{in}=_p\frac{(it)^p}{p!}|\mathrm{\Psi }_{in}`$ this implies that $`|\mathrm{\Psi }(t)`$ will be a superposition of $`|_l`$. Explicit calculation shows that
$`|_l`$ $`=`$ $`\gamma (\sqrt{(l+1)(Nl)(m+l+2)}|_{l+1}+\sqrt{l(Nl+1)(m+l+1)}|_{l1})1l<N`$ (25)
$`|_0`$ $`=`$ $`\gamma \sqrt{N(m+2)}|_1`$ (26)
$`|_N`$ $`=`$ $`\gamma \sqrt{N(m+N+1)}|_{N1}`$ (27)
which completes the proof.
Note that the form of the coefficients $`f_l(t)`$ didn’t play any role in our proof. Actually, the $`f_l`$ are in general hard to determine exactly. Solutions have been found in limiting cases. For the limit of a classical pump field ($`c`$ replaced by a c-number), the solution can be found by standard methods and is given in . The solution in the case of large incoming photon numbers ($`mN`$) is presented in Appendix B.
Let us pause here for a moment and summarize what we have found. Our system consisting of an ensemble of Lambda-atoms in the excited state is indeed equivalent to a superposition of optimal cloning machines a la Bužek-Hillery or Gisin-Massar, producing various numbers of clones. The atoms play the double role of photon source and of ancilla, the atomic ground states can be identified with the ancilla states in the qubit cloners. As for the corresponding qubit cloners, those ancillary atoms can also be seen as the output of a universal NOT gate. On the other hand, the atoms that end up in the excited state provide information about the number of clones that has actually been produced. This can be used to realize cloning with a fixed number of output clones by post-selection.
## III The equivalence between pairs of V-atoms and Lambda-atoms
In this section we present an alternative (but similar) way of realizing optimal universal cloning that uses entangled pairs of V-atoms instead of Lambda atoms. We prove optimality by showing that the system can be exactly mapped onto the system with Lambda atoms that we discussed above.
The two degenerate upper levels of each V-atom, $`|e_1`$ and $`|e_2`$, are coupled to the ground state $`|g`$ via the two orthogonal modes $`a_1`$ and $`a_2`$ respectively. The Hamiltonian describing the interaction of atom and field is:
$$_V=\gamma (a_1^{}\underset{k=1}{\overset{N}{}}|g^ke_1^k|+a_2^{}\underset{k=1}{\overset{N}{}}|g^ke_2^k|)+h.c.=\gamma (a_1^{}\underset{k=1}{\overset{N}{}}\sigma _{,1}^k+a_2^{}\underset{k=1}{\overset{N}{}}\sigma _{,2}^k)+h.c.$$
(28)
It arises from similar assumptions as (1). In contrast to before we now choose an entangled state of the atoms as the initial state. This is motivated by the fact that the initial atomic state has to be a singlet under polarization transformations in order for our cloning device to be again universal.
Let us first examine the simplest case of two entangled V-atoms, $`A`$ and $`B`$, and one incoming photon. The initial state of the system is
$$|\mathrm{\Psi }_{in}=\frac{1}{\sqrt{2}}(|e_1^Ae_2^B|e_2^Ae_1^B)|1,0$$
(29)
Developing the time evolution operator $`e^{i_𝒱t}`$ into a power series, one finds easily:
$`e^{i_𝒱t}|\mathrm{\Psi }_{in}`$ $`=`$ $`\mathrm{cos}(\gamma \sqrt{3}t){\displaystyle \frac{|e_1^Ae_2^B|e_2^Ae_1^B}{\sqrt{2}}}|1,0`$ (31)
$`i\mathrm{sin}(\gamma \sqrt{3}t)\left(\sqrt{{\displaystyle \frac{2}{3}}}{\displaystyle \frac{|g^Ae_2^B|e_2^Ag^B}{\sqrt{2}}}|2,0+\sqrt{{\displaystyle \frac{1}{3}}}{\displaystyle \frac{|e_1^Ag^B|g^Ae_1^B}{\sqrt{2}}}|1,1\right)`$
With the substitution
$`{\displaystyle \frac{|e_1^Ae_2^B|e_2^Ae_1^B}{\sqrt{2}}}|\stackrel{~}{e}`$ (32)
$`{\displaystyle \frac{|g^Ae_2^B|e_2^Ag^B}{\sqrt{2}}}|\stackrel{~}{g_1}`$ (33)
$`{\displaystyle \frac{|e_1^Ag^B|g^Ae_1^B}{\sqrt{2}}}|\stackrel{~}{g_2}`$ (34)
the state (31) has exactly the same form as the corresponding state (8) for one Lambda-atom, which implies that it also implements optimal universal $`12`$ cloning.
Actually, the correspondence goes much further. Consider an initial atomic state consisting of $`N`$ pairs of V-atoms, where each pair is in a singlet state:
$$|\psi _i=_{k=1}^N|\stackrel{~}{e}^k$$
(35)
with $`|\stackrel{~}{e}`$ as defined in (32).
It is easy to see that the action of the Hamiltonian (28) on each pair only generates one of the three antisymmetric atomic states in Eq. (32). Because of the invariance of the Hamiltonian under permutations, and in particular under the exchange of two atoms belonging to the same pair, transitions between states with different symmetry properties are impossible. In fact, with the identification (32) the Hamiltonian (28) has exactly the same form as the Hamiltonian for Lambda-atoms (1). The analysis made for Lambda atoms in Section II now goes through unchanged and we obtain the same cloning properties of a system of pairwise entangled V-atoms as we had before for Lambda-atoms, i.e. we have found another way of realizing optimal universal cloning. Although this scheme would without doubt be more difficult to realize experimentally, we believe that the underlying equivalence between the two systems is interesting and may be useful in other contexts as well.
## IV Cloning of photons versus cloning of qubits
In this section we are going to discuss the physical differences that exist in spite of the formal equivalence proven above between our photon cloners based on stimulated emission and the qubit cloners as usually considered . In particular, we will show that the claim that optimal cloning is realized by our devices is justified in spite of these differences.
In most of the previous work cloning was discussed in terms of quantum networks. In general, the situation considered in these papers is the following: one has a certain number of qubits that are localized in different positions, which makes them perfectly distinguishable. At the beginning, some of those qubits are the systems to be cloned, the others play the role of ancillas. After the cloning procedure, which consists of several joint operations on the qubits that can be expressed in terms of quantum gates, some of the qubits are the clones, the rest are ancillas, which for a specific form of the optimal cloning transformation can also be seen as outputs of the universal NOT operation. As a consequence of localization, it is possible to address individual clones.
In our stimulated emission cloners, the situation is different. All input systems (photons) are in the same spatial mode (called mode $`a`$ in this paper), and, even more importantly, all clones are produced into that mode. Note that this is completely unavoidable if stimulated emission is to be used. One can say that this is the price one has to pay for the great conceptual simplicity of the cloning procedure itself.
However, having all clones in the same spatial mode is not necessarily an important disadvantage. For example, if perfect cloning of that kind were possible, one could still determine the polarization of the original photon to arbitrary precision by performing measurements on the clones. This would still make superluminal communication possible . If one wants to distribute the clones to different locations, this can for example be achieved using an array of beam splitters. However, this does not lead to a situation where one can be sure to have exactly one photon in each mode. If one wants to have at most one photon in each mode, the array has to have many more output modes than there are photons.
Another distinguishing feature of our cloners compared to the usual qubit cloners is the fact that the same procedure is used to produce different numbers of clones. While in the qubit case the network to be used depends on the number of desired clones, in our case the final state is a superposition of states with different numbers of clones. Of course, the average number of clones produced depends on the number of atoms present in the system and the interaction time. As discussed in Sec. II cloning with a fixed number of output clones can be achieved by post-selection based on a measurement of the number of excited atoms in the final state.
The formal equivalence between the qubit cloners and our one-mode cloners can arise because the output state produced by the optimal qubit cloners is completely symmetric under the exchange of clones . Because of the bosonic nature of the photons there is a one-to-one-mapping between completely symmetric qubit states and photonic states. For a completely symmetric qubit state the two concepts of relative frequency of qubits in the “correct” basis state and of single-particle fidelity are equivalent. This can be seen in the following way. Let $`|\psi `$ denote the state that is to be copied. Then the usual definition of the (single-particle) cloning fidelity is
$$F=\psi |\rho _{red}|\psi ,$$
(36)
where $`\rho _{red}`$ is the reduced density matrix of one of the clones, say the first one, i.e.
$$\rho _{red}=\text{Tr}_{2,3,\mathrm{},N}\left[\rho \right]$$
(37)
Then $`F`$ can also be expressed as
$$F=\text{Tr}\left[\rho |\psi \psi |_1I_2\mathrm{}I_N\right].$$
(38)
On the other hand, the relative frequency of qubits in the state $`|\psi `$ can be written as
$$\frac{1}{N}\text{Tr}\left[\rho \left(|\psi \psi |_1I_2\mathrm{}I_N+I_1|\psi \psi |_2\mathrm{}I_N+\mathrm{}+I_1\mathrm{}|\psi \psi |_N\right)\right].$$
(39)
If $`\rho `$ is invariant under exchange of any two clones, it is obvious that (39) is equal to (38), i.e. for symmetric cloners the two concepts are completely equivalent. This justifies our definition of fidelity via the relative frequency in the case of photon cloning (cf. Sec. II).
Let us finally address the issue of optimality in the context of stimulated emission cloners. In this paper we have shown the formal equivalence of our scheme and the optimal schemes for qubit cloning. As a consequence, the fidelity of the clones saturates the bounds derived for the cloning of qubits. However, it is not entirely obvious that the bounds derived for the situation of distinct well-localized qubits also apply to our situation. Could one maybe achieve even higher fidelity in our one-mode case? The following argument shows that the bounds indeed apply in our situation as well, i.e. that photon cloning is not allowed to be better than qubit cloning.
Let us assume that we had a single-mode cloning machine that clones photons with a better fidelity than given by the bounds for qubits. Consequently, the relative frequency of “correct” photons has to exceed the bound for at least one value of the final total photon number $`M`$. This is obvious if $`M`$ has been fixed by post-selection. Otherwise the fidelity has to be defined as the average of the relative frequencies over all final total photon numbers. This average can only exceed the bound for qubits if the bound is violated for at least one particular value $`M`$ of the final photon number.
As a consequence, we have a universal map from the $`N`$-photon Hilbert space to the $`M`$-photon Hilbert space that achieves a relative frequency of correct photons in the final state that is higher than the qubit bound. But the existence of such a map is equivalent to the existence of a universal map from the totally symmetric $`N`$-qubit space to the totally symmetric $`M`$-qubit space with a single-particle fidelity equal to the relative frequency. The existence of the latter map is excluded by the theorems on cloning of qubits . This justifies our claim that the schemes presented in the previous sections realize optimal cloning of photons.
## V Conclusions and Outlook
In this paper we have shown that optimal universal cloning can be realized via stimulated emission in three-level systems. The permutation symmetry of the interaction allowed to map our system onto bosonic modes independent of the number of atoms used. Furthermore, we have found an equivalence between single Lambda atoms and entangled pairs of V systems, which might be fruitful in other contexts as well.
The connection between stimulated emission and optimal cloning is remarkable. Our results show that a task previously discussed in terms of rather complicated quantum networks can be realized in an elegant way using basic quantum systems and interactions. While it was clear from the beginning that perfect cloning is prohibited by fundamental principles, it is interesting to see how this impossibility arises in a concrete physical system. In our case, the physical process limiting the fidelity of the clones is spontaneous emission. It is fascinating that in this way spontaneous emission ensures that there cannot be any superluminal communication.
It might be interesting to investigate possible experimental realizations of our proposal, e.g. using a combination of cavity QED and Bose-Einstein condensates. This would potentially allow the creation of macroscopic numbers of clones.
Quantum cloners are often discussed in the context of eavesdropping in quantum cryptography. Currently all cryptography schemes rely on photons. Therefore devices based on the principles presented here could be useful to a future eavesdropper.
We would like to thank A. Zeilinger for stimulating discussions on various aspects of the present work, V. Bužek for pointing out that what we are doing can be seen as use of the Schwinger representation, and A. Karlsson for motivating us to discuss the physical peculiarities of our photon cloners. Furthermore we thank D. Bacon, Č. Brukner, I. Cirac, and S. Stenholm for useful discussions. This work was supported by the Austrian Science Foundation (FWF), project no. S6504 and F1506, by the U.S. ARO under DAAG55-98-1-0371 and by U.S. NSF DMS-9971169.
## A Schwinger representation
As noted above, the action of the Hamiltonian (1) on the initial state (2) only generates completely symmetric states of the atomic system. These states have the general form
$$\left(\genfrac{}{}{0pt}{}{N}{i,j}\right)^{1/2}\underset{\alpha }{}|g_1^{\alpha _1},g_1^{\alpha _2},\mathrm{},g_1^{\alpha _i},g_2^{\alpha _{i+1}},\mathrm{},g_2^{\alpha _{i+j}},e^{\alpha _{i+j+1}},\mathrm{},e^{\alpha _N}=:|i,j,Nij_{atoms}$$
(A1)
where the sum is over all arrangements $`\alpha `$ of the $`Nij`$ levels $`|e`$, the $`i`$ levels $`|g_1`$, and the $`j`$ levels $`|g_2`$ on the $`N`$ atoms, and $`\left(\genfrac{}{}{0pt}{}{N}{i,j}\right)=\frac{N!}{i!j!(Nij)!}`$ is the multinomial coefficient giving the number of such arrangements.
Now study the action of a typical term in the Hamiltonian (1) on the system whose state we will write as
$`|i,j,Nij_{atoms}|m+i,j_{photons}`$:
$`({\displaystyle \underset{k=1}{\overset{N}{}}}|g_1^ke^k|)a_1^{}|i,j,Nij_{atoms}|m+i,j_{photons}`$ (A2)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}|g_1^ke^k|\sqrt{{\displaystyle \frac{i!j!(Nij)!}{N!}}}{\displaystyle \underset{\alpha }{}}|g_1^{\alpha _1},\mathrm{},g_1^{\alpha _i},g_2^{\alpha _{i+1}},\mathrm{},e^{\alpha _N}a_1^{}|m+i,j_{field}`$ (A3)
$`=`$ $`(i+1)\sqrt{{\displaystyle \frac{i!j!(Nij)!}{N!}}}{\displaystyle \underset{\alpha }{}}|g_1^{\alpha _1},\mathrm{},g_1^{\alpha _i},g_1^{\alpha _{i+1}},g_2^{\alpha _{i+2}},\mathrm{},e^{\alpha _N}a_1^{}|m+i,j_{field}`$ (A4)
$`=`$ $`\sqrt{i+1}\sqrt{Nij}\sqrt{{\displaystyle \frac{(i+1)!j!(Nij1)!}{N!}}}{\displaystyle \underset{\alpha }{}}|g_1^{\alpha _1},\mathrm{},g_1^{\alpha _i},g_1^{\alpha _{i+1}},g_2^{\alpha _{i+2}},\mathrm{},e^{\alpha _N}a_1^{}|m+i,j_{field}`$ (A5)
$`=`$ $`\sqrt{i+1}\sqrt{Nij}|i+1,j,Nij1_{atoms}a_1^{}|m+i,j_{photons}`$ (A6)
Here the factor $`(i+1)`$ arises from the number of different configurations that a given arrangement $`\alpha `$ can be reached by. This shows that this term acts exactly like a term $`a_1^{}b_1^{}c`$. Similar calculations can be made for the other terms in the Hamiltonian. Together, they justify the Schwinger representation (15).
## B Limit of large photon number
Here we determine the coefficients $`f_l(t)`$ of Eq. (LABEL:inFs) in the limit of large $`m`$ ($`m>>N`$, many incoming photons, small number of atoms). For that case, the recursion (27) becomes
$`|_l`$ $`=`$ $`\gamma \sqrt{m}(\sqrt{(l+1)(Nl)}|_{l+1}+\sqrt{l(Nl+1)}|_{l1})1l<N`$ (B1)
$`|_0`$ $`=`$ $`\gamma \sqrt{m}\sqrt{N}|_1`$ (B2)
$`|_N`$ $`=`$ $`\gamma \sqrt{m}\sqrt{N}|_{N1}`$ (B3)
It is possible to diagonalize the “transfer” matrix $`A`$ acting on the vector $`(f_0,\mathrm{},f_N)`$ that corresponds to the action of $``$ on $`|\mathrm{\Psi }=_{l=0}^Nf_l|_l`$: $`A_{l,l+1}=\gamma \sqrt{m}\sqrt{(l+1)(Nl)}=A_{l+1,l}`$. This allows to exponentiate $`A`$ and to determine the final state of the system after a time $`t`$:
$$|\mathrm{\Psi }(t)=\underset{l=0}{\overset{N}{}}(i)^l\sqrt{\left(\genfrac{}{}{0pt}{}{N}{l}\right)}\mathrm{cos}^{Nl}(\gamma \sqrt{m}t)\mathrm{sin}^l(\gamma \sqrt{m}t)|_l$$
(B4)
Differentiating (B4) and using (B3) one can show that this state fulfills Schroedinger’s equation with the correct initial condition.
In this big-$`m`$-limit the probability to observe the system as an $`mm+l`$ cloner (i.e. the probability that $`l`$ additional photons are emitted) is
$$p(l)=\left(\genfrac{}{}{0pt}{}{N}{l}\right)\mathrm{cos}^{2(Nl)}(\gamma \sqrt{m}t)\mathrm{sin}^{2l}(\gamma \sqrt{m}t)$$
(B5)
This is a binomial distribution with a probability $`\mathrm{sin}^2(\gamma \sqrt{m}t)`$ for each atom to emit a photon. Setting $`N=1`$ or comparison with Eq. (8) shows that this is identical to the probability for the case of only one atom in the case of large $`m`$. This means that in this limit each atom interacts independently with the electromagnetic field, because the effect of the other atoms on the field is negligible. In the short-time limit $`p(l)=O(t^{2l})`$. Furthermore the expected average number of “clones” $`N_c=_{l=0}^Nlp(l)=N\mathrm{sin}^2(\gamma \sqrt{m}t)`$ oscillates with an $`m`$-dependent frequency. |
warning/0003/astro-ph0003186.html | ar5iv | text | # Accretion disc–stellar magnetosphere interaction: field line inflation and the effect on the spin-down torque.
## 1 Introduction
Accreting magnetised stars exist in the form of neutron stars in close binary systems and T Tauri stars with surrounding protostellar discs. In both cases an accretion disc is disrupted through interaction with the stellar magnetosphere and the accretion flow becomes channelled along field lines. Torques act between the star and the accreting material and they determine the spin history of the star and its equilibrium rotation period. A spin–up torque is produced by rapidly rotating material in the inner regions of the disc while a spin–down torque is produced by the more slowly rotating material in the outer parts. Calculation of the torques requires knowledge of the structure of the disc and the magnetosphere in which it is embedded.
Early work (\[Ghosh & Lamb (1979a)\]; \[Ghosh & Lamb (1979b)\]) suggested a direct dependence of the total torque on the mass accretion rate. However, recent observations of accreting neutron stars (\[Nelson et al. (1997)\]; \[Chakrabatry et al. (1997)\]) suggest that the torques can oscillate producing alternating spin–up and spin–down phases with no evidence of a consistent correlation between the torque and the mass accretion rate. ?) suggest that this may be due to a variable outer magnetosphere which causes a corresponding variation in the spin–down torque acting between the star and the disc. In this paper we investigate the structure of a magnetosphere corotating with the central star. We calculate the spin–down torque due to the star–disc interaction and we find that it is very sensitive to the amount of field line twisting that occurs through the coupling of the stellar field lines to the disc. The magnetospheric structure and the resulting torque are also sensitive to the applied boundary conditions. In the present work we neglect any torque that may arise from a stellar wind.
According to the Ghosh & Lamb model the magnetic field of the central star penetrates the accretion disc due to diffusion arising from the growth of various instabilities. In the outer parts of the disc, viscous stresses are more important than magnetic stresses so that the effects of the magnetic field on the disc structure can be ignored. In the inner parts, the flow becomes dominated by magnetic stresses which eventually lead to the truncation of the disc at an inner radius, $`R_{\mathrm{d},}`$ and the channeling of the flow to the star along stellar magnetic field lines. The twisting of the poloidal field by the vertical shear is assumed to be counteracted by the reconnection of the oppositely directed toroidal fields through the disc vertical extent. This occurs on a rapid timescale resulting to an equilibrium ratio for toroidal, $`B_\phi `$, over vertical, $`B_z`$ field components at the disc surface. Alternatively the growth of $`B_\phi `$ can be controlled by turbulent diffusion through the disc (\[Campbell (1987)\]; \[Yi (1994)\]) or reconnection of field lines in the magnetosphere (i.e. \[Livio & Pringle (1992)\]). In all these cases magnetic torques of a similar form are produced \[Wang (1995)\].
In all calculations of the magnetic torques resulting from disc-star interactions $`B_z`$ is often assumed to be approximately equal to the stellar dipole field. However, twisting magnetic field lines imply the existence of currents which will be distributed in the magnetically dominated magnetosphere which is assumed to corotate with the central star. The star-disc magnetosphere is expected to attain a force–free equilibrium in an Alfvén wave crossing time, as is the case in the solar corona. In the force–free equilibrium, the poloidal magnetic field will deviate from that of a dipole, the effect being especially significant in the case of large twisting (or $`|B_\phi /B_z|`$). We find that the magnetic torque acting on the disc may also be significantly reduced.
The related problem of the evolution of an initially potential field in response to an increasing footpoint shear through a sequence of quasi–static force–free equilibria has been addressed in the solar context. But note that, unlike in the work presented here, ideal MHD is assumed. Numerical calculations supported theoretical expectations (\[Aly (1984)\]; \[Aly (1985)\]; \[Aly (1991)\]; \[Sturrock (1991)\]) that the energy of the force–free field increases monotonically with increasing shear, approaching the energy of an open field (\[Klimchuck (1991)\]; \[Roumeliotis et al. (1994)\]; \[Mikic & Linker (1994)\]). Analytical sequences of quasi–static equilibria were also constructed using a self–similar approach (\[Newman et al. (1992)\]; \[Lynden-Bell & Boily (1994)\]; \[Wolfson (1995)\]) and it was shown that the complete opening of the initially closed field lines is possible at a finite shear. This is associated with the development of current sheets as was conjectured by ?).
?) have calculated magnetospheric force–free equilibria resulting from the twisting of an initially dipolar field anchored on a star which threads an embedded differentially rotating disc. Disc resistivity was taken to be of turbulent origin. However, because of problems associated with their numerical method these authors were not able to calculate equilibrium configurations for $`|B_\phi /B_z|`$ larger than $`2.55`$ at the surface of the disc. In most cases a maximum $`|B_\phi /B_z|1`$ was adopted. Nonetheless significant field line inflation was observed.
In the work presented here, we similarly find that moderate to large shearing of the stellar dipole field through interaction with the disc can lead to a partially open configuration, with expulsion of toroidal field, for favourable outer boundary conditions. Although the process can be inhibited by the presence of an outer boundary that confines the field, significant field line inflation may occur with a large reduction of the magnitude of the vertical field in comparison to the initial dipole value. This could have important consequences for studies of stellar spin evolution which have assumed a dipolar poloidal field at equilibrium (\[Ghosh & Lamb (1979b)\]; \[Cameron & Campbell (1993)\]; \[Yi (1995)\]; \[Ghosh (1995)\]; \[Armitage & Clarke (1996)\]).
After formulating the physical model and basic equations in Section 2, we describe our numerical method in Section 3. Using this we were able to perform calculations with a maximum value of $`|B_\phi /B_z|`$ at the disc surface of $`120`$ and with a variety of external disc radii. Results for Keplerian discs with different amounts of field twisting are given in sections 4. These calculations were performed with a field–confining, perfectly conducting outer boundary.
We also performed calculations with a partly sub–Keplerian accretion disc using the form of angular velocity, $`\mathrm{\Omega }`$ given by ?). This emulates the existence of an inner boundary layer as proposed in the Ghosh & Lamb model. Results are given in Section 4.3 for the conducting outer boundary condition. In Section 4.4 we investigate the effect of the outer boundary condition, showing that a boundary condition that allows field lines to penetrate it, leads to more open configurations.
In Section 4.5 we go on to evaluate the spin down torque acting between star and disc for the force free configurations. We show that this torque can be reduced by up to a factor of $`100,`$ in comparison to what would be obtained assuming an unmodified dipole field, in a configuration that has undergone large twisting and toroidal field expulsion.
Finally in Section 5 we summarize and discuss our results and their application.
## 2 Basic Equations
### 2.1 Force–free equilibria
We consider a magnetic field $`𝑩_{\mathrm{dp}}`$ anchored on a star that rotates with a rate $`\mathrm{\Omega }_{}.`$ We suppose that the magnetic field is axisymmetric with symmetry axis aligned with the rotation axis. A thin accretion disc orbits in the equatorial plane such that the system remains axisymmetric. We ignore the internal dynamics of the disc and assume that the material rotates with a prescribed angular velocity $`\mathrm{\Omega }.`$ For the most part this will equal or be close to the local Keplerian angular velocity $`\mathrm{\Omega }_\mathrm{K}\sqrt{GM/R^3}`$ where $`R`$ is the radial distance to the star measured in the midplane of the disc. We suppose that there is a highly conducting low density corona that corotates with the star and in which the disc is immersed.
However, the disc has non zero resistivity such that the stellar field is enabled to diffuse into the disc. Because the star/corona and disc rotate at different rates, field lines that permeate the disc are twisted in the azimuthal direction such that a toroidal component of the magnetic field is generated.
The disc is considered to be turbulent and manifesting an effective magnetic diffusivity, $`\eta .`$ This enables a steady state to be achieved in which the generation of toroidal magnetic field through differential rotation between the disc and star/corona is balanced by the effects of turbulent diffusion.
The corona responds to the twisting of the initially dipolar field by producing toroidal magnetic field and current density components. These adjust such that a force–free equilibrium is eventually attained. We expect this situation to occur when the coronal Alfvén speed sufficiently exceeds the characteristic rotational speeds. This requires a low density corona which is magnetically dominated. For the purposes of the work presented here, we assume that such a corona exists in which a force–free equilibrium is established in a time short compared to the stellar rotation period. Thus we look for steady state axisymmetric configurations involving the star, the disc and the corona.
An axisymmetric magnetic field $`𝑩`$ can be split into poloidal $`𝑩_\mathrm{p}`$ and toroidal $`B_\phi \widehat{\phi }`$ components such that:
$$𝑩_\mathrm{p}=\mathbf{}\mathbf{\times }(\mathrm{\Psi }\widehat{\phi }/(r\mathrm{sin}\theta )),$$
(1)
where $`\mathrm{\Psi }`$ is the magnetic stream function which is constant on field lines. Here we use spherical polar coordinates $`(r,\theta ,\phi )`$ based on the central star. Then
$$𝑩_\mathrm{p}=(\frac{1}{r^2\mathrm{sin}\theta }\frac{\mathrm{\Psi }}{\theta },\frac{1}{r\mathrm{sin}\theta }\frac{\mathrm{\Psi }}{r})$$
(2)
The force–free condition implies that the Lorentz force is zero everywhere, thus:
$$𝑱\mathbf{\times }𝑩=0$$
(3)
where $`𝑱=(c/4)\times 𝑩`$ is the current density (c.g.s units are used throughout). The toroidal component of equation (3) gives:
$$𝑩_\mathrm{p}\mathbf{}(r\mathrm{sin}\theta B_\phi )=0.$$
(4)
The definition of $`\mathrm{\Psi }`$ implies constancy on field lines or equivalently $`𝑩_\mathrm{p}\mathbf{}\mathrm{\Psi }=0`$. Thus (4) implies that:
$$r\mathrm{sin}\theta B_\phi =f(\mathrm{\Psi })f,$$
(5)
where $`f`$ is an arbitrary function of $`\mathrm{\Psi }.`$
Using (5) and the poloidal component of (3) we finally obtain the governing equation for $`\mathrm{\Psi }`$ in the form:
$$\mathrm{\Delta }\mathrm{\Psi }=ff^{}$$
(6)
where $`f^{}\mathrm{d}f/\mathrm{d}\mathrm{\Psi }`$ and $`\mathrm{\Delta }`$ is the differential operator defined through:
$$\mathrm{\Delta }=\frac{^2}{r^2}+\frac{1\mu ^2}{r^2}\frac{^2}{\mu ^2},$$
(7)
where $`\mu \mathrm{cos}\theta `$.
Equation (6) can be solved for $`\mathrm{\Psi }`$ if $`f`$ is known. Since $`f(\mathrm{\Psi })`$ is constant on poloidal field lines we need to calculate its value at only one point along each field line. This we do by consideration of the interaction of the field with the resistive disc.
### 2.2 Disc–stellar field interaction
The corona above the disc is assumed to corotate with the star and so a toroidal field will be generated due to the vertical shear. Neglecting the radial component of the field interior to the disc, we calculate the internal toroidal field using the toroidal component of the induction equation which in cylindrical coordinates $`(R,\phi ,z)`$ takes the form:
$$\frac{B_\phi }{t}=R𝑩_\mathrm{p}\mathrm{\Omega }+\frac{\eta }{R}\mathrm{\Delta }_\mathrm{c}(RB_\phi )+\frac{1}{R}\eta (RB_\phi )$$
(8)
where $`\mathrm{\Delta }_\mathrm{c}=^2(2/R)(/R)`$ and $`\eta =𝒟\nu `$ is the turbulent magnetic diffusivity which is assumed to be proportional to the turbulent viscosity $`\nu ,`$ the constant of proportionality being $`𝒟.`$ We use the Shakura & Sunyaev (1973) parameterization for $`\nu ,`$ $`\nu =\alpha _{\mathrm{SS}}\mathrm{\Omega }_\mathrm{K}H^2.`$ Here the viscosity parameter is $`\alpha _{\mathrm{SS}}`$ and $`H`$ is the disc semithickness. The gas velocity in the disc is taken to be $`𝒖=(0,r\mathrm{\Omega },0).`$
We consider that $`B_\phi `$ is generated by the shearing of the vertical magnetic field component arising from the $`z`$ dependence of $`\mathrm{\Omega }.`$ For simplicity we take $`\eta `$ to be a function of radial distance only. For a more general treatment see ?). Noting that the disc is thin we expect $`/z(R/H)/R`$ and so we neglect radial derivatives in equation (8). Then in a steady state we have:
$$\eta \frac{{}_{}{}^{2}B_{\phi }^{}}{z^2}=RB_z\frac{\mathrm{\Omega }}{z}$$
(9)
To find $`B_\phi `$ we need to specify $`\mathrm{\Omega }/z`$. We follow previous authors (\[Wang (1987)\]; \[Campbell (1987)\]; \[Yi (1994)\]) and take the vertical shear to be concentrated in a thin layer close to the disc surface. There the angular velocity changes between the disc midplane value and $`\mathrm{\Omega }_{}.`$ We take the vertical average of equation (9) the right hand of which gives:
$$\frac{RB_z}{H}_0^H\frac{\mathrm{\Omega }}{z}𝑑z=\frac{RB_z(\mathrm{\Omega }_{}\mathrm{\Omega })}{H},$$
(10)
where $`B_z`$ is assumed not to vary with $`z`$ thus being the disc midplane value. In the following $`\mathrm{\Omega }`$ will refer to the disc midplane value.
We assume a simple approximation for the vertical average of the dissipative term on the left-hand-side of equation (9) namely:
$$\frac{1}{H}_0^H\eta \frac{{}_{}{}^{2}B_{\phi }^{}}{z^2}𝑑z\eta \gamma \frac{B_\phi }{H^2},$$
(11)
where $`B_\phi `$ now applies to the value at the upper disc surface and $`\gamma `$ is a dimensionless parameter expected to be of order unity.
Combining equation (10) with equation (11) we obtain the following estimate for the equilibrium value of the ratio $`B_\phi /B_z`$ at the upper disc surface:
$$\frac{B_\phi }{B_z}=\pm \frac{RH}{\gamma \eta }(\mathrm{\Omega }_{}\mathrm{\Omega })$$
(12)
We have inserted the $`\pm `$ alternative to allow for different senses of rotation of the star and disc with respect to the right handed coordinate system with the negative sign corresponding to clockwise rotation. The angular velocities are then always positive. Using the turbulent diffusivity specified above and taking $`\gamma =1,`$ equation (12) gives for $`\mathrm{\Omega }=\mathrm{\Omega }_K`$:
$$\frac{B_\phi }{B_z}=\pm \frac{R}{𝒟\alpha _{\mathrm{SS}}H}\left[\left(\frac{R}{R_\mathrm{c}}\right)^{3/2}1\right]$$
(13)
From equation (13) we see that $`B_\phi `$ changes sign at the corotation radius, $`R_\mathrm{c}(GM/\mathrm{\Omega }_{}^2)^{1/3}`$. Similarly the vertically integrated $`z\phi `$ component of the magnetic torque per unit area, $`B_zB_\phi R^2/(2)`$, also changes sign such that the star transfers angular momentum to the disc outside the corotation radius while it gains angular momentum from disc material inside that radius.
In our calculations we shall assume that $`RH\mathrm{\Omega }_{}/(\gamma \eta )C=\text{constant}`$. Equation (13) then becomes (adopting the negative sign corresponding to clockwise rotation):
$$\frac{B_\phi }{B_z}=C\left[\left(\frac{R_\mathrm{c}}{R}\right)^{3/2}1\right].$$
(14)
Our assumption of $`C=\text{constant}`$ is such that for large values of $`R/R_\mathrm{c},`$ $`B_\phi /B_zC`$ which is constant. However, an important result of the calculations, namely the inflation of poloidal field lines such that $`B_z0`$ at the disc surface for large $`C`$, should not depend on having a dependence of $`C`$ on radius provided it remains large. For $`H/R=0.1`$, $`\alpha _{\mathrm{SS}}=0.01`$ and $`𝒟=\gamma =1`$ at the corotation radius, these being reasonable values for accretion discs, equation (14) gives $`|B_\phi /B_z|=C=10^3`$ for $`RR_\mathrm{c}.`$ We comment that equation (13) indicates that in a disc with constant $`H/R`$ and constant $`\alpha _{SS},`$
$$C\frac{1}{𝒟}\left(\frac{R}{R_\mathrm{c}}\right)^{3/2}.$$
(15)
Thus in this case if $`𝒟=1`$, $`C`$ would indeed increase with radius.
Our analysis is distinct from that of e.g. ?) who use a similar expression for $`|B_\phi /B_z|`$ to equation (14) (for $`RR_\mathrm{c}`$) but who assume that fast reconnection in the star–disc corona with the coronal Alfvén speed restricts $`|B_\phi /B_z|`$ to a maximum value of unity. The above discussion suggests that $`C`$ and hence $`|B_\phi /B_z|`$ at the disc surface is large. The actual magnitude of $`|B_\phi |`$ is then controlled by the requirement of a force free equilibrium in the corona. This results in a significant departure of the poloidal field from the original stellar dipole through a significant inflation and opening out of the field lines.
Equation (6) is solved for $`R_\mathrm{d}<r<R_{\mathrm{ext}}`$ where $`R_\mathrm{d}`$ and $`R_{\mathrm{ext}}`$ correspond to the inner and outer disc radii. $`R_\mathrm{d}`$ is usually taken to be small enough that the magnetic torque that is applied there overwelhms the viscous torque so that angular momentum transport is magnetically regulated.
### 2.3 The location of the inner disc radius
An accurate calculation of the inner disc radius would require the solution of the accretion disc structure problem with the inclusion of all magnetic stresses and allowing for sub–Keplerian rotation in the innermost disc region. In most studies an approximate value for $`R_\mathrm{d}`$ has been calculated based on the criterion of disc disruption at a radius where magnetic stresses start to dominate viscous stresses. In order for stellar accretion to proceed the magnetic field must not disrupt the disc exterior to the corotation radius since the magnetic torque for $`R>R_\mathrm{c}`$ imparts angular momentum to the disc gas which cannot connect to the stellar field. Assuming accretion takes place so that $`R_\mathrm{d}<R_\mathrm{c},`$ the inner disc radius can be estimated by requiring that the magnetic torque balances the rate of angular momentum advection by the flow or that:
$$\dot{M}\left[\frac{d\left(\mathrm{\Omega }R^2\right)}{dR}\right]_{R=R_\mathrm{d}}=[B_\phi B_z]_{r=R_\mathrm{d}}R_\mathrm{d}^2$$
(16)
where $`\dot{M}`$ is the mass accretion rate and the positive sign corresponds to clockwise rotation. Because of the rapid increase of the magnetic field strength inwards estimates of $`R_\mathrm{d}`$ using equation (16) (e.g. \[Wang (1995)\]; \[Yi (1995)\]) give $`R_\mathrm{d}`$ close to $`R_c`$ when the star is near its equilibrium spin rate.
However, we point out that this analysis is based on an assumed stellar dipole structure for $`B_z`$. When the effect of an increasing $`B_\phi `$ at the disc surface is taken into account we expect that both $`B_z`$ and the magnetic torques will be modified. We will return to this point below.
In our numerical calculations we have mostly set $`R_\mathrm{d}=R_\mathrm{c}`$ implicitly assuming that the star is close to its equilibrium spin rate. For completeness we also performed some calculations with a sub–Keplerian inner disc region (see Section 4.3).
## 3 Numerical solution
### 3.1 Previous work
Equation (6) has been solved by ?) using a relaxation method with a coordinate transformation $`r1/r`$ so that the solution can extend to arbitrary large radii in principle. ?) adopted equation (13) for $`B_\phi /B_z`$ with a maximum expected value of $`C3.2510^3(R_\mathrm{d}/R)^{1/8}.`$ However, in practice the numerical method would not converge unless $`C`$ was multiplied by a small parameter $`\lambda `$ whose value depends on $`R_\mathrm{c}/R_\mathrm{d}.`$ As a result the maximum value of $`|B_\phi /B_z|`$ at the disc surface that could be obtained in a converged calculation was $`2.55.`$ For this particular case $`\lambda =1.610^5`$ (A.Bardou, private communication). ?) found that the force–free equilibria consist of closed field lines whose footpoints are shifted radially outwards with respect to their position in the dipolar state and with field lines that are flattened at small $`\theta `$ as compared to dipolar ones. ?) found the same indication from a similarity solution. Numerical convergence problems were encountered also by ?) who used a relaxation method to solve equation (6) with a prescribed $`f`$ in the context of the solar corona. They found that no equilibrium solution could be calculated for $`|B_\phi /B_z|\begin{array}{c}>\\ \end{array}1`$. We comment that methods for solving equation (6) based on a simple iteration method for producing a contraction mapping (see also \[Wolfson & Verma (1991)\]) using successive previous approximations for the right hand side tend to fail to converge if $`ff^{}`$ is a sufficiently rapidly varying function of $`\mathrm{\Psi }.`$ By comparison with solutions of related algebraic problems, the lack of convergence in solving equation (6) is not necessarily related to the emergence of additional solutions through a bifurcation. On the other hand studies that employ different methods of solution do not seem to encounter convergence problems (i.e. \[Klimchuck (1991)\]; \[Roumeliotis et al. (1994)\]).
We conclude that previous studies of force–free equilibria related to star–disc interactions were restricted to modest values of $`|B_\phi /B_z|`$ at the disc surface due to restrictions enforced by the adopted numerical method.
### 3.2 Numerical method
In this paper we adopt a different method of solution of equation (6) to those indicated above. In our case equation (6) can be solved without the use of a controlling parameter $`\lambda `$ and in principle for any disc surface value of $`|B_\phi /B_z|.`$ In practice some restrictions also apply in this case but we were able to perform calculations for discs with a ratio of inner to outer radius of several decades and for a maximum value of $`|B_\phi /B_z|10^2`$.
In our method equation (6) is transformed into a parabolic PDE by moving the source term to the left of equation (6) and by equating the new expression to a multiple of the time derivative of $`\mathrm{\Psi }`$ which is expected to tend to zero close to an equilibrium configuration. We use an explicit finite–difference scheme to integrate the parabolic PDE forward in time as an initial value problem until a steady state is achieved. Following this approach rather than solving the elliptic equation (6) directly, we were able to obtain equilibria corresponding to a wide range of values of the source term.
The total magnetic stream function, hereafter denoted by $`\mathrm{\Psi }_\mathrm{t}`$, can be written as a sum of contributions arising from disc currents, $`\mathrm{\Psi }_\mathrm{d},`$ and the central dipole, $`\mathrm{\Psi }_{\mathrm{dp}}.`$ Thus $`\mathrm{\Psi }_\mathrm{t}=\mathrm{\Psi }_{\mathrm{dp}}+\mathrm{\Psi }_\mathrm{d}`$, where $`\mathrm{\Psi }_{\mathrm{dp}}`$ is given by:
$$\mathrm{\Psi }_{\mathrm{dp}}=K\frac{1\mu ^2}{r},$$
(17)
and $`K`$ is a normalizing constant. Note that $`\mathrm{\Delta }\mathrm{\Psi }_{\mathrm{dp}}=0`$ since $`f=0`$ for a dipole field. Equation (6) with $`f0`$ is then solved for $`\mathrm{\Psi }_\mathrm{d}`$ only from:
$$\mathrm{\Delta }\mathrm{\Psi }_\mathrm{d}=f(\mathrm{\Psi }_\mathrm{t})f^{}(\mathrm{\Psi }_\mathrm{t}).$$
(18)
We find it convenient to drop the subscript “$`\mathrm{d}`$” from the disc magnetic stream function and use dimensionless quantities, $`\overline{r}=r/r_{}`$, $`\overline{\mathrm{\Psi }}=\mathrm{\Psi }/\mathrm{\Psi }_0`$ and $`\overline{B}=B/B_0`$ where $`B`$ represents the magnitude of either the toroidal or the poloidal magnetic field components and we take $`\mathrm{\Psi }_0=B_0r_{}^2`$ such that $`B_0=B_{\mathrm{dp}}(r_{},0)`$, the magnitude of the dipole field on the equatorial plane. Equations (6) and (18) can then be read in dimensionless form if all quantities are replaced by their barred equivalents. For simplicity we will drop the bars from the dimensionless quantities from now on.
The computational domain is bounded by an inner radius $`r_{\mathrm{in}}`$, an outer radius $`r_{\mathrm{out}}=R_{\mathrm{ext}}`$, the disc midplane at $`\mu =0`$ and the symmetry axis, $`\mu =1.`$ The grid spacing is taken to be uniform in both $`r`$ and $`\mu .`$
At $`r=r_{\mathrm{in}}`$ and $`\mu =1`$ respectively we specify the boundary condition $`\mathrm{\Psi }=0`$. The first condition specifies only the dipole flux exists at the inner boundary while the second is a requirement of the coordinate system. On the disc midplane, $`\mu =0`$, we impose the symmetry condition, $`\mathrm{\Psi }/\mu =0`$.
We have considered two different outer boundary conditions. The first one is the Dirichlet condition: $`\mathrm{\Psi }=0`$ which forces the field due to the disc to be tangential there. Physically this corresponds to a conducting boundary which excludes the field arising from currents in the disc but not the original dipole field which may be assumed to have had infinite time to diffuse. The second outer boundary condition that we use is the Neumman condition: $`\frac{\mathrm{\Psi }}{r}=0.`$ This makes the disc field radial at the outer boundary as might be the case if there were a coronal wind there. The second boundary condition leads to much more open configurations than the first and hence to larger departures from the stellar dipole field.
We add a term $`F(\mathrm{\Psi }/\tau )`$ to the right-hand side of the dimensionless form of equation (18) where we are allowed to choose $`F`$ as an arbitrary function of position and solve the equation with a forward in time, $`\tau `$, and centered in space (FTCS) finite–difference scheme on a $`(Nr\times Nm)`$, $`(r,\mu )`$ computational grid. The finite–difference equation that we use is given at a grid point denoted with subscript $`(i,j)`$ by:
$`F_{i,j}{\displaystyle \frac{\mathrm{\Psi }_{i,j}^{n+1}\mathrm{\Psi }_{i,j}}{\mathrm{d}\tau }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{i+1,j}2\mathrm{\Psi }_{i,j}+\mathrm{\Psi }_{i1,j}}{(\mathrm{d}x)^2}}`$ (19)
$`+`$ $`{\displaystyle \frac{1\mu _j^2}{r_i^2}}{\displaystyle \frac{\mathrm{\Psi }_{i,j+1}2\mathrm{\Psi }_{i,j}+\mathrm{\Psi }_{i,j1}}{(\mathrm{d}\mu )^2}}+f\left(\mathrm{\Psi }_\mathrm{t}|_{i,j}\right)f^{}\left(\mathrm{\Psi }_\mathrm{t}|_{i,j}\right).`$
where $`\mathrm{d}r,`$ and $`\mathrm{d}\mu ,`$ are the uniform grid spacings in $`r,`$ and $`\mu ,`$ respectively. Variables without superscripts are at $`\tau `$ level $`n.`$ The time step is $`\mathrm{d}\tau .`$ We assume that the disc is truncated at $`R_\mathrm{d}`$ by an infinite conductor so we have $`B_\phi =0`$ for $`1<RR_\mathrm{d}.`$ For all other values of $`R`$ on the equator ($`\mu =0`$ denoted by subscript $`1`$) we use equation (12) approximated as follows:
$$B_\phi |_{i+1/2,1}=CB_z|_{i+1/2,1}\frac{\mathrm{\Omega }_{i+1/2}\mathrm{\Omega }_{}}{\mathrm{\Omega }_{}}.$$
(20)
In equation (20) $`\mathrm{\Omega }_{i+1/2}`$ is either given by the Keplerian rate or by the form proposed by ?) (see Section 4.3). $`\mathrm{\Omega }_{}`$ is calculated on the $`(i+1/2,1)`$ grid point that is closest to the prescribed value of $`R_\mathrm{c}`$. The vertical magnetic field at the disc surface is calculated at the $`(i+1/2,1)`$ points by numerical differentiation of equation (2).
Using equation (20) we calculate $`f_{i+1/2,1}=R_{i+1/2}B_\phi |_{i+1/2,1}`$ while $`(ff^{})_{i,1}`$ is calculated from:
$$(ff^{})_{i,1}=\frac{1}{2}\frac{f_{i+1/2,1}^2f_{i1/2,1}^2}{\mathrm{\Psi }_\mathrm{t}|_{i+1/2,1}\mathrm{\Psi }_\mathrm{t}|_{i1/2,1}}.$$
(21)
Since $`f`$ is a function of $`\mathrm{\Psi }_\mathrm{t}`$ only we construct a table of values of $`f`$ as function of $`\mathrm{\Psi }_\mathrm{t}|_{i,1}`$ on the equator. We then use this table to calculate $`f`$ for any required value of $`\mathrm{\Psi }_\mathrm{t}`$ in equation (19) using linear interpolation.
The main disadvantage of the method we use is the constraint on the time step coming from the requirements of numerical stability. We find that the maximum allowed timestep decreases rather steeply with $`C.`$
In order to achieve a more rapid evolution for the same computational time and always keeping in mind that we are only interested in final steady state equilibria, we have introduced a spatially varying “diffusion coefficient”, $`(1/F)`$, that multiplies the right-hand side of equation (19). This coefficient is taken to be equal to $`\mathrm{d}\tau _{i,j}/\mathrm{min}(\mathrm{d}\tau _{i,j})`$ and therefore allows for a more advanced evolution of the values at the grid points where $`\mathrm{d}\tau _{i,j},`$ the maximum allowed timestep based on local stability considerations, can be larger. Since the source term peaks near the corotation radius but it is significantly smaller elsewhere it is possible to reach the equilibrium in significantly smaller computational time in this way.
Time integrations have been performed until the right hand side of equation (19) becomes smaller than $`10^6`$ in dimensionless units although the magnetic field configuration is usually found to have converged to its final form well before this criterion is satisfied.
## 4 Numerical Results
The parameters used in the calculations presented here are summarised in Table 2. All the calculations were performed with $`r_{\mathrm{in}}=1`$. Different values were used for $`R_{\mathrm{ext}}`$ and $`C`$ and calculations were done using both outer boundary conditions on $`\mathrm{\Psi }.`$ All the calculations with $`\mathrm{\Omega }=\mathrm{\Omega }_\mathrm{K}`$ were performed with $`R_\mathrm{d}=R_\mathrm{c}`$. Our choice of $`R_\mathrm{c}=3`$ implies $`\mathrm{\Omega }_{}=0.19\mathrm{\Omega }_\mathrm{K}(R=R_{})`$ which corresponds to a rotation period of $`5\mathrm{days}`$ for $`R_{}=3R_{}`$ and $`M_{}=0.5M_{}`$. The observed rotation periods of CTTS are $`3`$$`12`$ days (\[Bouvier et al. (1997)\] and references therein).
The choice of $`R_{\mathrm{ext}}`$ was dictated mainly by the limitations of the computational method since a much larger value of $`R_{\mathrm{ext}}`$ would require increased values of $`Nr`$ and $`Nm`$ for the same resolution to be achieved as in the cases with smaller outer disc radii.
We chose $`C=3`$$`120`$ in order to examine the effect of an extended range of values of $`|B_\phi /B_z|`$ at the disc surface on the magnetospheric equilibrium. Various authors have assumed $`|B_\phi /B_z|1`$ at the disc surface in their analysis. ?) argue that anomalous resistivity in the disc would limit $`B_\phi /B_z.`$ ?) and ?) assert that fast reconnection outside the disc would lead to a similar result. These physical arguments have not so far been supplemented by detailed calculations in the context of star–disc interactions. Our maximum value of $`C`$ was chosen so that the computational times required are reasonable. Nevertheless much larger values of $`C`$ and therefore of $`B_\phi `$ on the disc surface could result in very large toroidal magnetic pressure on the disc surface and the subsequent destabilisation of the disc.
In Table 2 we have also included the parameters of calculations performed with a sub–Keplerian form of $`\mathrm{\Omega }`$. This is characterised by the “fastness” parameter $`\omega \mathrm{\Omega }_{}/\mathrm{\Omega }_\mathrm{K}(R_\mathrm{d})`$ (for more details see Section 4.3).
### 4.1 Qualitative considerations of star–disc interactions
The main result of the production of toroidal field in the star–disc corona is the presence of a feedback mechanism that reduces $`|B_z|`$ at the disc surface. In this way the magnitude of $`B_\phi `$ does not become too large to maintain a force free equilibrium when $`C`$ is very large. As $`C`$ is increased, the field becomes increasingly twisted. The result is the production of azimuthal currents as well as a toroidal field component in the corona. The azimuthal currents tend to reduce $`|B_z|`$ near the disc and hence the magnitude of $`B_\phi .`$ To see this in a qualitative way: the disc flux obeys equation (18) which, for the boundary conditions used, can be converted to an integral form:
$$\mathrm{\Psi }=_VG(r,\mu ,r^{},\mu ^{})f(\mathrm{\Psi }_\mathrm{t})f^{}(\mathrm{\Psi }_\mathrm{t})r^2𝑑\mu ^{}𝑑r^{}.$$
(22)
Here $`G`$ is an appropriate Green’s function and the stream function inside the integral, taken over the computational domain, is evaluated using the primed coordinates. Thus:
$$\mathrm{\Psi }_\mathrm{t}=_VG(r,\mu ,r^{},\mu ^{})f(\mathrm{\Psi }_\mathrm{t})f^{}(\mathrm{\Psi }_\mathrm{t})r^2𝑑\mu ^{}𝑑r^{}+\mathrm{\Psi }_{\mathrm{dp}}.$$
(23)
For small $`C,`$ $`f`$ is small and $`\mathrm{\Psi }_\mathrm{t}`$ is close to the dipole value $`\mathrm{\Psi }_{\mathrm{dp}}.`$ However, for large $`C`$ and hence $`f,`$ a solution can only be found (or the force free equilibrium can only be maintained) by otherwise reducing the magnitude of the integral. Note that if $`\mathrm{\Psi }_{\mathrm{dp}}`$ is scaled by multiplying by a constant value, both $`\mathrm{\Psi }_\mathrm{t}`$ and the integral above scale similarly. The magnitude of the integral can be reduced by having fewer field lines intersecting the disc and therefore more field lines for which $`f=0`$. If the outer boundary condition does not allow that, the magnitude of the integral can also be reduced by the field lines intersecting the disc at progressively larger radii (and smaller field magnitudes), as $`C`$ increases. In either case, the magnitude of $`B_\phi `$ is controlled and the poloidal field takes on either a more inflated or a partially open structure.
We now discuss our results for a disc with Keplerian rotation everywhere and $`R_\mathrm{d}=R_\mathrm{c}=3`$ (those cases are annotated with $`\omega =1`$ in Table 2).
### 4.2 Keplerian disc rotation profile
The calculations presented here were performed using five values of $`C`$ in the range $`3`$$`120`$ and using the Dirichlet outer boundary condition. The maximum value of $`R_{\mathrm{ext}}`$ used in these calculations decreases with $`C`$ (see Table 2). As $`C`$ increases the resolution required to resolve fully the magnetic field structure around the corotation radius increases. We have therefore restricted the external radius for $`C=30`$ or larger so that a reasonable resolution of the area around corotation can be achieved.
The maximum value of $`|B_\phi /B_z|`$ at the disc surface found in the calculations with a Keplerian disc rotation profile is $`100`$ being much larger than values considered by previous authors. Consequently we study the effect of increasing $`|B_\phi /B_z|`$ on the force–free magnetospheric field for a large region of parameter space.
Below we present results obtained with $`C=3`$ corresponding to maximum value of $`|B_\phi /B_z|=2.56`$$`2.97`$ for $`R_{\mathrm{ext}}=11`$$`80`$. When $`C=3`$ we were able to achieve the highest resolution and the largest range of $`R_{\mathrm{ext}}.`$
#### 4.2.1 Results with $`C=3`$.
All the calculations with $`C=3`$ were performed with $`Nr=200`$ and $`Nm=80`$ therefore cases with smaller $`R_{\mathrm{ext}}`$ have better resolution.
In the left panel of Fig. 2 we plot the equilibrium radial profiles of $`|B_\phi |`$ at the disc surface. Since these profiles are well resolved we can conclude that increasing the computational domain leads to smaller values of $`|B_\phi |`$. This follows from the dependence of $`B_z`$ on $`R_{\mathrm{ext}}.`$ In the right hand panel of Fig. 2 we plot the radial profile of $`B_z`$ at the disc surface for the different values of $`R_{\mathrm{ext}}`$.
The poloidal field is similar to the stellar dipole field inside corotation where $`B_\phi =0`$ but it deviates from that significantly at all other radii. Immediately outside corotation the field becomes smaller than the stellar dipole field as a result of field line inflation. At the largest radii the poloidal field eventually becomes larger than the dipole field locally. The radius at which this transition occurs moves to larger radii as $`R_{\mathrm{ext}}`$ increases indicating the influence of the outer boundary condition. The structure of the magnetic field is not radially self–similar. The departure from self–similarity in the inner part of the disc is due to the transition from $`B_\phi =0`$ inside the corotation radius to $`B_\phi `$ given by equation (14) for $`R>R_\mathrm{c}`$ while the outer part of the disc is influenced by the outer boundary condition.
Field line expansion as seen in Fig. (4) for all $`R_{\mathrm{ext}}`$ used is a consequence of field line twisting together with the assumption of a force–free equilibrium (see our discussion in Section 4.1). The inflation of field lines increases with increasing $`C`$ or field line footpoint shear until the field lines become partially or fully open (\[Aly (1985)\]; \[Wolfson & Low (1992)\]; \[Newman et al. (1992)\]; \[Lynden-Bell & Boily (1994)\]; \[Aly (1995)\]; \[Wolfson (1995)\]). In the fully open state line twisting disappears and $`B_\phi =0.`$
In Fig. 6 we plot the ratio of the final $`R_\mathrm{f}`$ to the initial or dipole $`R_\mathrm{i}`$ footpoint radius as a function of $`R_\mathrm{i}`$ for a group of field lines. The values of $`R_\mathrm{i}`$ chosen are the same as in Fig. 4 with $`R_\mathrm{i}4`$. Different curves correspond to different values of $`R_{\mathrm{ext}}`$. The ratio $`R_\mathrm{f}/R_\mathrm{i}`$ increases linearly with $`R_\mathrm{i}`$ until a turnover radius is reached where the outer boundary starts affecting the rate of inflation. For $`R_\mathrm{i}`$ larger than the turnover radius $`R_\mathrm{f}/R_\mathrm{i}`$ decreases with $`R_\mathrm{i}`$. For the cases with $`R_{\mathrm{ext}}=60`$ and $`80`$ the values of $`R_\mathrm{f}/R_\mathrm{i}`$ are similar for $`R_\mathrm{i}10`$. We therefore conclude that the maximum inflation has been reached for the innermost disc region in these cases. The maximum overall value of $`R_\mathrm{f}/R_\mathrm{i}`$ is $`3.2`$ and it corresponds to $`|B_\phi /B_z|`$ of $`2.73`$. The maximum value of $`R_\mathrm{f}/R_\mathrm{i}`$ quoted by ?) is 1.9 but this is an underestimate because of resolution problems (A.Bardou, private communication). Due to the effect of the outer boundary in our calculations it is difficult to predict how inflated the configuration would become for $`R_{ext}\mathrm{}.`$ However, we expect the field lines will remain closed within a radius of approximately $`3R_\mathrm{c}`$ for $`C=3`$.
#### 4.2.2 Results for larger values of $`C`$
We performed calculations for $`C=10`$ with $`R_{\mathrm{ext}}=60`$ and $`80.`$ Both these have $`Nr=100`$ and $`Nm=50`$ so the resolution is reduced with respect to the $`C=3`$ case. In Fig. 8 we plot the radial profiles of $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ and $`B_z(r,0)`$ for $`C=10`$ and also for $`C=3`$ and $`R_{\mathrm{ext}}=80`$. For $`R_{\mathrm{ext}}=80`$ we find that $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ is smaller for $`C=10`$ than for $`C=3`$ everywhere in the disc as expected. For the smaller $`R_{\mathrm{ext}}`$ value the effect of the outer boundary is stronger as it was found also for $`C=3`$. We find that $`B_z`$ decreases with $`C`$ for $`R<15`$. For larger radii the trend is reversed because of the effect of the outer boundary.
In the left panel of Fig. 10 we plot dipole and equilibrium poloidal field lines for $`C=10`$ with $`R_{\mathrm{ext}}=80`$. Field lines interior to $`2R_{ext}/3`$ are characterised by the largest flattening at small $`\theta `$ while the outer field lines are more spherical due to the effects of the outer boundary. A value of $`R_{\mathrm{ext}}`$ exceeding the outer disc radius by a factor $`10`$$`100`$ might be necessary in order to remove the effects of the outer boundary on the calculation of field line inflation (see \[Roumeliotis et al. (1994)\]).
In the right panel of Fig. 10 we plot $`R_\mathrm{f}/R_\mathrm{i}`$ as function of $`R_\mathrm{i}`$ for $`C=3`$ and $`C=10`$, each with $`R_{\mathrm{ext}}=60`$ and $`80`$. For $`R_\mathrm{i}8`$ the ratio $`R_\mathrm{f}/R_\mathrm{i}`$ is $`1.75`$ times larger for $`C=10`$ than for $`C=3`$ showing greater inflation. But for larger values of $`R_\mathrm{i}`$ this factor drops to $`1.25`$. The larger influence of the outer boundary on the equilibrium for $`C=10`$ is manifested in the smaller turnover radius observed in this case. It is not clear if the ratio of final to initial footpoint radius found for $`C=10`$ at $`R_\mathrm{i}=6`$ would continue to increase linearly with $`R_\mathrm{i}`$ if the outer boundary effect was not present. In the ?) results $`R_\mathrm{f}/R_\mathrm{i}`$ increases linearly with $`R_\mathrm{i}`$ for $`RR_\mathrm{c}.`$ In our case $`|B_\phi /B_z|`$ tends to a constant for $`RR_\mathrm{c}`$ and $`R_\mathrm{f}/R_\mathrm{i}`$ could behave similarly.
We now discuss results obtained for $`C=30`$$`120`$. These calculations were performed with $`Nr=50`$ and $`Nr=30`$ and $`R_{\mathrm{ext}}=11`$ except for one case with $`C=30`$ and $`R_{\mathrm{ext}}=21`$ which had $`Nr=120`$ and $`Nm=50`$. The disc surface equilibrium profiles of $`B_\phi `$ for $`C=3`$$`120`$ and $`R_{\mathrm{ext}}=11`$ are compared in Fig. 12. Note that the magnitude of $`B_\phi `$ varies only by a factor of $`3`$ as $`C`$ varies between $`3`$ and $`120.`$
The effect of the growing $`|B_\phi |`$ on the stream function for $`C30`$ and for $`R_{\mathrm{ext}}=11`$ is to make it larger than the dipole value, more so around $`R=R_\mathrm{c},`$ leading to an enhancement of $`B_z(r,0)`$ there. The disc surface radial profiles of $`B_z`$ for $`C=3`$$`120`$ are compared in the left panel of Fig. 14. Because the toroidal magnetic pressure increases sharply just outside $`R=R_\mathrm{c}`$ it is expected that a larger poloidal field just interior to $`R_c`$ is required for equilibrium. Thus field line deflation results. This occurs in this model because the field lines are not sheared inside corotation. For $`R>R_\mathrm{c}`$ $`B_z(r,0)`$ becomes smaller than the dipole value before increasing close to the outer boundary. Due to the increased $`|B_\phi /B_z|`$ at the disc surface, $`B_z`$ decreases with $`C`$ outside corotation as expected.
To study the dependence on the outer boundary radius for large $`C`$, we performed a calculation with $`C=30`$ and $`R_{\mathrm{ext}}=21`$ and compared the result to the case where $`R_{\mathrm{ext}}=11`$. The radial resolution is approximately the same for these cases (see Table 2). The corresponding radial profiles of $`B_z(r,0)`$ are plotted in the right panel of figure 14.
When $`R_{\mathrm{ext}}`$ is increased $`B_z(r,0)`$ values are unmodified inside $`R=10`$ so the deflation of the field lines there is independent of the outer disc radius. The exterior point at which $`B_z(r,0)`$ exceeds the dipole value moves to a larger radius for $`C=30`$ as compared to that when $`C=3`$. There is also a change to the radial dependence of $`B_z(r,0)`$ which becomes flatter at large $`R`$ for $`C30`$.
In Fig. 16 we plot the magnetic field lines for $`C=30`$ with $`R_{\mathrm{ext}}=11`$ (left panel) and $`R_{\mathrm{ext}}=21`$ (right panel). We see that the deflation of field lines for $`r<10`$ is the same in these cases. At larger radii inflation rather than deflation is observed for $`R_{\mathrm{ext}}=21`$. The field line inflation observed at $`r>10`$ for $`C=30`$ with $`R_{\mathrm{ext}}=21`$ is smaller than that found for $`C=3`$ and the same $`R_{\mathrm{ext}}`$ (compare with Fig. 4). This result is unexpected given the larger $`|B_\phi /B_z|`$ at the disc surface in the $`C=30`$ case. This is probably the result of the increased influence of the outer boundary condition on the equilibrium configuration for increasing $`C`$ values.
In the next section we will discuss results obtained using a sub–Keplerian form of $`\mathrm{\Omega }`$ and which results to a smoother radial variation of $`B_\phi /B_z`$ on the disc surface.
### 4.3 Sub–Keplerian disc rotation profile
As the magnetic field strength increases with decreasing disc radius and the magnetic stresses start to affect the radial and vertical disc equilibrium the disc is expected to ultimately corotate with the star. Interior to the corotation radius the disc will then be sub–Keplerian. In the present study we investigate the effect of sub–Keplerian rotation on $`B_\phi /B_z`$ and the equilibrium poloidal magnetic field. We modify $`\mathrm{\Omega }`$ so that it becomes sub–Keplerian for $`RR_\mathrm{c}.`$ The form of $`\mathrm{\Omega }`$ that we will use is that of ?) which is given, for $`R>R_\mathrm{d}`$, by:
$$\mathrm{\Omega }=\mathrm{\Omega }_\mathrm{K}(R_\mathrm{d})\left\{\left(\frac{R_\mathrm{d}}{R}\right)^{3/2}(1\omega )\mathrm{exp}\left[\frac{3}{2}(1\omega )^1\left(\frac{R}{R_\mathrm{d}}1\right)\right]\right\}.$$
(24)
Here $`\omega =\mathrm{\Omega }_{}/\mathrm{\Omega }_\mathrm{K}(R_\mathrm{d})`$. At $`R=R_\mathrm{d}`$ we have $`\mathrm{\Omega }=\mathrm{\Omega }_{}`$ and $`\mathrm{d}\mathrm{\Omega }/\mathrm{d}R=0`$. Interior to $`R=R_\mathrm{d}`$, the material is taken to rotate with $`\mathrm{\Omega }=\mathrm{\Omega }_{}.`$ We plot $`\mathrm{\Omega }`$ as a function of $`R`$ in the left panel of Fig. 18 for $`R_\mathrm{c}=3`$ and $`\omega =0.875`$ and $`0.2`$. The condition: $`\mathrm{\Omega }=\mathrm{\Omega }_{}`$ at $`R=R_\mathrm{c}`$ determines $`R_\mathrm{d}`$ through the chosen value of $`\omega `$.
For the largest value of $`\omega `$, $`R_\mathrm{d}`$ almost coincides with $`R_\mathrm{c}`$ and $`\mathrm{\Omega }`$ remains Keplerian until very close to $`R_\mathrm{c}`$. For smaller values of $`\omega `$ we have $`R_\mathrm{d}<R_\mathrm{c}`$ and in the case where $`\omega =0.2,`$ $`R_\mathrm{d}`$ almost reaches the stellar surface ($`R=1`$). Here as in the previous sections we have $`B_\phi =0`$ for $`R<R_\mathrm{d}`$ but a part of the differentially rotating disc may exist inside corotation and that can affect the equilibrium magnetic field. The variation of $`B_\phi /B_z`$ with $`r`$ at the disc surface can be seen in the right panel of Fig. 18 where we took $`C=30.`$
For $`\omega =0.875`$ the profile of $`B_\phi /B_z`$ is virtually unchanged with respect to the Keplerian case with $`R_\mathrm{d}=R_\mathrm{c}.`$ A larger effect on $`B_\phi /B_z`$ occurs when $`\omega =0.2.`$ In that case there is a region inside corotation where $`B_\phi /B_z`$ reverses sign.
#### 4.3.1 Results with $`\omega =0.875`$ and $`C=30`$$`120`$
These calculations have $`R_{\mathrm{ext}}=60.`$
In the left panel of Fig. 20 we plot $`B_z(r,0)`$ for $`C=30`$$`120`$ with $`\omega =0.875`$ and $`R_{\mathrm{ext}}=60`$ with the corresponding $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ profiles shown in the right panel. We also plot the case $`C=3`$ with $`\mathrm{\Omega }=\mathrm{\Omega }_\mathrm{K}`$. Apart from boundary effects, $`|B_z|`$ decreases with $`C`$, the decrease being faster between $`C=60`$ and $`120`$ than between $`C=30`$ and $`60`$. These results reinforce our expectation: $`B_z0`$ for $`C\mathrm{}`$.
As seen in the right panel of Fig. 20 the disc values of $`\mathrm{\Psi }_\mathrm{t}`$ increase with $`C`$ between $`C=30`$ and $`120`$. This is contrary to theoretical prediction and to the trend observed between $`C=3`$ and $`30`$. As we will see in Section 4.4 this result depends strongly on the outer boundary condition.
The profiles obtained in the non–Keplerian case with $`\omega =0.875`$ are smoother around corotation than in the Keplerian case. This is due to the smoother radial variation of $`B_\phi /B_z`$. However, $`B_z(r,0)`$ is still larger than the dipole value there and some deflation is observed interior to $`r=R_\mathrm{c}`$. At larger radii field lines are inflated as we would expect given the increase in $`R_{\mathrm{ext}}`$.
#### 4.3.2 $`\omega =0.2`$
In this case $`R_\mathrm{d}/R_\mathrm{c}=0.34`$ and an inversion of the sign of $`B_\phi `$ may significantly affect the variation of $`B_z`$ with radius. We adopted $`C=3`$ with $`R_{\mathrm{ext}}=21`$ and $`C=30`$ with $`R_{\mathrm{ext}}=11`$ with two different resolutions (see Table 2 for details). For $`C=3`$ the difference between Keplerian and non–Keplerian cases is not significant.
In the left panel of Fig. 22 we plot the disc surface profiles of $`B_\phi `$ for $`C=30`$ and $`\omega =0.2`$ obtained with two different resolutions and for $`C=30`$ with $`\mathrm{\Omega }=\mathrm{\Omega }_\mathrm{K}`$. For $`C=30`$ the larger radial variation of $`|B_\phi |`$ in the non–Keplerian case has a more significant effect on the equilibrium. As the radial gradient of $`B_\phi ^2`$ increases for small radii a small deflation of the field lines is observed there. Eventually $`B_\phi ^2`$ decreases with radius and inflation is observed for $`R>2`$. The reverse was observed in the Keplerian case where field lines were found to be deflated everywhere in the disc (always for $`R_{\mathrm{ext}}=11`$).
In the right panel of Fig. 22 we plot the radial profile of $`B_z(r,0)`$ for the same parameters as in the left panel. We note that when $`\omega =0.2`$ the field deviates more from the dipole form around corotation, becoming larger than the dipole value locally, than in the case with $`\mathrm{\Omega }=\mathrm{\Omega }_\mathrm{K}.`$ This difference is mainly due to the larger radial variation of $`|B_\phi |`$ with radius in the non–Keplerian case. The difference observed between the profiles obtained in the high and low resolution cases is mostly due to the radial shift of the gridpoint corresponding to $`R=R_\mathrm{c}`$.
We conclude that the deflation observed in previous calculations with a Keplerian rotation profile and with $`\omega =0.875`$ is related to the restricted twisting of field lines inside the corotation radius. In the case with $`\omega =0.2`$ we have $`B_\phi 0`$ almost everywhere in the disc. Consequently field lines are much more inflated in this case.
### 4.4 A different outer boundary condition
To study the effect of the outer boundary condition on our results we have also performed calculations with the Neumann condition: $`\frac{\mathrm{\Psi }}{r}=0`$ and for $`C=3,10,60`$ and $`120`$. The parameters used in these calculations are given in Table 2. All cases have $`\omega =0.875`$ except when $`C=3`$ where a Keplerian rotation profile was used. The new boundary condition makes the disc field normal to the outer boundary. Thus we expect some of the field lines to be open at the equilibrium, at least at large disc radii far from the disc midplane.
Indeed the results show much larger inflation and opening of the field lines in this case than in the cases presented in previous sections where the Dirichlet condition was used. Also almost no deflation is observed inside corotation showing the global effect of the outer boundary condition. As seen in Fig. 24 both $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ and $`B_z(r,0)`$ decrease with $`C`$ monotonically. For $`C=120`$ the equilibrium seems to have reached a limiting configuration with almost all field lines open.
Comparing the disc radial profiles of $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ depicted here with those in Fig. 20 we note that $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ decreases considerably when the second boundary condition is used. Also $`\mathrm{\Psi }_\mathrm{t}(r,0)`$ becomes approximately independent of $`r`$ everywhere outside the corotation area which is the limiting behaviour that we expect for large values of $`C`$. Comparing the disc radial profiles of $`B_z(r,0)`$ between figures 20 and 24 we note that for $`C=3`$ the magnetic field is similar for the two boundary conditions for $`r10`$. Further out $`B_z(r,0)`$ is smaller for the second boundary condition. This difference increases with $`C`$. We also note that $`B_z(r,0)`$ has a stronger dependence on $`r`$ when the second boundary condition is used.
In the first three panels of Fig. 26 we plot the magnetic field lines for $`C=3,10`$ and $`60`$ respectively, using the same initial (dipole) field lines as in Fig. 4 (with $`R_{\mathrm{ext}}=60`$). For $`C=3`$ the level of inflation is comparable for the two boundary conditions for $`r\begin{array}{c}<\\ \end{array}10`$. For larger initial footpoint radii field lines become much more inflated for the second outer boundary condition and eventually they open up for $`r15`$. As $`C`$ increases more field lines become open as seen by comparing the cases with $`C=10`$ and $`60`$. The case with $`C=120`$ is virtually identical to the $`C=60`$ case. Some field lines remain closed even when $`C=60`$ as can be seen by plotting more contour levels (see lower right panel of Fig. 26). The closed field lines correspond to initial radii which are very close to the inner disc radius. Consequently $`ff^{}`$ is non-zero for an increasingly smaller range of $`\mathrm{\Psi }_\mathrm{t}`$ values as $`C`$ increases. Calculations for larger values of $`C`$ could only be performed accurately using considerably larger numerical resolution. However, our present calculations have already shown that field lines at equilibrium will be approximately open for $`|B_\phi /B_z|10^2`$ if a Neumman outer boundary condition is used.
The ratio $`|B_\phi |/B_\mathrm{p}`$ is found to be less than unity for $`\mu <0.1`$. Therefore $`B_\phi 0`$ in most of the corona as expected on open field lines. In the limit of a large twist the disc field acts as to expel the poloidal magnetic field from the disc and consequently to annihilate the toroidal magnetic field, as long as the disc field is permitted to penetrate the outer boundary. In the limit of large $`C`$ the poloidal field resembles the solution $`Z_1`$ given by ?). This corresponds to a perfectly conducting disc with central hole immersed in a central dipole field where some flux escapes to infinity such that the field is non singular at the disc inner edge.
### 4.5 Modification of the spin–down and spin–up torques
For the force–free equilibria calculated above, the poloidal magnetic field differs significantly from the original stellar dipole form. As a consequence, the local magnetic torque acting on the disc may also differ significantly from what one obtains assuming $`B_zB_{\mathrm{dp}}`$. Since this torque regulates the spin evolution of the central star it is important to examine the possible effect of our results on the calculation of the total torque experienced by the star.
The total magnetic torque $`N`$ can be split into two components of opposite sign. The first arises inside corotation where angular momentum is transfered from the disc to spin up the star. The second arises outside corotation where angular momentum is transferred to the disc. In order for stellar accretion to occur we require $`R_\mathrm{d}<R_\mathrm{c}`$. In one set of calculations we assumed $`R_\mathrm{d}=R_\mathrm{c}`$ and adopted Keplerian rotation. In that case the spin–up torque would be zero. In reality we expect $`R_\mathrm{d}/R_\mathrm{c}0.91`$$`0.97`$ (i.e. \[Wang (1995)\]) giving rise to a relatively small spin–up torque. This torque is probably small compared to that arising from direct accretion onto the star.
We now calculate the total torques for the case where the disc is assumed to be Keplerian and truncated at the corotation radius.
#### 4.5.1 The Keplerian case
The total magnetic torque on the star is in general given by (e.g. \[Ghosh & Lamb (1979b)\]):
$$N=_{R_\mathrm{d}}^{R_{\mathrm{ext}}}B_\phi B_zR^2dR.$$
(25)
We adopt dimensionless parameters such that the torque is given in units of $`N_{}B_{}^2r_{}^3`$. Of course all field values are calculated on the disc surface ($`\mu =0`$). For the cases presented here the first boundary condition has been used. When $`R_\mathrm{d}=R_\mathrm{c},`$ as is considered here, $`N`$ is positive definite and only acts to spin down the star.
The total torque $`N`$ can be calculated analytically for a dipole field with $`B_\phi `$ given by equation (14) and is given by:
$$N^{\mathrm{dp}}=\frac{C}{3}\left[\frac{1}{3}+\frac{2}{3}\left(\frac{R_{\mathrm{ext}}}{R_\mathrm{c}}\right)^{9/2}\left(\frac{R_{\mathrm{ext}}}{R_\mathrm{c}}\right)^3\right]R_\mathrm{c}^3.$$
(26)
In the left hand panel of Fig. 28 we plot $`N/C`$ and $`N^{\mathrm{dp}}/C`$ as a function of $`C`$ for $`R_{\mathrm{ext}}=11`$. Although $`N^{\mathrm{dp}}/C`$ is constant with $`C`$, $`N/C`$ decreases with increasing $`C`$ as the poloidal magnetic field in the force–free equilibria decreases with increasing $`C`$. The ratio of force–free to dipole torque $`N/N^{\mathrm{dp}}`$ for $`C30`$ varies between $`0.14`$$`0.013`$ which is significantly smaller than the same ratio in the case with $`C=3`$. Although the effect of the outer boundary condition is to suppress inflation in all cases with large $`C`$ values and $`R_{\mathrm{ext}}=11`$, the effect of the twisting of the field lines on the final equilibrium is an increased reduction of the total torque with respect to its corresponding dipolar value.
In the right panel of Fig. 28 we plot $`N`$ and $`N^{\mathrm{dp}}`$ where we see that the behaviour of $`N`$ as function of $`C`$ is very different (in fact almost reversed) to that expected for a dipolar poloidal field. From this figure we see that $`N`$ increases slowly for $`C15`$ and it decreases for larger values of $`C`$. As the values of $`N`$ for $`C30`$ are between $`10`$ and $`100`$ times smaller than the dipolar values the increasing field line twisting has an important effect on the spin down of the star. The magnetic breaking timescales calculated are longer and therefore magnetic breaking models might have to be reevaluated. However, we have not taken into account here the disc region inside corotation. Also the results described here for $`C30`$ have $`R_{\mathrm{ext}}=11`$ which is relatively small and therefore the effect of the outer boundary condition is important.
#### 4.5.2 The sub–Keplerian case
In the case where $`\mathrm{\Omega }\mathrm{\Omega }_\mathrm{K}`$ the total magnetic torque includes both spin–up and spin–down components. For the case with $`\omega =0.875`$ the spin–up component is of no significance. In the following we will compare the total magnetic torques calculated using a sub–Keplerian rotation profile with $`\omega =0.875`$ to the corresponding dipolar values. Calculations with similar parameters were performed with both outer boundary conditions.
In Fig. 30 we plot $`N/C`$ and $`N^{\mathrm{dp}}/C`$ as a function of $`C`$. The values in the left panel of this figure are from calculations performed with the first outer boundary condition. Here, as for the Keplerian case, we find that $`N/C`$ decreases with $`C`$ and that the ratio $`N/N^{\mathrm{dp}}`$ ranges from $`0.56`$ for $`C=3`$ to $`0.05`$ for $`C=120`$. The larger ratios here are mainly a consequence of the use of a sub–Keplerian rotation profile. However, For $`C30`$ we have again $`N/C1/C`$ approximately. Therefore, the trend for a smaller total torque is clear and the difference between force–free and dipole cases is again significant for $`C1`$. For $`B_\phi B_z`$ the difference is not as large although we have to reserve final judgement until calculations with much larger values of $`R_{\mathrm{ext}}`$ are performed.
For calculations performed with the second outer boundary condition we found significantly smaller values of $`B_z(r,0)`$ which decreases faster with $`r`$ than when the first boundary condition is used. The steeper decline of $`B_z(r,0)`$ with $`r`$ is more noticeable in the cases with large $`C`$ values. For large $`C`$ values we therefore expect much smaller values of $`N/C`$ to arise when the second boundary condition is used. In the right panel of Fig. 30 we plot $`N/C`$ and $`N^{\mathrm{dp}}/C`$ as a function of $`C`$ for the calculations presented in Section 4.4. The ratio $`N/N^{\mathrm{dp}}`$ varies in the range $`0.45`$$`0.003`$. Although this ratio is almost unity for $`C=3,`$ as was the case for the first boundary condition, it becomes progressively smaller for larger $`C`$ values. For large values of $`C`$ we have $`N/C1/C^{3/2}`$ when the second boundary condition is used.
In summary, when $`|B_\phi /B_z|`$ on the disc surface is large the magnetic torque acting on it will be much smaller than that estimated assuming $`B_zB_{\mathrm{dp}},`$ especially when external conditions enable the opening of field lines.
## 5 Discussion
We first summarise our results. For a Keplerian rotation profile with the first boundary condition the maximum inflation observed in our calculations with $`C=3`$ is such that $`R_\mathrm{f}/R_\mathrm{i}=3.2.`$ This is a confirmation of the results of ?). The expansion of the field lines is restricted by the outer boundary in this case. However, for $`|B_\phi /B_z|3`$ at the disc surface the field lines with $`r3R_\mathrm{c}`$ do not vary much once $`R_{\mathrm{ext}}>10.`$ Thus although significantly inflated they will remain closed in the force–free equilibrium. Inflation increases with $`C`$ so that when $`|B_\phi /B_z|=10,`$ a maximum $`R_\mathrm{f}/R_\mathrm{i}5`$ is found for $`R_\mathrm{i}=8.`$
In calculations with $`C=30`$$`120`$ and $`R_{\mathrm{ext}}=11`$ a new effect was found. Inner field lines become deflated rather than inflated with respect to the initial dipolar configuration. We expect the level of deflation, and consequent increase in poloidal magnetic pressure, to depend on the rate of increase of toroidal magnetic pressure, which it has to balance around $`R=R_\mathrm{c}`$. We demonstrated, by varying $`R_{\mathrm{ext}}`$, that the deflation is not related to the outer boundary. We found that the level of field line deflation for $`r<10`$ was restricted when a smoother non Keplerian rotation profile was used. However, some level of deflation could be typical of the field structure in the region near the corotation radius. That applies especially to discs with $`\eta /\nu 1`$ and with $`R_\mathrm{d}R_\mathrm{c}`$ when a Dirichlet outer boundary condition is used.
In all the calculations presented here the equilibrium poloidal field is more nearly radial above the disc surface than the stellar dipole field. The lengthening of the field lines is accompanied by a reduction of $`|B_\phi |/B_\mathrm{p}`$ for increasing $`\mu .`$ We expect that the force–free condition constrains $`|B_\phi |/B_\mathrm{p}`$ to remain of order unity even for $`|B_\phi /B_z|1`$ at the disc surface.
We also performed calculations with a Neumman outer boundary condition that required the disc field to be radial at the outer boundary as it might occur if a wind existed. In these cases, other parameters being equal, the equilibrium field structure was more open. The ratio $`|B_\phi |/B_\mathrm{p}`$ was found to be less than unity for $`\mu <0.1`$ and $`0`$ in most of the corona as expected on open field lines. In the limit of a large twist the poloidal magnetic field tends to be expelled from the disc and consequently the toroidal magnetic field is annihilated, as long as the disc field is permitted to penetrate the outer boundary.
We also found that for large values of $`|B_\phi /B_z|`$ on the disc surface the total magnetic torque is much smaller than that estimated assuming $`B_zB_{\mathrm{dp}}`$. For $`C=120`$ the torque is only $`0.003`$ of the dipole value when the second boundary condition is used and $`0.0130.05`$ times the dipole value when the first boundary condition is used (depending on the form of $`\mathrm{\Omega }`$). Since reasonable values of $`\eta `$ for circumstellar discs tend to produce very large values of $`|B_\phi /B_z|`$ on the disc surface, we believe that the modification of the total torque discussed here will be significant if a force free equilibrium can be attained. We therefore expect significant modification of the magnetic breaking times of T Tauri stars to result from the large twisting of the magnetospheric field lines.
### 5.1 The spin–down of neutron stars
Recent observations of accreting neutron stars (\[Nelson et al. (1997)\]; \[Chakrabatry et al. (1997)\]) suggest that the total torque between star and disc can oscillate producing alternating spin–up and spin–down phases with little evidence of any correlation with the mass accretion rate. ?) have suggested that this may be due to variations in the structure of the magnetosphere.
The work presented here indicates that a large range in the value of the spin–down torque may be obtained depending on the outer boundary condition and effective disc resistivity. In cases favouring large field line inflation and open field lines the spin–down torque resulting from the disc-magnetosphere interaction becomes very small. Thus if the magnetosphere oscillates between such a state and one with smaller inflation, oscillations in the spin–down torque and hence the total torque may be produced. The change from open to closed field lines may also affect the magnetic torque that arises from a wind, if such a wind is produced. Changes in the magnetospheric configuration might occur through reconnection through the disc midplane, variations in outer conditions being more or less favourable for open field configurations, or variations in the internal disc resistivity.
## Acknowledgments
This work was supported by the European Union grant ERBFMRX-CT98-0195. V.A. acknowledges support by the State Scholarships Foundation (IKY) of the Hellenic Republic through a postgraduate studentship. The authors are grateful to Anne Bardou for useful discussions. |
warning/0003/astro-ph0003232.html | ar5iv | text | # COSMIC VELOCITIES 2000: A REVIEW To Appear in Proceedings of the XXXVth Rencontres de Moriond: Energy Densities in the Universe
## 1 Background
The study of cosmic flows emerged as a distinct subfield of cosmology in the late 1970s, spurred by the discovery of the Tully-Fisher (TF) and Faber-Jackson relations for spiral and elliptical galaxies, respectively. With these empirical correlations between galaxy luminosity and internal velocity as distance indicators, one could measure redshift-independent distances for galaxies out to tens or even hundreds of megaparsecs. Both the TF relation and the successors of Faber-Jackson—the $`D_n`$-$`\sigma `$ and Fundamental Plane (FP) relations—yield distances with $`15`$$`20\%`$ accuracy. Although hopes for further improving the accuracy of TF and FP have proved unfounded, these distance indicators have remained the workhorses of cosmic velocity analysis for two decades. Only in recent years have newer and more accurate distance indicator methods—in particular Type Ia Supernovae and Surface Brightness Fluctuations—begun to complement (not supplant) TF and FP, as discussed further below.
The early scientific emphasis in flow studies was on determining the amplitude of Virgo infall (e.g., Tonry & Davis 1981; Aaronson et al. 1982). The Virgo cluster and its environs were then thought to dominate the local flow field. It has since been recognized that the local velocity field is more complex. There are a number of nearby attractors and voids, as well as the tidal effect of distant mass concentrations. The most sophisticated models of the local velocity field now use gravitational instability theory to predict peculiar velocities on the basis of the galaxy density field observed in redshift surveys. The redshift survey most widely used for this purpose is the IRAS redshift survey, both in its older, 1.2 Jy (Fisher et al. 1995) and newer PSCz (Saunders et al. 2000) incarnations. If one assumes IRAS galaxies trace mass and adopts the approximations of linear theory, comparison of predicted velocities with the observed ones constrains the parameter $`\beta _I=\mathrm{\Omega }_m^{0.6}/b_I.`$ Here, $`b_I`$ is the biasing parameter for IRAS galaxies, a measure of whether IRAS galaxies are more ($`b_I>1`$) or less ($`b_I<1`$) clustered than mass. (The subscript is needed because different redshift samples have different clustering amplitudes, and thus different biasing parameters.)
Measurement of $`\beta _I`$ has thus become one of the main thrusts of cosmic flow analysis in the 1990s. Early in the decade, when there was a widespread theoretical prejudice in favor of an $`\mathrm{\Omega }_m=1`$ universe, it was thought peculiar velocities might be the key to proving it. On large scales, the argument went, galaxies should trace the dominant dark matter, and the large-scale peculiar field should therefore reflect the underlying Einstein-de Sitter nature of the universe. The earliest efforts in this direction indeed seemed to bear out this suspicion, finding $`\beta _I1.3,`$ suggestive of $`\mathrm{\Omega }_m=1`$ (Dekel et al. 1993). More recently, however, cosmic flow-based estimates of $`\beta _I`$ have often, though not invariably, produced values in the $`0.4`$$`0.6`$ range, consistent with a low-density cosmology. In §3 I will summarize recent work done on this problem, and explain why I believe the low $`\beta _I`$ values are more likely to be correct.
Another aim of cosmic flow studies first arose serendipitously: efforts to detect bulk flow on very large ($`\stackrel{>}{}100h^1\mathrm{Mpc}`$) scales. This work was propelled by the discovery, in 1987, of a bulk flow stretching across the sky by the “7-Samurai” (7S) group (Dressler et al. 1987). Working with the newly discovered $`D_n`$-$`\sigma `$ relation, the 7S found that elliptical galaxies out to $`3000\mathrm{km}\mathrm{s}^1`$ redshift exhibited fairly uniform Hubble expansion from the vantage point of the Local Group (LG) barycenter. But the LG is known to move at $`630\mathrm{km}\mathrm{s}^1`$ with respect to the Cosmic Microwave Background radiation (CMB), as indicated by the CMB dipole anisotropy. Thus, the implication of the 7S findings was that the ellipticals were streaming at $`600\mathrm{km}\mathrm{s}^1`$ relative to the CMB. If there is any validity to Big Bang cosmology, though, the CMB defines an absolute standard of rest—the “cosmic rest frame,” as it were. The 7S finding thus inaugurated a long-standing puzzle in cosmology: how can large-scale, coherent bulk flows exist in a universe that seems to be so uniform on large scales?
The puzzle deepened in the following decade, with a number of groups confirming a 7S-like bulk flow, and, in several cases, finding that it continued to scales three or four times the 7S volume. The current controversy may be framed as “what is the scale of the largest bulk flow,” or, equivalently, as one of convergence scale: at what distance are the galaxies within a spherical shell finally at rest in the CMB frame? Some astronomers (this author included) have been led to wonder whether such a convergence scale existed; perhaps the CMB-defined “cosmic rest frame” was offset by $`600\mathrm{km}\mathrm{s}^1`$ from the frame in which uniform Hubble expansion is observed—a conjecture which, if correct, would call into question some of the fundamental tenets of cosmology. Fortunately—at least if you like agreement between theory and observation—it now appears that the convergence to the CMB has been detected at a distance of $`50`$$`60h^1\mathrm{Mpc}.`$ Or at least, that is what I will argue in §2, when I summarize results from newly completed surveys.
## 2 The Scale of the Largest Bulk Flows
First we should ask, from the perspective of cosmology, Why are bulk flows interesting? In particular, what does their convergence scale tell us?
The answer lies in the sensitivity of bulk flows to long-wavelength modes of the mass fluctuation power spectrum. Mass conservation in the linear regime of gravitational instability tells us that
$$𝐯=\mathrm{\Omega }_m^{0.6}\delta ,$$
(1)
where $`𝐯`$ is the peculiar velocity vector and $`\delta `$ is the mass density contrast. The corresponding equation in Fourier space is $`𝐤𝐯_k=\mathrm{\Omega }_m^{0.6}\delta _k,`$ where the subscript denotes Fourier transform. Thus, $`|𝐯_k|\delta _k/k,`$ which is to say, long-wavelength perturbations (small $`k`$) have a larger impact on large-scale peculiar velocities than they do on mass fluctuations.
One can flesh out these ideas by calculating the mean square bulk velocity on a scale $`R:`$
$$v^2(R)=\frac{\mathrm{\Omega }_m^{1.2}}{2\pi ^2}_0^{\mathrm{}}P(k)\stackrel{~}{W}^2(kR)𝑑k,$$
(2)
where $`P(k)`$ is the mass fluctuation power spectrum and $`\stackrel{~}{W}^2(kR)`$ is the Fourier transform of a top-hat window of radius $`R.`$ In Figure 1 the rms expected bulk velocity, $`V_{rms}(R)=\sqrt{v^2(R)},`$ is plotted against $`R`$ (left panel) and $`\mathrm{\Omega }_m`$ (right panel) for COBE-normalized CDM power spectra. The left panel assumes a canonical $`\mathrm{\Omega }_m=0.3,`$ $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ universe. For comparison, the rms density fluctuation $`\sigma __M(R)=(\delta M)/M)^2(R)`$ is also plotted. Note that $`\sigma __M`$ drops much more rapidly with scale than does $`V_{rms}.`$ This is a result of the sensitivity of large-scale bulk flow to long-wavelength modes of the power spectrum. An observational consequence is that while redshift surveys have difficulty probing the power spectrum on scales $`\stackrel{>}{}100h^1\mathrm{Mpc},`$ bulk flow studies can in principle do so.
Closer inspection of the left panel also shows why large-amplitude ($`V\stackrel{>}{}500\mathrm{km}\mathrm{s}^1`$), large-scale $`R\stackrel{>}{}60h^1\mathrm{Mpc}`$ are potentially problematic for standard structure-formation scenarios. In the CDM-type models shown here, there simply isn’t enough large-scale power to drive such flows. Stated another way, the CDM universe approaches homogeneity sufficiently rapidly with increasing scale that the coherence scale of bulk flows should be a few tens of megaparsecs at most. Moreover, the right panel shows, perhaps counterintuitively, that boosting the matter density doesn’t help. (This is a consequence of imposing COBE-normalization; normalizing the power spectrum to the cluster abundance does not substantially change our conclusions.)
However, high-amplitude bulk motions on scales $`\stackrel{>}{}100h^1\mathrm{Mpc}`$ are just what were found by three surveys conducted in the early and mid 1990s. The first of these, and the one which has achieved the greatest notoriety, was the Brightest Cluster Galaxy (BCG) survey by Lauer & Postman (1994; LP). LP used the BCGs as standard candles, thus obtaining distances to over 100 Abell clusters out to $`15,000\mathrm{km}\mathrm{s}^1`$ redshift, and found them to be coherently moving relative to the CMB at $`700\mathrm{km}\mathrm{s}^1.`$ The combination of large scale and high amplitude places the LP result far above the expected $`V_{rms}`$ values in Figure 1.
More recently, two large surveys of cluster galaxies appeared to confirm the scale and amplitude (but not the direction) of the LP result. The Streaming Motions of Abell Clusters (SMAC) survey of Hudson et al. (1999) used the FP relation to measure cluster ellipticals to about the same depth as the LP survey, and found a $`600\mathrm{km}\mathrm{s}^1`$ flow, in roughly the same direction as the that found by 7S a decade earlier. And in a Tully-Fisher survey of cluster galaxies in a shell between $`9000`$ and $`13,000\mathrm{km}\mathrm{s}^1,`$ I (Willick 1999, LP10K) found a streaming motion of about $`700\mathrm{km}\mathrm{s}^1`$ in roughly the same direction as SMAC. The above results are summarized in Table I, along with another bulk flow measurement based on more nearby galaxies from the Mark III Catalog of TF and $`D_n`$-$`\sigma `$ data (Willick et al. 1997), as measured by the POTENT algorithm (on which more in § 3).
Table I. Recent Bulk Flow Measurements
| Survey | $`R`$ ($`\mathrm{km}\mathrm{s}^1`$) | $`V_B`$ ($`\mathrm{km}\mathrm{s}^1`$) | Comments |
| --- | --- | --- | --- |
| Lauer-Postman (LP) | 15000 | 700 | BCG |
| Willick (LP10K) | 12000 | 700 | TF |
| Hudson et al. (SMAC) | 14000 | 600 | FP |
| Dekel et al. (POTENT) | 6000 | 350 | TF+$`D_n`$-$`\sigma `$ |
The four results above argue for large bulk flows, but are not fully consistent. The SMAC, LP10K, and POTENT/Mark III flows agree in direction, but the LP flow is nearly orthogonal. Also, the smaller amplitude of the smaller scale POTENT/Mark III measurement is puzzling; one would expect bulk flow amplitude to diminish monotonicaly with scale (Figure 1). Even without evidence to the contrary, then, the above results are less than convincing.
More importantly, however, 1999 saw the announcements of new survey results that contradict the findings of Table 1. The nearby Type Ia Supernova (SN Ia) data have accumulated, and as reported by Riess (2000), the sample of SN Ia distances within $`10,000\mathrm{km}\mathrm{s}^1`$ redshift show no evidence for bulk flow. This is quite important, because the SN Ia data are of a fundamentally different nature than the other distance indicators employed in cosmic flow studies. Also, as shown in the cosmological context, SN Ia have small scatter, $`0.15`$ mag. The EFAR FP survey of Colless et al. (2000) similarly finds no bulk motion on a large scale, and in particular is inconsistent with LP at better than 99% confidence. Similar results have been obtained by Dale, Giovanelli, and coworkers from their extensive TF surveys (as summarized by Dale & Giovanelli 2000). The key findings of these surveys, plus that of the Shellflow survey, to which I turn next, are summarized in Table 2.
Table II. More Recent Bulk Flow Measurements
| Survey | $`R`$ ($`\mathrm{km}\mathrm{s}^1`$) | $`V_B`$ ($`\mathrm{km}\mathrm{s}^1`$) | Comments |
| --- | --- | --- | --- |
| Riess et al. | $`10000`$ | $`0`$ | SN Ia |
| Courteau et al. (SHELLFLOW) | $`6000`$ | $`0`$ | TF |
| Colless et al. (EFAR) | $`10000`$ | $`0`$ | FP |
| Dale & Giovanelli (SFI) | $`6000`$ | $`0`$ | TF |
| Dale & Giovanelli (SCI/SCII) | $`15000`$ | $`0`$ | TF |
A group of us who were active in cosmic velocity measurement and analysis (M. Strauss, S. Courteau, M. Postman, D. Schlegel, and myself) realized in 1995 that a critical issue had not been properly addressed in the then-extant Tully-Fisher surveys: the need for extremely uniform data across the sky. We proposed for and were granted extensive NOAO time for the Shellflow project, a TF survey of a shell of 300 galaxies between $`5000`$ and $`7000\mathrm{km}\mathrm{s}^1`$ redshift. By using identical observational setups from northern and southern hemisphere NOAO telescopes we ensured that data nonuniformity could not produce spurious peculiar velocities.
In a recent paper (Courteau et al. 2000) we reported our main result: the bulk flow of the shell centered at $`6000\mathrm{km}\mathrm{s}^1`$ is $`70_{70}^{+100}\mathrm{km}\mathrm{s}^1,`$ i.e., is consistent with being at rest in the CMB frame. The residuals with respect to a fit that assumes pure Hubble flow in the CMB frame are shown in Figure 2. The points are everywhere consistent with being due to TF scatter, not coherent peculiar velocities. Inflowing and outflowing points are well mixed at all positions, indicating the absence of coherent motions.
Having worked with the Shellflow data myself, I am confident that it is of very good quality, and am convinced that systematic errors have a very small effect, if any on our results. I therefore find the conclusion of insignificant bulk flow highly persuasive. The Shellflow results refer to a specific scale, namely, a $`60h^1\mathrm{Mpc}`$ sphere. They do not directly test the LP, SMAC, and LP10K results of Table 1. However, because abundant evidence indicates that the universe approaches homogeneity monotonically with increasing scale size, I believe that the Shellflow result, if correct, is physically inconsistent with the LP, SMAC, and LP10K findings, and that the latter are therefore not to be taken at face value.
One should not, however, judge the LP, SMAC, and LP10K authors harshly (and I have an obvious reason for hoping you don’t!). Measuring distortions of the Hubble expansion that are at the level of a few percent is a challenging task, given the limitations of the distance indicator techniques we work with. False detections, if that is what they are, will most likely be seen in hindsight as inevitable products of initial efforts at a very difficult measurement.
## 3 The Value of $`\beta _I`$
In considering the scientific implications of large-scale bulk flow surveys, I made no mention of comparison with redshift surveys. That is because the full-sky redshift surveys we have are not reliable enough, at distances $`\stackrel{>}{}150h^1\mathrm{Mpc},`$ to predict what the bulk flows should be on such scales. The situation changes when we talk about the velocity field within $`50h^1\mathrm{Mpc}.`$ At these distances both the peculiar velocity and galaxy density fields are mapped with enough accuracy to do a detailed comparsion. The goals of this comparison are (1) to verify that the two maps are compatible with the gravitational instability paradigm, and (2) assuming they are, to measure $`\beta _I=\mathrm{\Omega }_m^{0.6}/b_I,`$ and if possible to go further and measure $`\mathrm{\Omega }_m`$ itself.<sup>a</sup><sup>a</sup>aIt will be assumed in what follows that, unless otherwise specified, the redshift survey in question is one of IRAS galaxies.
The theoretical bases of the comparison are either of two forms of the linear velocity-density relation. One is the differential form, Eq. (1), which for comparison of observables is written as
$$𝐯(𝐫)=\beta _I\delta _{I,g}(𝐫),$$
(3)
where $`\delta _{I,g}(𝐫)`$ is the galaxy overdensity as determined by an IRAS redshift survey. The other is the corresponding integral form,
$$𝐯(𝐫)=\frac{\beta _I}{4\pi }d^3𝐫^{}\frac{\delta _{I,g}(𝐫^{})(𝐫^{}𝐫)}{\left|𝐫^{}𝐫\right|^3}.$$
(4)
In both Eqs. (3) and (4), the spatial positions $`𝐫`$ are assumed to be measured in velocity units—i.e., they are distances in units of the Hubble velocity. In this way, the Hubble constant itself is removed from the analysis, which is (needless to say) a useful simplification.
Although the two forms of the velocity-density relation are equivalent, they lead to two rather different analytical approaches to measuring $`\beta _I.`$ Because the controversy between the “high” and the “low” (see §1) values of $`\beta _I`$ appears to revolve around this distinction, it is worth taking a moment to understand it. To apply Eq. (3), one needs to map out the three-dimensional velocity field $`𝐯(𝐫),`$ differentiate it, and finally compare it to the galaxy density field to determine $`\beta _I.`$ Because only the radial component of $`𝐯`$ is observable, one first needs an ansatz for “three-dimensionalizing” the inherently one-dimensional velocity data. An elegant approach to this problem was developed by Bertschinger & Dekel (1989; see Dekel 1994 for a review), who argued that the large-scale velocity field should be irrotational and thus expressible as the gradient of a potential function, which could itself be computed by integrating the observed, radial velocities along rays. This algorithm, known as POTENT, thus produces a 3D velocity field, smoothed on a rather large (typically $`12h^1\mathrm{Mpc}`$) scale, which can be differentiated and then used in Eq. (3).
Eq. (4) suggests a different approach. Rather than heavily processing the velocity data, one carries out the indicated integration using the redshift survey data. One thus obtains a predicted peculiar velocity field as a function of $`\beta _I,`$ of which only the radial component, $`u(𝐫)=𝐯(𝐫)𝐫/r,`$ is used in the subsequent analysis. The predicted $`u(𝐫)`$ is compared with the observed radial peculiar velocities from the TF (FP, etc.) data sets; the final estimate of $`\beta _I`$ is that which yields the closest match predictions and observations.
The two approaches to measuring $`\beta _I`$ are known as the density-density (d-d) and velocity-velocity (v-v) comparisons. It is notable that d-d comparisons, all done using POTENT to reconstruct the 3D velocity field, have consistently produced values of $`\beta _I`$ consistent with unity—and thus, to the extent $`b_I`$ is itself not so different from one, implicit estimates of $`\mathrm{\Omega }_m`$ near unity as well. The 1993 paper by Dekel et al. was already mentioned (see §1). A more recent application of POTENT, using improved peculiar velocity data, was that of Sigad et al. (1998), who found $`\beta _I=0.89\pm 0.1.`$
Since about 1995, several v-v alternatives to the POTENT approach have been developed for measuring $`\beta _I.`$ They have differed in the way in which the IRAS (or other) redshift data are used to predict peculiar velocities, and in the way the predicted and observed peculiar velocities are compared. But as v-v methods they have more in common with one another than with POTENT; in particular, the heavy computational work is done with the redshift data, while the TF (FP, etc.) data are used only in a limited, statistical sense. These methods include the Least Action Principle approach of Shaya, Tully, & Peebles (1995), who obtain $`\beta _I0.35\pm 0.1`$), the VELMOD method of Willick et al. (1997) and Willick & Strauss (1998), who obtain $`\beta _I=0.5\pm 0.05,`$ and the ITF method (Davis, Nusser, & Willick 1996) which, applied to the Type Ia Supernova data produced a value of $`\beta _I=0.4\pm 0.1.`$
Figure 3 shows representative results from two v-v analyses. The left panel shows the VELMOD results from Willick & Strauss (1998). The statistic plotted is $`=2\mathrm{ln}P,`$ where $`P`$ is the probability of observing the TF data given the IRAS velocity model. (The calculation is done for the “inverse” TF relation, which is immune to selection biases.) The minimum of the curve occurs at the maximum likelihood value of $`\beta _I,`$ and its curvature yields the $`1\sigma `$ error, as indicated on the plot.
The right panel shows the first use of the new Surface Brightness Fluctuation (SBF) data set in $`\beta `$-measurement, as reported by Blakeslee et al. (2000). As was stated above in connection with SN Ia data, it is essential to carry out these experiments with independent distance indicators, and the SBF method is quite distinct from Tully-Fisher. Moreover, like SN Ia, SBF distances are on average about twice as accurate as TF distances, and have the additional advantage of having well-determined errors. In Figure 3, the SBF data are compared with velocities predicted by IRAS and the Optical Redshift Survey (ORS). The minimum $`\chi ^2`$ is achieved for $`\beta _I=0.44\pm 0.08,`$ consistent with the VELMOD TF results. (The ORS galaxies are more clustered than IRAS galaxies, and therefore yield a lower $`\beta `$ value.)
A truly convincing explanation for this discrepancy between the v-v and d-d $`\beta `$ values has not yet been found. It seems to me, however, that the v-v comparison is a more robust procedure. In it, the intensive data manipulation is done on the redshift survey data, which is by its nature more reliable than the distance indicator data (redshifts errors are fractionally very small; redshift-independent distance errors are always $`20\%`$). In the d-d comparison, it is the data set with much larger scatter and non-Gaussian errors that is subjected to intensive manipulation. In particular, the d-d comparison requires that noisy data first be smoothed, then integrated, and then differentiated in three dimensions. There is ample opportunity in this procedure for errors to propagate. The same noisy data in the v-v studies are left virtually untouched, save for the statistical comparison with their predicted values. This procedure is far more stable. For this reason I consider the low $`\beta `$ values derived from the v-v analyses to be more reliable.
## 4 Summary
I have argued that the two major controversies in cosmic flow analysis have been largely resolved in the last few years. With regard to bulk flows, most recent surveys show convergence to the CMB frame by a distance of $`60h^1\mathrm{Mpc}.`$ A corollary is that the observed bulk motions do not require more large-scale power than is provided by COBE-normalized CDM density fluctuation spectra. With regard to the value of $`\beta _I=\mathrm{\Omega }_m^{0.6}/b_I,`$ a number of independent analyses now suggest a low value, $`\beta _I0.4`$$`0.5.`$ If the IRAS galaxies are nearly unbiased with respect to mass, a reasonable if not airtight hypothesis, these $`\beta `$-values imply a density parameter $`\mathrm{\Omega }_m0.2`$$`0.3.`$
A word of caution is in order, however. The above conclusions represent a consensus view, not a unanimous one. The bulk flow detections listed in the first three rows of Table 1 have not been in any sense “refuted,” which is to say as far as we know there is nothing wrong with the data. New surveys, such as FP200 (see http://astro.uwaterloo.ca/$``$mjhudson/fp200 for details) will, it is hoped, settle the issue definitively. Similarly, while a majority of recent velocity-density comparisons favor low $`\beta _I,`$ the reason for the discrepant POTENT result, $`\beta _I0.9,`$ is not well understood. Tests of methods such as POTENT and VELMOD using N-body simulations are under way, and may clarify things. Moreover, the coming decade will bring much larger TF data sets from the DENIS and 2MASS infrared surveys. As always, these new data sets, if they live up to their promise, are our best hope for putting any remaining contoversy to rest.
## Acknowledgments
I am grateful to my collaborators on the Shellflow project, Stéphane Courteau, Michael Strauss, Marc Postman, and David Schlegel, as well as to Michael Hudson, Marc Davis, Avishai Dekel, John Tonry and Alan Dressler for enlightening discussions over the last several years. Special thanks go to Stéphane Courteau for organizing the highly successful Cosmic Flows 99 conference in Victoria, B.C. last summer. My research is supported by a Cottrell Scholarship of Research Corporation, NSF grant AST96-17188, and a Terman Fellowship from Stanford University. |
warning/0003/hep-ph0003016.html | ar5iv | text | # Collisional Energy Loss of Fast Charged Particles in Relativistic Plasmas
## Abstract
Following an argument by Kirzhnits we rederive an exact expression for the energy loss of a fast charged particle in a relativistic plasma using the quantum field theoretical language. We compare this result to perturbative calculations of the collisional energy loss of an energetic electron or muon in an electron-positron plasma and of an energetic parton in the quark-gluon plasma.
Keywords: Energy loss; Relativistic plasmas; Thermal Field Theory
The energy loss of a fast charged particle in a medium is a well studied subject . Recently the energy loss of energetic particles, such as leptons and partons, in relativistic plasmas has attracted great interest. In relativistic heavy ion collisions the energy loss of a high energy quark or gluon coming from primary hard collisions in the fireball may serve as a signature for the quark-gluon plasma formation . In Supernovae explosions the energy loss of neutrinos, having a weak charge, in the plasma surrounding the stellar core might be an important mechanism for triggering the explosion .
The total energy loss of a particle in a medium can be decomposed into a collisional and a radiative contribution. While the first one originates from the energy transfer to the medium particles, the latter one is caused by radiation from the fast particle. Here we want to consider only the collisional component. Whereas the radiative energy loss dominates in the case of partons or charged leptons , the collisional one is dominant for neutrinos in a Supernova plasma due to the small coupling of the neutrinos to the medium.
In quantum field theory the collisional energy loss per unit length is defined as
$$\frac{dE}{dx}=\frac{1}{v}𝑑\mathrm{\Gamma }\omega ,$$
(1)
where $`v`$ is the velocity of the incident particle with energy $`E`$ and $`\omega =EE^{}`$ the energy transfer to the medium. The interaction rate $`\mathrm{\Gamma }`$ can be calculated either from the matrix element of the process responsible for the energy loss or equivalently from the imaginary part of the self energy of the particle with four momentum $`P=(E,𝐩)`$, and mass $`M`$, ($`p=|𝐩|`$)
$$\mathrm{\Gamma }(E)=\frac{1}{2E}[1n_F(E)]tr[(P/+M)Im\mathrm{\Sigma }(E,p)],$$
(2)
where $`n_F(E)=1/[\mathrm{exp}(E/T)+1]`$ is the Fermi distribution in the case of a fermion propagating through a plasma of temperature $`T`$. In the following we restrict ourselves first to electrons or muons with high energies $`ET`$ in an electron-positron plasma. Furthermore we assume first only small momentum and energy transfers, $`\omega `$, $`kT`$. Assuming a one-loop approximation for $`\mathrm{\Sigma }`$ but allowing for the most general photon propagator, indicated by the blob in Fig.1, we find
$$\left(\frac{dE}{dx}\right)_{soft}=\frac{e^2}{2\pi v^2}_0^k^{}𝑑kk_{vk}^{vk}𝑑\omega [1+n_B(\omega )]\left[\rho _l(\omega ,k)+\left(v^2\frac{\omega ^2}{k^2}\right)\rho _t(\omega ,k)\right],$$
(3)
where $`k^{}T`$ is the separation scale, decomposing the energy loss into a soft and a hard part. $`n_B(\omega )=1/[\mathrm{exp}(\omega /T)1]`$ is the Bose distribution and $`\rho _{l,t}`$ are the spectral functions of the full photon propagator, defined as
$$D_{l,t}(k_0,k)=_{\mathrm{}}^{\mathrm{}}𝑑\omega \frac{\rho _{l,t}(\omega ,k)}{k_0\omega +i\epsilon }.$$
(4)
At finite temperature the photon propagator has only two independent components , given in Coulomb gauge by the longitudinal and transverse propagators
$`D_l(k_0,k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2\mathrm{\Pi }_l(k_0,k)+i\epsilon }},`$ (5)
$`D_t(k_0,k)`$ $`=`$ $`{\displaystyle \frac{1}{k_0^2k^2\mathrm{\Pi }_t(k_0,k)+i\epsilon }},`$ (6)
where $`\mathrm{\Pi }_{l,t}`$ are the longitudinal and transverse components of the polarization tensor. It should be noted that the soft collisional energy loss, discussed here, follows according to (3) only from the exchange of one dressed space-like ($`\omega ^2k^2<0`$) photon from the particle to the medium according to linear response theory. However, the medium particles may undergo further interactions. The physical process corresponding to the imaginary part of the self energy of Fig1. can be found by using cutting rules. An example is shown in Fig.2. There is no diagram, where to or more photons are emitted from the fast particle, as it is the case e.g. for bremsstrahlung. In the case of a neutrino, however, diagrams containing two gauge boson lines are suppressed anyway.
The spectral functions can be expressed by the imaginary part of the photon propagator according to
$$\rho _{l,t}(\omega ,k)=\frac{1}{\pi }ImD_{l,t}(\omega ,k).$$
(7)
For soft energy transfers $`\omega T`$, as we assumed above, the photon distribution can be expanded, leading to
$$1+n_B(\omega )\frac{T}{\omega }+\frac{1}{2}.$$
(8)
Substituting (8) into (3) only the second term in (8) contributes since the spectral functions are odd functions of $`\omega `$ .
Alternatively the soft energy loss can also be derived from classical plasma physics arguments. It follows from the induced electric field of the fast charged particle in the plasma, which reacts on the incident particle by the Lorentz force, causing the energy loss . This process is known as the Fermi density effect . Introducing the dielectric functions of the medium, the soft energy loss can be written as
$$\left(\frac{dE}{dx}\right)_{soft}=\frac{e^2}{4\pi ^2v^2}_0^k^{}𝑑kk_{vk}^{vk}𝑑\omega \omega \left[k^2Imϵ_l(\omega ,k)^1+\left(v^2\frac{\omega ^2}{k^2}\right)Im(\omega ^2ϵ_t(\omega ,k)k^2)^1\right].$$
(9)
This expression is equivalent to (3), since the dielectric functions are related to the polarization tensor via
$`ϵ_l(\omega ,k)`$ $`=`$ $`1{\displaystyle \frac{\mathrm{\Pi }_l(\omega ,k)}{k^2}},`$ (10)
$`ϵ_t(\omega ,k)`$ $`=`$ $`1{\displaystyle \frac{\mathrm{\Pi }_t(\omega ,k)}{\omega ^2}}`$ (11)
and therefore also to the spectral functions, which are given by the imaginary part of the propagators (6), by
$`Imϵ_l(\omega ,k)^1=\pi k^2\rho _l(\omega ,k),`$ (12)
$`Im(\omega ϵ_t(\omega ,k)k^2)^1=\pi \rho _t(\omega ,k),`$ (13)
where only the discontinuous part of the spectral functions coming from the imaginary part of the polarization tensor contributes.
Now we want to derive an exact result for the soft collisional energy loss using a generalized Kramers-Kronig relation and the asymptotic behavior of the dielectric functions. For this purpose we introduce the response function of the medium $`R(k_0,k)=R_l(k_0,k)+R_t(k_0,k)`$, given by
$`R_l(k_0,k)`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_l(k_0,k)}},`$ (14)
$`R_t(k_0,k)`$ $`=`$ $`{\displaystyle \frac{k^2k_0^2}{k^2k_0^2ϵ_t(k_0,k)}}.`$ (15)
Using (6) and (11) we obtain
$$R(k_0,k)=k^2D_l(k_0,k)+(k_0^2k^2)D_t(k_0,k).$$
(16)
Replacing the spectral functions in (3) or the dielectric functions in (9) by the quantity $`R`$, making the substitution $`kq=\sqrt{k^2\omega ^2}`$, i.e. introducing the magnitude of the four momentum of the exchanged photon, and using $`ImR(\omega )=ImR(\omega )`$ we find
$$\left(\frac{dE}{dx}\right)_{soft}=\frac{e^2}{2\pi ^2}_0^q^{}𝑑qq_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,\sqrt{q^2+\omega ^2})}{q^2+\omega ^2}.$$
(17)
Here we restricted ourselves to ultrarelativistic particles, $`v=1`$, and $`q^{}T`$. Eq. (17) agrees with Ref., if we replace there $`Q^2`$ by $`e^2/4\pi `$.
The response function $`R`$ fulfills the following Kramers-Kronig relation
$$R(k_0,k)=\stackrel{~}{R}+\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,k)}{\omega ^2k_0^2i\epsilon },$$
(18)
which can be shown to be equivalent to the definition of the spectral functions (4), if we use $`\rho _{l,t}(\omega )=\rho _{l,t}(\omega )`$. Here $`\stackrel{~}{R}=lim_{k_0\mathrm{}}R(k_0,k)lim_{k_0\mathrm{}}1/k_0^2=0`$. The relation (18) can be generalized to the so-called Leontovich relation exploiting causality, from which one obtains
$$R(k_0,\sqrt{k^2+k_0^2})=R_{\mathrm{}}+\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,\sqrt{k^2+\omega ^2})}{\omega ^2k_0^2i\epsilon },$$
(19)
where $`R_{\mathrm{}}=lim_{k_0\mathrm{}}R(k_0,\sqrt{k^2+k_0^2})`$.
The $`\omega `$-integral
$$I=\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,\sqrt{q^2+\omega ^2})}{q^2+\omega ^2}$$
(20)
appearing in the energy loss (17) agrees with the integral on the right hand side of the Leontovich relation, if we replace $`\omega `$ by $`iq`$ and $`\sqrt{k^2+\omega ^2}`$ by 0, i.e. $`k^2=q^2`$, in (19). Therefore we can write
$$I=R(iq,0)R_{\mathrm{}}.$$
(21)
Since the longitudinal and the transverse dielectric functions are identical at zero momentum , $`ϵ_l(k_0,0)=ϵ_t(k_0,0)`$, we have $`R(iq,0)=0`$. $`R_{\mathrm{}}`$ is related to the high frequency and momentum limit of the dielectric functions, which agrees with the vacuum result $`ϵ_l=1`$. For the transverse part we have to consider corrections to the vacuum value. From the second equation of (11) we get
$$\underset{k_0\mathrm{}}{lim}ϵ_t(k_0,\sqrt{q^2+k_0^2})=1\frac{\omega _0^2}{k_0^2},$$
(22)
where
$$\omega _0^2\underset{k_0\mathrm{}}{lim}\mathrm{\Pi }_t(k_0,\sqrt{q^2+k_0^2})$$
(23)
Using the Kramers-Kronig relation for the transverse dielectric functions it can be shown , that $`\omega _0`$ is independent of $`q`$. It can be considered as the effective thermal mass of the transverse high frequency plasma excitations, which is given by $`\omega _0^2=e^2n1/\mathrm{\Omega }`$ in the relativistic limit . Here $`n`$ is the number density of the medium and $`\mathrm{\Omega }`$ the energy of the plasma particles. In the non-relativistic limit $`\omega _0`$ is identical to the plasma frequency .
Using (15) together with the high frequency and momentum limit of the dielectric functions we get
$$I=R_{\mathrm{}}=\frac{\omega _0^2}{q^2+\omega _0^2}.$$
(24)
Combining this result for $`I`$ with (17) we end up with
$$\left(\frac{dE}{dx}\right)_{soft}=\frac{e^2}{4\pi }\omega _0^2\mathrm{ln}\frac{q^{}}{\omega _0},$$
(25)
where we assumed $`q^{}\omega _0`$.
The unknown parameter $`\omega _0`$ following from the full transverse polarization tensor serves as an infrared cutoff for the photon exchange. Since the total collisional energy loss has to be independent of the arbitrary separation scale $`q^{}`$, the hard part has to assume the form
$$\left(\frac{dE}{dx}\right)_{hard}=\frac{e^2}{4\pi }\omega _0^2\mathrm{ln}\frac{q_{max}}{q^{}},$$
(26)
where $`q_{max}`$ is proportional to the maximum energy transfer, i.e. $`q_{max}\sqrt{ET}`$ in the relativistic limit $`E\mathrm{\Omega }`$ , which we have considered here. Note that the hard contribution to the energy loss contains besides $`t`$-channel diagrams, as the one in Fig.2, also $`s`$\- and $`u`$-channel ones, which, however, do not contribute to the leading logarithm.
Hence we obtained a very simple expression
$$\frac{dE}{dx}=\frac{e^2}{4\pi }\omega _0^2\mathrm{ln}\frac{q_{max}}{\omega _0},$$
(27)
for the exact result of the collisional energy loss, independent of any approximation to the full photon propagator or the dielectric functions of the medium, respectively. To logarithmic accuracy the final result just depends on the parameter $`\omega _0`$.
As an example we consider the high temperature limit of the energy loss. There it can be calculated to leading order perturbation theory using the Hard Thermal Loop (HTL) resummation technique . Computing the soft energy loss using the HTL resummed photon propagator in Fig.1 and calculating the hard part from the tree level scattering matrix elements one finds in the limit $`v=1`$
$$\frac{dE}{dx}=\frac{e^2}{4\pi }\omega _0^2\left(\mathrm{ln}\frac{\sqrt{ET}}{\omega _0}+0.120\right),$$
(28)
where $`\omega _0^2=3m_\gamma ^2/2`$. The thermal photon mass $`m_\gamma `$, which is equivalent to the plasma frequency, is given by $`m_\gamma =eT/3`$. Indeed $`\omega _0`$ is given by the high frequency and momentum limit (23) of the transverse HTL polarization tensor
$$\mathrm{\Pi }_t^{HTL}(k_0,k)=\frac{3}{2}m_\gamma ^2\frac{\omega ^2}{k^2}\left[1\left(1\frac{k^2}{k_0^2}\right)\frac{k_0}{2k}\mathrm{ln}\frac{k_0+k}{k_0k}\right]$$
(29)
confirming the general result (27) in the HTL limit. Also $`\omega _0^2=e^2n1/\mathrm{\Omega }`$ holds for the HTL case, where
$$m_\gamma ^2=\frac{4e^2}{3\pi ^2}_0^{\mathrm{}}𝑑kkn_F(k),$$
(30)
because
$$n=4\frac{d^3k}{(2\pi )^3}n_F(k),\frac{1}{\mathrm{\Omega }}=\frac{\frac{d^3k}{(2\pi )^3}\frac{1}{k}n_F(k)}{\frac{d^3k}{(2\pi )^3}n_F(k)}.$$
(31)
In the case of the collisional energy loss of quarks or gluons in a quark-gluon plasma we simply have to replace the factor $`e^2`$ in (27) by $`C_Fg^2=4g^2/3`$ for quarks and by $`C_Ag^2=3g^2`$ for gluons, where $`g`$ is the strong coupling constant. Furthermore $`\omega _0^2`$ is now the high frequency and momentum limit of the transverse gluon polarization tensor, which is given by $`3m_g/2`$, where the effective gluon mass reads $`m_g^2=g^2T^2(1+n_f/6)/3`$ in a QGP containing $`n_f`$ thermalized quark flavors. Using the HTL limit for $`\omega _0`$ in (27) we reproduce again the result obtained from an explicit HTL resummed calculation within the logarithmic approximation. Another application of (27) has been discussed in Ref. in connection with the neutrino energy loss in matter.
Summarizing, we have shown that different definitions of the collisional energy loss, based either on quantum field theory or on plasma physics, are equivalent. Translating the arguments, based on a generalized Kramers-Kronig relation, given by Kirzhnits to a quantum field theoretical language, using self energies, propagators and their spectral functions, we gave an exact result for the collisional energy loss (27). Within the logarithmic approximation it contains the effective mass of the high frequency transverse plasma mode as the only parameter. Assuming that this mass is given approximately by the high temperature result, we obtain an simple estimate for the collisional energy loss for energetic electrons and muons in a QED plasma and for partons in the quark-gluon plasma. Finally we showed that the perturbative result for the collisional energy loss, obtained within the HTL resummation method, is in agreement with the general result found by Kirzhnits.
ACKNOWLEDGMENTS
The author is grateful to G. Raffelt for drawing his attention to the paper by D.A. Kirzhnits and for helpful discussions and to the Max-Planck-Institut für Physik (Werner-Heisenberg-Institut) for their hospitality. |
warning/0003/math-ph0003017.html | ar5iv | text | # Quadratic Poisson algebras for two dimensional classical superintegrable systems and quadratic associative algebras for quantum superintegrable systems.
## I Introduction
In classical mechanics, integrable system is a system possessing more constants of motion in addition to the energy. A comprehensive review of the two-dimensional integrable classical systems is given by Hietarinta , where the space was assumed to be flat. The case of non flat space is under current investigation . An interesting subset of the totality of integrable systems is the set of systems, which possess a maximum number of integrals, these systems are termed as superintegrable ones. The Coulomb and the harmonic oscillator potentials are the most familiar classical superintegrable systems, whose their quantum counterpart has nice symmetry properties, which are described by the $`so(N+1)`$ and $`su(N)`$ Lie algebras respectively.
The Hamiltonian of a classical system is generally a quadratic function of the momenta. In the case of the flat space, all the known two dimensional superintegrable systems with quadratic integrals of motion are simultaneously separable in more than two orthogonal coordinate systems . The integrals of motion of a two dimensional superintegrable system in flat space close in a restrained classical Poisson algebra . The study of the quadratic Poisson algebras is a matter under investigation, related to several branches of physics as: the solution of the classical Yang \- Baxter equation , the two dimensional superintegrable systems in flat space or on the sphere , the statistics or the case of ”exactly solvable” classical problems .
The quantization of a classical integrable system corresponds generally to a quantum integrable system. The mechanism of quantum deformation of a classical system to a quantum one is not fully understood. Initially the problem of quantization of classical superintegrable system was viewed as a relatively simple and somehow trivial problem , but several authors have proved that this quantization procedure has to add correction terms to the integrals of motion or to the Hamiltonian, these correction terms are of order $`𝒪(\mathrm{}^2)`$ . The result of the quantum deformation of a superintegrable system is realized by the shift of the classical Poisson algebra to some quantum polynomial associative algebra. The same fact is true in the case of quadratic Poisson algebra corresponding to the Yang - Baxter equation , which is turned to a quantum quadratic associative algebra with four generators. The same idea was discussed in reference , where the classical problems, which are expressed by a quadratic Poisson algebra are mapped to quantum ones described by the corresponding quantum operator quadratic algebra. The same shift is indeed true for the superintegrable systems, where the classical ones correspond to the quantum ones and the classical quadratic Poisson algebra is mapped to a quadratic associative algebra.
In this paper we show that the deformation of the classical Poisson algebra to a quadratic associative algebra implies a deformation of the parameters of the quadratic algebra. The general form of the quadratic algebras, which are encountered in the case of the two dimensional quantum superintegrable systems, is investigated. In references was conjectured that, the energy eigenvalues correspond to finite dimensional representations of the latent quadratic algebras. Granovskii et al in studied the representations of the quadratic Askey - Wilson algebras $`QAW(3)`$. Using there the proposed ladder representation, the finite dimensional representations are calculated and this method was applied to several superintegrable systems . Another method for calculating the finite dimensional representations is the use of the deformed oscillator algebra and their finite dimensional version which are termed as generalized deformed parafermionic algebras. The main task of this paper is to reduce the calculations of eigenvalues to a system of two algebraic equations with two parameters to be determined. These equations are universal equations, which are valid of all superintegrable systems, with quadratic integrals of motion.
This paper is organized as follows: In section II the general form of the quadratic Poisson algebra for a two dimensional system with quadratic integrals of motion is derived. In section III the special form of the Poisson algebra of the known two dimensional superintegrable systems in flat space is written. In section IV the quantum version of the Poisson quadratic algebra is studied. The deformed parafermionic oscillator algebra is reviewed and the oscillator realization of the quadratic algebras is realized. The finite dimensional representations of the quadratic algebras are generated by using the technique of deformed parafermionic algebras. The problem is reduced to the solution of a system of two algebraic equations in section VI. In section VII the energy eigenvalues of all the known superintegrable systems in the flat two dimensional space are determined by solving the appropriate algebraic equations. Finally in section VIII there is a discussion of the results of this paper.
## II Quadratic Poisson Algebras
Let consider a two dimensional superintegrable system. The general form of the Hamiltonian is:
$$H=a(q_1,q_2)p_1^2+2b(q_1,q_2)p_1p_2+c(q_1,q_2)p_2^2+V(q_1,q_2)$$
(1)
this Hamiltonian is a quadratic form of the momenta. The system is superintegrable, therefore there are two additional integrals of motion $`A`$ and $`B`$. In that section, we consider that, these integrals of motion are quadratic functions of the momenta, i.e. they are given by the general forms:
$$\begin{array}{cc}\hfill A=& A(q_1,q_2,p_1,p_2)=c(q_1,q_2)p_1^2+2d(q_1,q_2)p_1p_2+e(q_1,q_2)p_2^2+\hfill \\ & +f(q_1,q_2)p_1+g(q_1,q_2)p_2+Q(q_1,q_2)\hfill \end{array}$$
The integral of motion $`B`$ is assumed to be indeed a quadratic form, which is analogous to above one.
$$\begin{array}{cc}\hfill B=& B(q_1,q_2,p_1,p_2)=h(q_1,1_2)p_1^2+2k(q_1,q_2)p_1p_2+l(q_1,q_2)p_2^2+\hfill \\ & +m(q_1,q_2)p_1+n(q_1,q_2)p_2+S(q_1,q_2)\hfill \end{array}$$
By definition the following relations are satisfied:
$$\{H,A\}_P=\{H,B\}_P=0$$
(2)
where $`\{.,.\}_P`$ is the usual Poisson bracket.
From the integrals of motion $`A,B`$, we can construct the integral of motion:
$$C=\{A,B\}_P$$
(3)
The integral of motion $`C`$ is not a new independent integral of motion, which is a cubic function of the momenta. The integral $`C`$ is not independent from the integrals $`H,A`$ and $`B`$ as it will be shown later. The fact that, the integral $`C`$ is a cubic function of momenta, implies the impossibility of expressing $`C`$ as a polynomial function of the other integrals, which are quadratic functions of momenta. Starting from the integral of motion $`C`$, we can construct the (non independent) integrals $`\{A,C\}_P`$ and $`\{B,C\}_P`$. These integrals are quartic functions of the momenta, i.e. functions of fourth order. Therefore these integrals could be expressed as quadratic combinations of the integrals $`H,A`$, and $`B`$. Therefore the following relations are assumed to be valid:
$$\{A,C\}_P=\alpha A^2+\beta B^2+2\gamma AB+\delta A+ϵB+\zeta $$
(4)
and
$$\{B,C\}_P=aA^2+bB^2+2cAB+dA+eB+z$$
(5)
We can take appropriate a linear combination of the integrals $`A`$ and $`B`$ and we can always consider the case $`\beta =0`$.
The Jacobi equality for the Poisson brackets induces the relation
$$\{A,\{B,C\}_P\}_P=\{B,\{A,C\}_P\}_P$$
The following relations
$$b=\gamma ,c=\alpha \text{and}e=\delta $$
must be satisfied.
The integrals $`A,B`$ and $`C`$ satisfy the quadratic Poisson algebra:
$$\begin{array}{ccc}\hfill \{A,B\}_P& =& C\hfill \\ \hfill \{A,C\}_P& =& \alpha A^2+2\gamma A,B+\delta A+ϵB+\zeta \hfill \\ \hfill \{B,C\}_P& =& aA^2\gamma B^22\alpha AB+dA\delta B+z\hfill \end{array}$$
(6)
where $`\alpha ,\gamma ,a`$ are constants and
$$\begin{array}{c}\delta =\delta (H)=\delta _0+\delta _1H\hfill \\ ϵ=ϵ(H)=ϵ_0+ϵ_1H\hfill \\ \zeta =\zeta (H)=\zeta _0+\zeta _1H+\zeta _2H^2\hfill \\ d=d(H)=d_0+d_1H\hfill \\ z=z(H)=z_0+z_1H+z_2H^2\hfill \end{array}$$
where $`\delta _i,ϵ_i,\zeta _i,d_i`$ and $`z_i`$ are constants. The associative algebra, whose the generators satisfy equations (6), is the general form of the closed Poisson algebra of the integrals of superintegrable systems with integrals quadratic in momenta.
The quadratic Poisson algebra (6) possess a Casimir which is a function of momenta of degree 6 and it is given by:
$$\begin{array}{cc}\hfill K=& C^22\alpha A^2B2\gamma AB^22\delta AB\hfill \\ & ϵB^22\zeta B+\frac{2}{3}aA^3+dA^2+2zA=\hfill \\ \hfill =& k_0+k_1H+k_2H^2+k_3H^3\hfill \end{array}$$
(7)
Obviously
$$\{K,A\}_P=\{K,B\}_P=\{K,C\}_P=0$$
Therefore the integrals of motion of a superintegrable two dimensional system with quadratic integrals of motion close a constrained classical quadratic Poisson algebra (6), corresponding to a Casimir equal at most to a cubic function of the Hamiltonian (7).
In the general case of a superintegrable system the integrals are not necessarily quadratic functions of the momenta, but rather polynomial functions of the momenta. The case of the systems with a quadratic and a cubic integral of motion are recently studied by Tsiganov . The general form of the Poisson algebra of generators $`A,B,`$ and $`C`$ is characterized by a polynomial function $`h(A,B)`$, which satisfy the following equations:
$$\begin{array}{ccc}\hfill \{A,B\}_P& =& C\hfill \\ \hfill \{A,C\}_P& =& h/B\hfill \\ \hfill \{B,C\}_P& =& h/A\hfill \end{array}$$
(8)
and the Casimir of the algebra is given by
$$K=K(H)=C^22h(A,B),\{K,A\}_P=\{K,B\}_P=0$$
(9)
where $`h(A,B)`$ is a polynomial function of the integrals of motion $`A`$ and $`B`$. In the case of the quadratic Poisson algebra (6) the form of the function $`h(A,B)`$ is given by equation (7):
$$\begin{array}{cc}\hfill h(A,B)=& \frac{a}{3}A^3+\alpha A^2B+\gamma AB^2\hfill \\ & \frac{d}{2}A^2+\delta AB+\frac{ϵ}{2}B^2\hfill \\ & zA+\zeta B\hfill \end{array}$$
In the general case of a two dimensional superintegrable system, with quadratic Hamiltonian, one integral $`A`$ of order $`m`$ in momenta and one integral $`B`$ of order $`n`$ ($`nm`$), the general form of the function $`h(A,B)`$ can be given by the general form:
$$h(A,B)=h_0(A)+h_1(A)B+h_2(A)B^2$$
where $`h_i(A)`$ are polynomials of the integrals $`A`$ and $`H`$. The proof of this assumption is based on the dependence of the integrals of motion on the momenta. For simplicity reasons, the proof of this proposition will not be given here.
## III Poisson algebras for superintegrable systems
Let consider the superintegrable systems with quadratic integrals of motion, these potentials are given by several authors starting from different but comparable points of view. In references the integrals of motion are generated by solving the Darboux conditions for integrability of quaratic integrals. In the Hamilton - Jacobi equation is solved by separation of variables nad the two dimensional Hamiltonians which are separable in more than one coordinate system are obtained. The separation of variables is essential for solving the quantum counterpart of the superintegrable system and the solution of the associate Schrödinger equations is given in . Using this method the quantum superintegrable systems have been solved on the sphere and the hyperboloid. From classical point of view the super integrable are given in , while the case of a pseudo Euclidean kinetic term has been studied in . The extension on the systems with a quadratic and a cubic integral of motion is sytematized in .
In this section we consider the case superintegrable systems gin in ref , because in the next sections we study the quantum counterparts of these potentials. In this paper the following superintegrable systems are considered:
Potential i):
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\omega ^2r^2+\frac{\mu _1}{x^2}+\frac{\mu _2}{y^2}\right)$$
This potential has the following independent integrals of motion:
$$A=p_x^2+\omega ^2x^2+\frac{\mu _1}{x^2}$$
and
$$B=\left(xp_yyp_x\right)^2+r^2\left(\frac{\mu _1}{x^2}+\frac{\mu _2}{y^2}\right)$$
The constants, which characterize the corresponding quadratic algebra (6), are given by:
$$\begin{array}{ccc}\alpha =8,\hfill & \gamma =0,\hfill & \delta =16H,\hfill \\ ϵ=16\omega ^2,\hfill & \zeta =16(\mu _1+\mu _2)\omega ^2\hfill & \\ a=0,\hfill & d=0,\hfill & z=16(\mu _2\mu _1)\omega ^2\hfill \end{array}$$
the value of the Casimir (7) is:
$$K=16\left((\mu _2\mu _1)^2\omega ^2+4\mu _1H^2\right)$$
Potential ii):
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\omega ^2\left(4x^2+y^2\right)+\frac{\mu }{y^2}\right)$$
This potential has the following independent integrals of motion:
$$A=p_x^2+4\omega ^2x^2$$
and
$$B=\left(xp_yyp_x\right)p_y+\frac{\mu x}{y^2}\omega ^2xy^2$$
The constants, which characterize the corresponding quadratic algebra (6), are given by:
$$\begin{array}{ccc}\alpha =0,\hfill & \gamma =0,\hfill & \delta =0,\hfill \\ ϵ=16\omega ^2,\hfill & \zeta =0,\hfill & \\ a=6,\hfill & d=16H,\hfill & z=8\mu _2\omega ^28H^2\hfill \end{array}$$
the value of the Casimir (7) is:
$$K=0$$
Potential iii):
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\frac{k}{r}+\frac{1}{r}\left(\frac{\mu _1}{r+x}+\frac{\mu _2}{rx}\right)\right)$$
This potential has the following independent integrals of motion:
$$A=(xp_yyp_x)^2+r\left(\frac{\mu _1}{r+x}+\frac{\mu _2}{rx}\right)$$
and
$$B=\left(xp_yyp_x\right)p_y\frac{\mu _1}{2r}\frac{rx}{r+x}+\frac{\mu _2}{2r}\frac{r+x}{rx}+\frac{kx}{2r}$$
The constants, which characterize the corresponding quadratic algebra (6), are given by:
$$\begin{array}{ccc}\alpha =0,\hfill & \gamma =2,\hfill & \delta =0,\hfill \\ ϵ=0,\hfill & \zeta =k(\mu _1\mu _2),\hfill & \\ a=0,\hfill & d=8H,\hfill & z=4(\mu _1+\mu _2)Hk^2/2\hfill \end{array}$$
the value of the Casimir (7) is:
$$K=2(\mu _1\mu _2)^2Hk^2(\mu _1+\mu _2)$$
Potential iv):
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\frac{k}{r}+\mu _1\frac{\sqrt{r+x}}{r}+\mu _2\frac{\sqrt{rx}}{r}\right)$$
This potential has the following independent integrals of motion:
$$A=(yp_xxp_y)p_y+\frac{\mu _1(rx)\sqrt{r+u}}{2r}\frac{\mu _2(r+x)\sqrt{ru}}{2r}\frac{kx}{2r}$$
and
$$B=\left(xp_yyp_x\right)p_x\frac{\mu _1x\sqrt{ru}}{2r}+\frac{\mu _2x\sqrt{r+u}}{2r}\frac{ky}{2r}$$
The constants, which characterize the corresponding quadratic algebra (6), are given by:
$$\begin{array}{ccc}\alpha =0,\hfill & \gamma =0,\hfill & \delta =0,\hfill \\ ϵ=2H,\hfill & \zeta =\mu _1\mu _2/2,\hfill & \\ a=0,\hfill & d=2H,\hfill & z=\frac{\mu _1^2\mu _2^2}{4}\hfill \end{array}$$
the value of the Casimir (7) is:
$$K=k^2H/2k(\mu _1^2+\mu _2^2)/4$$
## IV The quadratic associative algebra
The quantum counterparts of the classical systems, which have been studied in section II, are quantum superintegrable systems. The quadratic classical Poisson algebra (6) possesses a quantum counterpart, which is a quadratic associative algebra of operators. The form of the quadratic algebra is similar to the classical Poisson algebra, the involved constants are generally functions of $`\mathrm{}`$ and they should coincide with the classical constants in the case $`\mathrm{}0`$. Let consider the quadratic associative algebra generated by the generators $`\{A,B,C\}`$, which satisfy the commutation relations
$$\begin{array}{c}[A,B]=C\hfill \\ [A,C]=\alpha A^2+\beta B^2+\gamma \{A,B\}+\delta A+ϵB+\zeta \hfill \\ [B,C]=aA^2+bB^2+c\{A,B\}+dA+eB+z\hfill \end{array}$$
(10)
After rotating the generators $`A`$ and $`B`$, we can always consider the case $`\beta =0`$.
The Jacobi equality for the commutator induces the relation
$$[A,[B,C]]=[B,[A,C]]$$
the following relations
$$b=\gamma ,c=\alpha \text{and}e=\delta $$
must be satisfied, and consequently the general form of the quadratic algebra (10) can be explicitly written as follows:
$$[A,B]=C$$
(11)
$$[A,C]=\alpha A^2+\gamma \{A,B\}+\delta A+ϵB+\zeta $$
(12)
$$[B,C]=aA^2\gamma B^2\alpha \{A,B\}+dA\delta B+z$$
(13)
The Casimir of this algebra is given by:
$$\begin{array}{cc}\hfill K=& C^2\alpha \{A^2,B\}\gamma \{A,B^2\}+(\alpha \gamma \delta )\{A,B\}+\hfill \\ \hfill +& (\gamma ^2ϵ)B^2+(\gamma \delta 2\zeta )B+\hfill \\ \hfill +& \frac{2a}{3}A^3+(d+\frac{a\gamma }{3}+\alpha ^2)A^2+(\frac{aϵ}{3}+\alpha \delta +2z)A\hfill \end{array}$$
(14)
another useful form of the Casimir of the algebra is given by:
$$\begin{array}{cc}\hfill K=& C^2+\frac{2a}{3}A^3\frac{\alpha }{3}\{A,A,B\}\frac{\gamma }{3}\{A,B,B\}+\hfill \\ & +\left(\frac{2\alpha ^2}{3}+d+\frac{2a\gamma }{3}\right)A^2+\left(ϵ+\frac{2\gamma ^2}{3}\right)B^2+\hfill \\ & +\left(\delta +\frac{a\gamma }{3}\right)\{A,B\}+\left(\frac{2\alpha \delta }{3}+\frac{aϵ}{3}+\frac{d\gamma }{3}+2z\right)A+\hfill \\ & +\left(\frac{\alpha ϵ}{3}+\frac{2\delta \gamma }{3}2\zeta \right)B+\frac{\gamma z}{3}\frac{\alpha \zeta }{3}\hfill \end{array}$$
(15)
where
$$\{A,B,C\}=ABC+ACB+BAC+BCA+CAB+CBA$$
This quadratic algebra has many similarities to the Racah algebra $`QR(3)`$, which is a special case of the Askey - Wilson algebra $`QAW(3)`$. The algebra (1113) does not coincide with the Racah algebra $`QR(3)`$, if $`a0`$ in the relation (13). Unless this difference between (10) and $`QR(3)`$ algebra a representation theory can be constructed by following the same procedures as they were described by Granovskii, Lutzenko and Zhedanov in ref. . In this paper we shall give a realization of this algebra using the deformed oscillator techniques. The finite dimensional representations of the algebra (10) will be constructed by constructing a realization of the algebra with the generalized parafermionic algebra introduced by Quesne.
## V Deformed Parafermionic Algebra
Let now consider a realization of the algebra (1113), by using of the deformed oscillator technique, i.e. by using a deformed oscillator algebra $`\{b^{},b,𝒩\}`$, which satisfies the
$$[𝒩,b^{}]=b^{},[𝒩,b]=b,b^{}b=\mathrm{\Phi }\left(𝒩\right),bb^{}=\mathrm{\Phi }\left(𝒩+1\right)$$
(16)
where the function $`\mathrm{\Phi }(x)`$ is a ”well behaved” real function which satisfies the boundary condition:
$$\mathrm{\Phi }(0)=0,\text{and}\mathrm{\Phi }(x)>0\text{for}x>0$$
(17)
As it is well known this constraint imposes the existence a Fock type representation of the deformed oscillator algebra, which is bounded by bellow, i.e. there is a Fock basis $`|n>,n=0,1,\mathrm{}`$ such that
$$\begin{array}{c}𝒩|n>=n|n>\hfill \\ b^{}|n>=\sqrt{\mathrm{\Phi }\left(n+1\right)}|n+1>,n=0,1,\mathrm{}\hfill \\ b|0>=0\hfill \\ b|n>=\sqrt{\mathrm{\Phi }\left(n\right)}|n1>,n=1,2,\mathrm{}\hfill \end{array}$$
(18)
The Fock representation (18) is bounded by bellow. The generalized deformed algebra given in ref is equivalent to several deformed oscillator schemes as the Odaka- Kishi - Kamefuchi unification scheme , the Beckers- Debergh unification scheme , The Fibonacci oscillator , for a discussion of deformation schemes see
In the case of nilpotent deformed oscillator algebras, there is a positive integer $`p`$, such that
$$b^{p+1}=0,\left(b^{}\right)^{p+1}=0$$
the above equations imply that
$$\mathrm{\Phi }(p+1)=0,$$
(19)
In that case the deformed oscillator (16) has a finite dimensional representation, with dimension equal to $`p+1`$, this kind of oscillators are called deformed parafermion oscillators of order $`p`$.
An interesting property of the deformed parafermionic algebra is that the existence of a faithful finite dimensional representation of the algebra implies that:
$$𝒩\left(𝒩1\right)\left(𝒩2\right)\mathrm{}\left(𝒩p\right)=0$$
(20)
The above restriction and the constraints (17) and (19) imply that the general form of the structure function $`\mathrm{\Phi }(𝒩)`$ has the general form:
$$\mathrm{\Phi }(𝒩)=𝒩(p+1𝒩)(a_0+a_1𝒩+a_2𝒩^2+\mathrm{}a_{p1}𝒩^{p1})$$
A systematic study and applications of the parafermionic oscillator is given in references .
We shall show, that there is a realization of the quadratic algebra , such that
$`A=A\left(𝒩\right)`$ (21)
$`B=b\left(𝒩\right)+b^{}\rho \left(𝒩\right)+\rho \left(𝒩\right)b`$ (22)
where the $`A[x],b[x]`$ and $`\rho (x)`$ are functions, which will be determined. In that case (11) implies:
$$C=[A,B]C=b^{}\mathrm{\Delta }A\left(𝒩\right)\rho \left(𝒩\right)\rho \left(𝒩\right)\mathrm{\Delta }A\left(𝒩\right)b$$
(23)
where
$$\mathrm{\Delta }A\left(𝒩\right)=A\left(𝒩+1\right)A\left(𝒩\right)$$
Using equations (21), (22) and (12) we find:
$$\begin{array}{cc}\hfill [A,C]=& [A\left(𝒩\right),b^{}\mathrm{\Delta }A\left(𝒩\right)\rho \left(𝒩\right)\rho \left(𝒩\right)\mathrm{\Delta }A\left(𝒩\right)b]=\hfill \\ \hfill =& b^{}\left(\mathrm{\Delta }A\left(𝒩\right)\right)^2\rho \left(𝒩\right)+\rho \left(𝒩\right)\left(\mathrm{\Delta }A\left(𝒩\right)\right)^2b=\hfill \\ \hfill =& \alpha A^2+\gamma \{A,B\}+\delta A+ϵB+\zeta =\hfill \\ \hfill =& b^{}\left(\gamma \left(A\left(𝒩+1\right)+A\left(𝒩\right)\right)+ϵ\right)\rho \left(𝒩\right)+\hfill \\ & +\rho \left(𝒩\right)\left(\gamma \left(A\left(𝒩+1\right)+A\left(𝒩\right)\right)+ϵ\right)b+\hfill \\ & +\alpha A\left(𝒩\right)^2+2\gamma A\left(𝒩\right)b\left(𝒩\right)+\delta A\left(𝒩\right)+ϵB\left(𝒩\right)+\zeta \hfill \end{array}$$
(24)
therefore we have the following relations:
$`\left(\mathrm{\Delta }A\left(𝒩\right)\right)^2=\gamma \left(A\left(𝒩+1\right)+A\left(𝒩\right)\right)+ϵ`$ (25)
$`\alpha A\left(𝒩\right)^2+2\gamma A\left(𝒩\right)b\left(𝒩\right)+\delta A\left(𝒩\right)+ϵB\left(𝒩\right)+\zeta =0`$ (26)
while the function $`\rho \left(𝒩\right)`$ can be arbitrarily determined. In fact this function can be fixed, in order to have a polynomial structure function $`\mathrm{\Phi }(x)`$ for the deformed oscillator algebra (16).
The solutions of equation (25) depend on the value of the parameter $`\gamma `$, while the function $`b(𝒩)`$ is uniquely determined by equation (26) (provided that almost one among the parameters $`\gamma `$ or $`ϵ`$ is not zero). At this stage, the cases $`\gamma 0`$ or $`\gamma =0`$, should be treated separately. We can see that:
* $`\gamma 0`$
In that case the solutions of equations (25) and (26) are given by:
$$A\left(𝒩\right)=\frac{\gamma }{2}\left((𝒩+u)^21/4\frac{ϵ}{\gamma ^2}\right)$$
(27)
$$\begin{array}{cc}\hfill b\left(𝒩\right)=& \frac{\alpha \left((𝒩+u)^21/4\right)}{4}+\frac{\alpha ϵ\delta \gamma }{2\gamma ^2}\hfill \\ & \frac{\alpha ϵ^22\delta ϵ\gamma +4\gamma ^2\zeta }{4\gamma ^4}\frac{1}{\left((𝒩+u)^21/4\right)}\hfill \end{array}$$
(28)
* $`\gamma =0,ϵ0`$
The solutions of equations (25) and (26) are given by:
$$A(𝒩)=\sqrt{ϵ}\left(𝒩+u\right)$$
(29)
$$b(𝒩)=\alpha \left(𝒩+u\right)^2\frac{\delta }{\sqrt{ϵ}}\left(𝒩+u\right)\frac{\zeta }{ϵ}$$
(30)
The constant $`u`$ will be determined later.
Using the above definitions of equations $`A(𝒩)`$ and $`b(𝒩)`$, the left hand side and right hand side of equation (13) gives the following equation:
$$\begin{array}{c}2\mathrm{\Phi }(𝒩+1)\left(\mathrm{\Delta }A\left(𝒩\right)+\frac{\gamma }{2}\right)\rho (𝒩)2\mathrm{\Phi }(𝒩)\left(\mathrm{\Delta }A\left(𝒩1\right)\frac{\gamma }{2}\right)\rho (𝒩1)=\hfill \\ =aA^2\left(𝒩\right)\gamma b^2(𝒩)2\alpha A\left(𝒩\right)b(𝒩)+dA\left(𝒩\right)\delta b(𝒩)+z\hfill \end{array}$$
(31)
Equation (14) gives the following relation:
$$\begin{array}{cc}\hfill K=& \\ \hfill =& \mathrm{\Phi }(𝒩+1)\left(\gamma ^2ϵ2\gamma A\left(𝒩\right)\mathrm{\Delta }A^2\left(𝒩\right)\right)\rho (𝒩)+\hfill \\ & +\mathrm{\Phi }(𝒩)\left(\gamma ^2ϵ2\gamma A\left(𝒩\right)\mathrm{\Delta }A^2\left(𝒩1\right)\right)\rho (𝒩1)\hfill \\ & 2\alpha A^2\left(𝒩\right)b(𝒩)+\left(\gamma ^2ϵ2\gamma A\left(𝒩\right)\right)b^2(𝒩)+\hfill \\ & +2\left(\alpha \gamma \delta \right)A\left(𝒩\right)b(𝒩)+\left(\gamma \delta 2\zeta \right)b(𝒩)+\hfill \\ & +\frac{2}{3}aA^3\left(𝒩\right)+\left(d+\frac{1}{3}a\gamma +\alpha ^2\right)A^2\left(𝒩\right)+\hfill \\ & +\left(\frac{1}{3}aϵ+\alpha \delta +2z\right)A\left(𝒩\right)\hfill \end{array}$$
(32)
Equations (31) and (32) are linear functions of the expressions $`\mathrm{\Phi }\left(𝒩\right)`$ and $`\mathrm{\Phi }\left(𝒩+1\right)`$, then the function $`\mathrm{\Phi }\left(𝒩\right)`$ can be determined, if the function $`\rho (𝒩)`$ is given. The solution of this system, i.e. the function $`\mathrm{\Phi }\left(𝒩\right)`$ depends on two parameters $`u`$ and $`K`$ and it is given by the following formulae:
* $`\gamma 0`$
$$\rho (𝒩)=\frac{1}{32^{12}\gamma ^8(𝒩+u)(1+𝒩+u)(1+2(𝒩+u))^2}$$
and
$$\begin{array}{c}\mathrm{\Phi }(𝒩)=3072\gamma ^6K(1+2(𝒩+u))^2\hfill \\ 48\gamma ^6(\alpha ^2ϵ\alpha \delta \gamma +aϵ\gamma d\gamma ^2)\hfill \\ (3+2(𝒩+u))(1+2(𝒩+u))^4(1+2(𝒩+u))+\hfill \\ +\gamma ^8(3\alpha ^2+4a\gamma )(3+2(𝒩+u))^2(1+2(𝒩+u))^4(1+2(𝒩+u))^2+\hfill \\ +768(\alpha ϵ^22\delta ϵ\gamma +4\gamma ^2\zeta )^2+\hfill \\ +32\gamma ^4(1+2(𝒩+u))^2(112(𝒩+u)+12(𝒩+u)^2)\hfill \\ (3\alpha ^2ϵ^26\alpha \delta ϵ\gamma +2aϵ^2\gamma +2\delta ^2\gamma ^24dϵ\gamma ^2+8\gamma ^3z+4\alpha \gamma ^2\zeta )\hfill \\ 256\gamma ^2(1+2(𝒩+u))^2\hfill \\ (3\alpha ^2ϵ^39\alpha \delta ϵ^2\gamma +aϵ^3\gamma +6\delta ^2ϵ\gamma ^23dϵ^2\gamma ^2+2\delta ^2\gamma ^4+\hfill \\ +2dϵ\gamma ^4+12ϵ\gamma ^3z4\gamma ^5z+12\alpha ϵ\gamma ^2\zeta 12\delta \gamma ^3\zeta +4\alpha \gamma ^4\zeta )\hfill \end{array}$$
(33)
* $`\gamma =0,ϵ0`$
$$\rho (𝒩)=1$$
$$\begin{array}{c}\mathrm{\Phi }(𝒩)=\hfill \\ =\frac{1}{4}\left(\frac{K}{ϵ}\frac{z}{\sqrt{ϵ}}\frac{\delta }{\sqrt{ϵ}}\frac{\zeta }{ϵ}+\frac{\zeta ^2}{ϵ^2}\right)\hfill \\ \frac{1}{12}\left(3da\sqrt{ϵ}3\alpha \frac{\delta }{\sqrt{ϵ}}+3\left(\frac{\delta }{\sqrt{ϵ}}\right)^26\frac{z}{\sqrt{ϵ}}+6\alpha \frac{\zeta }{ϵ}6\frac{\delta }{\sqrt{ϵ}}\frac{\zeta }{ϵ}\right)(𝒩+u)\hfill \\ +\frac{1}{4}\left(\alpha ^2+da\sqrt{ϵ}3\alpha \frac{\delta }{\sqrt{ϵ}}+\left(\frac{\delta }{\sqrt{ϵ}}\right)^2+2\alpha \frac{\zeta }{ϵ}\right)(𝒩+u)^2\hfill \\ \frac{1}{6}\left(3\alpha ^2a\sqrt{ϵ}3\alpha \frac{\delta }{\sqrt{ϵ}}\right)(𝒩+u)^3+\frac{1}{4}\alpha ^2(𝒩+u)^4\hfill \end{array}$$
(34)
The above formula is valid for $`ϵ>0`$.
## VI Finite dimensional representations of quadratic algebras
Let consider a representation of the quadratic algebra , which is diagonal to the generator $`A`$ and the Casimir $`K`$. Using the parafermionic realization defined by equations (21) and (22), we see that this a representation diagonal to the parafermionic number operator $`𝒩`$ and the Casimir $`K`$. The basis of a such representation corresponds to the Fock basis of the parafermionic oscillator, i.e. the vectors $`|k,n>,n=0,1,\mathrm{}`$of the carrier Fock space satisfy the equations
$$𝒩|k,n>=n|k,n>,K|k,n>=k|k,n>$$
The structure function (33) (or respectively (33) ) depend on the eigenvalues of the of the parafermionic number operator $`𝒩`$ and the Casimir $`K`$. The vectors $`|k,n>`$ are also eigenvectors of the generator $`A`$, i.e.
$$A|k,n>=A(k,n)|k,n>$$
In the case $`\gamma 0`$ we find from equation (27)
$$A(k,n)=\frac{\gamma }{2}\left((n+u)^21/4\frac{ϵ}{\gamma ^2}\right)$$
In the case $`\gamma =0,ϵ0`$ we find from equation (29)
$$A(k,n)=\sqrt{ϵ}\left(n+u\right).$$
If the deformed oscillator corresponds to a deformed Parafermionic oscillator of order $`p`$ then the two parameters of the calculation $`k`$ and $`u`$ should satisfy the constrints (17) and (19) the system:
$$\begin{array}{c}\mathrm{\Phi }(0,u,k)=0\\ \mathrm{\Phi }(p+1,u,k)=0\end{array}$$
(35)
then the parameter $`u=u(k,p)`$ is a solution of the system of equations (35).
Generally there are many solutions of the above system, but a unitary representation of the deformed parafermionic oscillator is restrained by the additional restriction
$$\mathrm{\Phi }(x)>0,\text{for}x=1,2,\mathrm{},p$$
We must point out that the system (35) corresponds to a representation with dimension equal to $`p+1`$.
The proposed method of calculation of the representation of the quadratic algebra is an alternative to the method given by Granovskii et al and reduces the search of the representations to the solution of a system of polynomial equations (35). Also its is applied to an algebra not included in the cases of the algebras, which are treated in the above references. We must point out, that there are several papers on the representations of quadratic ( or generally polynomial algebras) , these algebras are deformations of the su(2) or osp(1/2) algebras. The general form of the quadratic algebra, which is studied in this paper, is different by definition from the deformed versions of su(2) or osp(1/2).
## VII Quadratic algebras for the quantum superintegrable systems
In this section, we shall give an example of the calculation of eigenvalues of a superintegrable two-dimensional system, by using the methods of the previous section. The calculation by an empirical method was performed in and the solution of the same problem by using separation of variables was studied in . Here in order to show the effects of the quantization procedure we don’t use $`\mathrm{}=1`$ as it was considered in references and . That means that the following commutation relations are taken in consideration:
$$[x,p_x]=i\mathrm{},[y,p_y]=i\mathrm{}$$
### VII-a Potential i)
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\omega ^2r^2+\frac{\mu _1}{x^2}+\frac{\mu _2}{y^2}\right)$$
This potential has the following independent integrals of motion:
$$A=p_x^2+\omega ^2x^2+\frac{\mu _1}{x^2}$$
and
$$B=\left(xp_yyp_x\right)^2+r^2\left(\frac{\mu _1}{x^2}+\frac{\mu _2}{y^2}\right)$$
The constants, which characterize the corresponding quadratic algebra (10), are given by:
$$\begin{array}{ccc}\alpha =8\mathrm{}^2,\hfill & \gamma =0,\hfill & \delta =16h^2H,\hfill \\ ϵ=16\mathrm{}^2\omega ^2,\hfill & \zeta =16\mathrm{}^2(\mu _1+\mu _2)\omega ^2+8\mathrm{}^4\omega ^2\hfill & \\ a=0,\hfill & d=16\mathrm{}^4,\hfill & z=16\mathrm{}^2(\mu _2\mu _1)\omega ^216\mathrm{}^4H\hfill \end{array}$$
the value of the Casimir (14) is:
$$K=16\mathrm{}^2\left((\mu _2\mu _1)^2\omega ^2+4\mu _1H^2\right)16\mathrm{}^4\left(3H^2+2\mathrm{}^2\omega ^22(\mu _1+\mu _2)\right)$$
For simplicity reasons we introduce the positive parameters $`k_1`$ and $`k_1`$, which are related to the potential parameters $`\mu _1`$ and $`\mu _2`$ by the relations:
$$\mu _1=\left(k_1^2\frac{1}{4}\right)\mathrm{}^2\mu _2=\left(k_2^2\frac{1}{4}\right)\mathrm{}^2$$
This quadratic algebra corresponds to the case $`\gamma =0`$ and $`ϵ>0`$ of the algebra given by equations (1113). In that case, the structure function (34) of the deformed parafermionic algebra of Section V can be given by the simple form:
$$\begin{array}{cc}\hfill \mathrm{\Phi }(x)=& 16\mathrm{}^4(x+u\frac{1}{2}\frac{k_1}{2})(x+u\frac{1}{2}+\frac{k_1}{2})\hfill \\ & \left(x+u\frac{1}{2}\frac{k_2}{2}\frac{E}{2\mathrm{}\omega }\right)\left(x+u\frac{1}{2}+\frac{k_2}{2}\frac{E}{2\mathrm{}\omega }\right)\hfill \end{array}$$
In the above formula $`E`$ is the eigenvalue of the energy. The values of the parameters $`u`$ and $`E`$ corresponding to the a representation of the parafermionic algebra of dimension $`p+1`$ are determined by the restrictions (35), which are transcribed as:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(p+1)=0$$
One should notice, that only the solutions which correspond to positive eigenvalues of the integral $`A`$ must be retained. The acceptable solutions are four and correspond to the following values of the parameters $`u`$ and $`E`$:
$$u=\frac{1}{2}+\frac{ϵ_1k_1}{2},E=2\mathrm{}\omega \left(p+1+\frac{ϵ_1k_1+ϵ_2k_2}{2}\right)$$
where $`ϵ_i=\pm 1`$. The corresponding structure function is
$$\mathrm{\Phi }(x)=16\mathrm{}^4x\left(p+1x\right)\left(x+ϵ_1k_1\right)\left(p+1x+ϵ_2k_2\right)$$
The corresponding eigenvalues of the operator $`A`$ are given by:
$$A(m)=4\mathrm{}\omega \left(m+\frac{ϵ_1k_1+ϵ_2k_2}{2}\right),m=0,1,\mathrm{},p$$
The structure function $`\mathrm{\Phi }(x)`$ should be a positive function, for $`x=1,2,\mathrm{},p`$ therefore the constants $`k_i`$ are restricted by the relations:
$$ϵ_1k_1>1,ϵ_2k_2>1$$
### VII-b Potential ii)
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\omega ^2\left(4x^2+y^2\right)+\frac{\mu }{y^2}\right)$$
This potential has the following independent integrals of motion:
$$A=p_x^2+4\omega ^2x^2$$
and
$$B=\frac{1}{2}\{xp_yyp_x,p_y\}+\frac{\mu x}{y^2}\omega ^2xy^2$$
The constants, which characterize the corresponding quadratic algebra (10), are given by:
$$\begin{array}{ccc}\alpha =0,\hfill & \gamma =0,\hfill & \delta =0,\hfill \\ ϵ=16\mathrm{}^2\omega ^2,\hfill & \zeta =0,\hfill & \\ a=6\mathrm{}^2,\hfill & d=16\mathrm{}^2H,\hfill & z=8\mathrm{}^2(\mu \omega ^2H^2)+6\mathrm{}^4\omega ^2\hfill \end{array}$$
the value of the Casimir (14) is:
$$K=64\mathrm{}^4\omega ^2H$$
For simplicity reasons we introduce the positive parameters $`k`$, which is related to the potential parameter $`\mu `$ by the relation:
$$\mu =\left(k^2\frac{1}{4}\right)\mathrm{}^2$$
This quadratic algebra corresponds to the case $`\gamma =0`$ and $`ϵ>0`$ of the algebra given by equations (1113). In that case, the structure function (34) of the deformed parafermionic algebra of Section V can be given by the simple form:
$$\mathrm{\Phi }(x)=8\mathrm{}^3\omega \left(x+u\frac{1}{2}\right)\left(x+u\frac{1}{2}\frac{k}{2}\frac{E}{2\mathrm{}\omega }\right)\left(x+u\frac{1}{2}+\frac{k}{2}\frac{E}{2\mathrm{}\omega }\right)$$
In the above formula $`E`$ is the eigenvalue of the energy. The values of the parameters $`u`$ and $`E`$ corresponding to the a representation of the parafermionic algebra of dimension $`p+1`$ are determined by the restrictions (35), which are transcribed as:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(p+1)=0$$
One should notice, that only the solutions which correspond to positive eigenvalues of the integral $`A`$ must be retained. The acceptable solutions are four and correspond to the following values of the parameters $`u`$ and $`E`$:
$$u=\frac{1}{2},E=2\mathrm{}\omega \left(p+1+\frac{ϵk}{2\mathrm{}}\right)$$
where $`ϵ=\pm 1`$. The corresponding structure function is
$$\mathrm{\Phi }(x)=4\mathrm{}^3x\left(p+1x\right)\left(p+1x+ϵk\right)$$
The structure function should be a positive function, therefore the values of the parameter $`k`$ are restrained by
$$ϵk>1$$
The eigenvalues of the operator $`A`$ are given by:
$$A(m)=4\mathrm{}\omega (m+\frac{1}{2}),m=0,1,\mathrm{},p$$
### VII-c Potential iii)
$$H=\frac{1}{2}\left(p_x^2+p_y^2+\frac{k}{r}+\frac{1}{r}\left(\frac{\mu _1}{r+x}+\frac{\mu _2}{rx}\right)\right)$$
In ref the parabolic coordinates have been used:
$$\begin{array}{cc}x=\frac{1}{2}\left(\xi ^2\eta ^2\right),\hfill & p_x=\frac{\xi }{\xi ^2+\eta ^2}p_\xi \frac{\eta }{\xi ^2+\eta ^2}p_\eta ,\hfill \\ y=\xi \eta ,\hfill & p_y=\frac{\eta }{\xi ^2+\eta ^2}p_\xi +\frac{\xi }{\xi ^2+\eta ^2}p_\eta ,\hfill \\ [\xi ,p_\xi ]=i\mathrm{},\hfill & [\eta ,p_\eta ]=i\mathrm{}\hfill \end{array}$$
For comparison reasons we quote all the relations in both, cartesian and parabolic systems, so
$$H=\frac{1}{\xi ^2+\eta ^2}\left(\frac{1}{2}\left(p_\xi ^2+p_\eta ^2\right)+k+\frac{\mu _1}{\xi ^2}+\frac{\mu _2}{\eta ^2}\right)$$
This potential has the following independent integrals of motion:
$$\begin{array}{cc}\hfill A=& (xp_yyp_x)^2+r\left(\frac{\mu _1}{r+x}+\frac{\mu _1}{rx}\right)=\hfill \\ \hfill =& \frac{1}{2}\left(\frac{1}{2}\left(\eta p_\xi \xi p_\eta \right)^2+\left(\xi ^2+\eta ^2\right)\left(\frac{\mu _1}{\xi ^2}+\frac{\mu _2}{\eta ^2}\right)\right)\hfill \end{array}$$
and
$$\begin{array}{cc}\hfill B=& \frac{1}{2}\left(\{xp_yyp_x,p_y\}\frac{\mu _1}{r}\frac{rx}{r+x}+\frac{\mu _2}{r}\frac{r+x}{rx}+\frac{kx}{r}\right)\hfill \\ \hfill =& \frac{1}{\xi ^2+\eta ^2}\left(\frac{1}{2}\left(\xi ^2p_\eta ^2\eta ^2p_\xi ^2\right)+\mu _2\frac{\xi ^2}{\eta ^2}\mu _1\frac{\eta ^2}{\xi ^2}+\frac{k}{2}\frac{\xi ^2\eta ^2}{\xi ^2+\eta ^2}\right)\hfill \end{array}$$
The constants, which characterize the corresponding quadratic algebra (10), are given by:
$$\begin{array}{ccc}\alpha =0,\hfill & \gamma =2\mathrm{}^2,\hfill & \delta =0,\hfill \\ ϵ=\mathrm{}^4,\hfill & \zeta =\mathrm{}^2k(\mu _1\mu _2),\hfill & \\ a=0,\hfill & d=8\mathrm{}^2H,\hfill & z=\mathrm{}^2\left(4(\mu _1+\mu _2)Hk^2/2\right)+\mathrm{}^4H\hfill \end{array}$$
the value of the Casimir (14) is:
$$K=\mathrm{}^2\left(2(\mu _1\mu _2)^2Hk^2(\mu _1+\mu _2)\right)2\mathrm{}^4\left((\mu _1+\mu _2)H\frac{k^2}{4}\right)+\mathrm{}^6H$$
For simplicity reasons we introduce the positive parameters $`k_1`$ and $`k_1`$, which are related to the potential parameters $`\mu _1`$ and $`\mu _2`$ by the relations:
$$\mu _1=\frac{\mathrm{}^2}{2}\left(k_1^2\frac{1}{4}\right)\mu _2=\frac{\mathrm{}^2}{2}\left(k_2^2\frac{1}{4}\right)$$
This quadratic algebra corresponds to the case $`\gamma 0`$ of the algebra given by equations (1113). In that case, the structure function (33) of the deformed parafermionic algebra of Section V can be given by the simple form:
$$\begin{array}{ccc}\hfill \mathrm{\Phi }(x)=& \hfill 32^{14}\mathrm{}^{16}& (2x1+k_1+k_2)(2x1+k_1k_2)(2x1k_1+k_2)\hfill \\ & \hfill & \left(2x1k_1k_2\right)\left(8\mathrm{}^2Hx^28\mathrm{}^2Hx+2\mathrm{}^2H+k^2\right)\hfill \end{array}$$
In the above formula $`E`$ is the eigenvalue of the energy. The values of the parameters $`u`$ and $`E`$ corresponding to the a representation of the parafermionic algebra of dimension $`p+1`$ are determined by the restrictions (35), which are transcribed as:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(p+1)=0$$
One should notice, that only the solutions which correspond to positive eigenvalues of the integral $`A`$ must be retained. The acceptable solutions are four and correspond to the following values of the parameters $`u`$ and $`E`$:
$$u=\frac{1}{2}\left(2+ϵ_1k_1+ϵ_2k_2\right),E=\frac{k^2}{2\mathrm{}^2\left(2(p+1)+ϵ_1k_1+ϵ_2k_2\right)^2}$$
where $`ϵ_i=\pm 1`$. The corresponding structure function is
$$\begin{array}{ccc}\hfill \mathrm{\Phi }(x)=& \hfill 32^{20}k^2\mathrm{}^{16}& x(p+1x)(x+ϵ_1k_1)(x+ϵ_2k_2)\hfill \\ & \hfill & \left(x+ϵ_1k_1+ϵ_2k_2\right)\frac{\left(x+p+1+ϵ_1k_1+ϵ_2k_2\right)}{\left(2(p+1)+ϵ_1k_1+ϵ_2k_2\right)^2}\hfill \end{array}$$
The eigenvalues of the operator $`A`$ are given by the formula:
$$A(m)=\mathrm{}^2\left(m+ϵ_1k_1+ϵ_2k_2+\frac{3}{2}\right)^2,m=0,1,\mathrm{},p$$
The positive sign of the structure function for $`x=1,2,\mathrm{},p`$ is obtained when:
$$ϵ_1k_1>1,ϵ_2k_2>1,\text{and}ϵ_1k_1+ϵ_2k_2>1$$
### VII-d Potential iv)
$$\begin{array}{cc}\hfill H=& \frac{1}{2}(p_x^2+p_y^2++\frac{k}{r}+\mu _1\frac{\sqrt{r+x}}{r}+\mu _2\frac{\sqrt{rx}}{r})=\hfill \\ \hfill =& \frac{1}{\xi ^2+\eta ^2}\left(\frac{1}{2}\left(p_\xi ^2+p_\eta ^2\right)+k+\mu _1\xi +\mu _2\eta \right)\hfill \end{array}$$
This potential has the following independent integrals of motion:
$$\begin{array}{cc}\hfill A=& \frac{1}{2}\left(\{(yp_xxp_y),p_y\}+\frac{\mu _1(rx)\sqrt{r+u}}{r}\frac{\mu _2(r+x)\sqrt{ru}}{r}\frac{kx}{r}\right)=\hfill \\ \hfill =& \frac{1}{2\left(\xi ^2+\eta ^2\right)}\left(\eta ^2p_\xi ^2\xi ^2p_\eta ^2+k\left(\eta ^2\xi ^2\right)+2\xi \eta \left(\mu _1\eta \mu _2\xi \right)\right)\hfill \end{array}$$
and
$$\begin{array}{cc}\hfill B=& \frac{1}{2}\left(\{xp_yyp_x,p_x\}\frac{\mu _1x\sqrt{ru}}{r}+\frac{\mu _2x\sqrt{r+u}}{r}\frac{ky}{r}\right)=\hfill \\ \hfill =& \frac{1}{2\left(\xi ^2+\eta ^2\right)}\left(\xi \eta \left(p_\xi ^2+p_\eta ^2\right)\left(\xi ^2+\eta ^2\right)p_\xi p_\eta +2k\xi \eta +\left(\mu _2\xi \mu _1\eta \right)\left(\eta ^2\xi ^2\right)\right)\hfill \end{array}$$
The constants, which characterize the corresponding quadratic algebra (10), are given by:
$$\begin{array}{ccc}\alpha =0,\hfill & \gamma =0,\hfill & \delta =0,\hfill \\ ϵ=2\mathrm{}^2H,\hfill & \zeta =\mathrm{}^2\mu _1\mu _2/2,\hfill & \\ a=0,\hfill & d=2\mathrm{}^2H,\hfill & z=\mathrm{}^2(\mu _1^2\mu _2^2)/4\hfill \end{array}$$
the value of the Casimir (14) is:
$$K=\mathrm{}^2k^2H/2+\mathrm{}^2k(\mu _1^2+\mu _2^2)/4+\mathrm{}^4H^2$$
This quadratic algebra corresponds to the case $`\gamma =0`$ and $`ϵ>0`$ of the algebra given by equations (1113). It is worth noticing that the algebra is extremely simple, which can be reduced to the usual $`su(2)`$. We prefer to treat this algebra with the proposed methods in this paper for pedagogical reasons. The existence of the finite dimensional representations of this algebra implies that, the coefficient $`ϵ`$ in equation (12) should be positive, therefore the energy operator $`H`$ must have energy eigenvalues $`E<0`$. For simplicity reasons we introduce the new parameters:
$$\begin{array}{ccc}\epsilon =\sqrt{2E}/\mathrm{},\hfill & \lambda =k/\mathrm{}^2,\hfill & \\ \nu _1=\mu _1/\mathrm{}^2,\hfill & \nu _2=\mu _2/\mathrm{}^2,\hfill & \nu ^2=\nu _1^2+\nu _2^2\hfill \end{array}$$
The structure function (34) of the deformed parafermionic algebra of Section V can be given by the form:
$$\mathrm{\Phi }(x)=\frac{\mathrm{}^4}{16\epsilon ^4}\left(\nu _1^2\lambda \epsilon ^2+2(x+u\frac{1}{2})\epsilon ^3\right)\left(\nu _2^2\lambda \epsilon ^22(x+u\frac{1}{2})\epsilon ^3\right)$$
In the above formula the parameter $`\epsilon `$ is related to the the eigenvalue $`E`$ of the energy. The values of the parameters $`u`$ and $`\epsilon `$, corresponding to the a representation of the parafermionic algebra of dimension $`p+1`$, are determined by the restrictions (35), which are transcribed as:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(p+1)=0$$
The first condition can be used for determining the acceptable values of the parameter $`u`$. Two possible solutions are found:
$`u=u_1={\displaystyle \frac{\nu _2^2\lambda \epsilon ^2+\epsilon ^3}{2\epsilon ^3}}`$ (36)
$`u=u_2={\displaystyle \frac{\nu _1^2\lambda \epsilon ^2\epsilon ^3}{2\epsilon ^3}}`$ (37)
Using these solutions and the condition $`Phi(p+1)=0`$, we find that the $`\epsilon `$ must satisfy two possible cubic equations:
$`u_1\mathrm{\hspace{0.33em}2}(p+1)\epsilon ^32\lambda \epsilon ^2+\nu ^2=0`$ (38)
$`u_2\mathrm{\hspace{0.33em}2}(p+1)\epsilon ^3+2\lambda \epsilon ^2\nu ^2=0`$ (39)
If $`\epsilon `$ is a solution of equation (38) then $`\epsilon `$ is the solution of the other equation (39), therefore there is almost one solution which is positive. This solution leads to the structure function:
$$\mathrm{\Phi }(x)=\frac{\epsilon ^2}{4}x\left(p+1x\right)$$
which is positive for $`x=1,2,\mathrm{},p`$.
## VIII Discussion
If we compare the quadratic associative algebra, introduced in section IV with the corresponding Poisson algebra given in section II, we see that in general the quantum constants are similar to the classical ones up to an factor equal to $`h^2`$, but there are quantum corrections of order $`h^4`$ and $`h^6`$. The knowledge of the classical constants of the Poisson algebra is not sufficient to reproduce the rules of quantum associative operator algebra. Therefore, the passage from the classic Poisson algebra to the non commutative quantum algebra can not be realized by simple replacements of the Poisson brackets by commutators and by a symmetrization procedure.
The energy eigenvalues of section VII corroborate the results of reference (the differences in the case of the potential iv are due to some misprints in that reference). The calculation of the energy eigenvalues in reference was achieved by solving the corresponding Schroedinger differential equations, while in this paper the energy eigenvalues are obtained by algebraic methods. The advantage of the proposed method is that, the energy eigenvalues are reduced to simple algebraic calculations of the roots of polynomial equations, whose the form is universally determined by the the structure functions (33), (33) and the system (35). These equations are valid for any two dimensional superintegrable system with integrals of motion, which are quadratic functions of the momenta. The same equations should be valid in the case of two dimensional superintegrable systems in curved space . The superintegrable systems bring up for discussion the open problem of the quantization of a Poisson algebra in a well determined context, because these systems and their quantum counterparts are explicitly known.
From the above discussion several open problems are risen:
* The calculation of the classical Poisson algebras and their quantum counterparts for the totality of the two dimensional problems in curved space. This study will lead to the calculation of the energy eigenvalues by algebraic methods.
* The two dimensional superintegrable systems are classified by the values of the constants of the Darboux conditions . The relation of these constants with the constants of the quadratic Poisson algebra is not yet known.
* The Poisson algebras for the Drach superintegrable systems with a cubic integral of motion were written by using a classical analogue of the deformed parafermionic algebra . Their quantum counterparts and the calculation of their energy eigenvalues are is a topic under investigation.
* The Poisson algebras and the associated quantum counterparts for the three dimensional superintegrable systems are not yet fully studied. Recently the quantum quadratic algebras have been written down, but a systematic calculation of energy eigenvalues was not yet performed.
The above points show that, the study of non linear Poisson algebras and their quantum counterparts is a topic of interest. |
warning/0003/hep-th0003164.html | ar5iv | text | # 1 Introduction
## 1 Introduction
There are great progresses in studying worldsheet instanton corrections by the discovery of the mirror symmetry ,. When we study topological sigma model , physical observables are constructed by combining couplings, correlation functions and metrics of operators. Typical constituent blocks are three-point and two-point functions (metrics). Their behaviors are controlled by moduli parameters and these observables contain information about moduli space of the sigma model -. In this paper, we investigate a Kähler potential of the B-model moduli space from the point of view of the topological sigma model together with results by CFT at Gepner point.
## 2 Calabi-Yau $`d`$-fold
We concentrate on a one-parameter family of Calabi-Yau $`d`$-fold $`M`$ realized as a zero locus of a hypersurface in a projective space $`CP^{d+1}`$
$`W;\widehat{\{p=0\}/𝐙_N^{(N1)}},`$
$`p=X_1^N+X_2^N+\mathrm{}+X_N^NN\psi X_1X_2\mathrm{}X_N=0,`$
where we introduce a number $`N=d+2`$. Hodge numbers $`h^{p,q}`$ of this $`d`$ fold are calculated as
$`h^{p,q}=\delta _{p+q.d},(0pd,0qd,pq),`$
$`h^{p,p}=\delta _{2p,d}+{\displaystyle \underset{m=0}{\overset{p}{}}}(1)^m\left(\begin{array}{c}d+2\\ m\end{array}\right)\left(\begin{array}{c}(p+1m)(d+1)+p\\ d+1\end{array}\right),`$
and deformation of complex structure is parametrized by a complex parameter $`\psi `$.
In order to describe the complex structure moduli space near Gepner point ($`\psi 0`$), we uses a basis of periods
$`\stackrel{~}{\varpi }_k(\psi )=\left[\mathrm{\Gamma }\left({\displaystyle \frac{k}{N}}\right)\right]^N{\displaystyle \frac{(N\psi )^k}{\mathrm{\Gamma }(k)}}G_k(\psi ),`$
$`G_k:=1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{k}{N}+n\right)}{\mathrm{\Gamma }\left(\frac{k}{N}\right)}}\right]^N{\displaystyle \frac{\mathrm{\Gamma }(k)}{\mathrm{\Gamma }(Nn+k)}}(N\psi )^{Nn}.`$
They behave around $`\psi 0`$
$`\stackrel{~}{\varpi }_k(\psi )=\left[\mathrm{\Gamma }\left({\displaystyle \frac{k}{N}}\right)\right]^N{\displaystyle \frac{(N\psi )^k}{\mathrm{\Gamma }(k)}}\left[1+𝒪(\psi ^N)\right].`$
This set is a projective coordinate of the B-model moduli space. Also we use a canonical basis $`\{\varpi _k^{(0)}\}`$ of the coordinates at the Gepner point as
$`\stackrel{~}{\varpi }_k=\varpi _k^{(0)}\left[\mathrm{\Gamma }\left({\displaystyle \frac{k}{N}}\right)\right]^N,`$
and take ratios to construct normalized coordinates
$`\omega _{k1}:=\varpi _k^{(0)}/\varpi _1^{(0)}={\displaystyle \frac{(N\psi )^{k1}}{(k1)!}}{\displaystyle \frac{G_k}{G_1}}(k=1,2,\mathrm{},N1).`$
Some part of the cohomology classes of $`H^d(W)`$ is associated with the complex structure deformation with Hodge numbers $`h^{\mathrm{},d\mathrm{}}=1`$ $`(\mathrm{}=0,1,2,\mathrm{},d)`$. Then associated operators are represented as
$`𝒪^{(k)}:=(X_1X_2\mathrm{}X_N)^k(k=0,1,2,\mathrm{},N2),`$
and each $`𝒪^{(k)}`$ is related to a canonical period $`\omega _k`$. Also a coupling $`\stackrel{~}{\psi }:=N\psi `$ is associated with an operator $`𝒪^{(1)}=X_1X_2\mathrm{}X_N`$. When we choose a canonical set of periods, these operators are normalized appropriately and fusion couplings $`\kappa _{\mathrm{}}`$ are evaluated at $`\psi =0`$
$`\stackrel{}{\omega }:={}_{}{}^{t}(\begin{array}{cccc}\omega _0& \omega _1& \mathrm{}& \omega _{N2}\end{array}),`$
$`𝐊:=\left(\begin{array}{cccccc}0& \kappa _0& & & & \\ & 0& \kappa _1& & & \\ & & 0& \kappa _2& & \\ & & & \mathrm{}& \mathrm{}& \\ & & & & 0& \kappa _{d1}\\ 0& 0& 0& \mathrm{}& 0& 0\end{array}\right)`$
$`_{\stackrel{~}{\psi }}\stackrel{}{\omega }=𝐊\stackrel{}{\omega },`$
$`𝒪^{(1)}𝒪^{(\mathrm{})}=\kappa _{\mathrm{}}𝒪^{(\mathrm{}+1)}(\mathrm{}=0,1,\mathrm{},d1).`$
$`\kappa _0=1,\kappa _1=\omega _1=1+𝒪(\psi ^N),`$
$`\kappa _{\mathrm{}}={\displaystyle \frac{1}{\kappa _\mathrm{}1}}{\displaystyle \frac{1}{\kappa _\mathrm{}2}}\mathrm{}{\displaystyle \frac{1}{\kappa _0}}\omega _{\mathrm{}}=1+𝒪(\psi ^N),(\mathrm{}2).`$
In this canonical set of periods, the three-point couplings are normalized to units up to terms with order $`𝒪(\psi ^N)`$.
## 3 Kähler potential
Next let us consider a Kähler potential $`K`$ of the moduli space. We can construct $`K`$ as a quadratic form of periods
$`e^K=i^{d^2}{\displaystyle _M}\mathrm{\Omega }\overline{\mathrm{\Omega }}={\displaystyle \underset{m,n=1}{\overset{N1}{}}}I_{m,n}\stackrel{~}{\varpi }_m^{}\stackrel{~}{\varpi }_n.`$
The matrix $`I=\{I_{m,n}\}`$ is determined by properties of intersection numbers of homology cycles and does not change under any infinitesimal deformation of continuous moduli parameters. That is, the $`I_{m,n}`$s cannot depend on any moduli parameters and turns out to be constant numbers.
On the other hand, global structures of the moduli space are encoded in monodromies and it is important to see property of $`K`$ under a global monodromy transformation. Our basis $`\{\stackrel{~}{\varpi }_k\}`$ diagonalizes a cyclic $`𝐙_N`$ monodromy $`\psi \alpha \psi `$ ($`\alpha =e^{2\pi i/N}`$) at $`\psi =0`$
$`\stackrel{~}{\varpi }_k(\alpha \psi )=\alpha ^k\stackrel{~}{\varpi }_k(\psi )(k=1,2,\mathrm{},N1).`$
$`\alpha =e^{2\pi i/N}.`$
But this transformation induces a change of the $`I_{m,n}`$
$`I_{m,n}I_{m,n}\alpha ^{m+n}(m,n=1,2,\mathrm{},N1).`$
Because the Kähler potential is physical quantity and should not depend on monodromies, it is invariant under the transformation. Only invariant parts of this transformation are diagonal ones $`I_{m,m}`$ $`(m=1,2,\mathrm{},N1)`$ and we find that the matrix $`I`$ is diagonal one:
$`e^K={\displaystyle \underset{k=1}{\overset{N1}{}}}I_k\stackrel{~}{\varpi }_k^{}\stackrel{~}{\varpi }_k.`$ (1)
## 4 Determination of $`I_k`$
Next, in order to determine the Kähler potential, all we have to do is to fix the diagonal matrix. We use a method of the $`tt^{}`$ fusion in fixing the $`I_k`$ ($`k=1,2,\mathrm{},N1`$). First note that the Kähler potential is related to a moduli space metric $`g_{\stackrel{~}{\psi }\overline{\stackrel{~}{\psi }}}`$
$`\overline{}K=g_{\stackrel{~}{\psi }\overline{\stackrel{~}{\psi }}},={\displaystyle \frac{}{\stackrel{~}{\psi }}},\overline{}={\displaystyle \frac{}{\overline{\stackrel{~}{\psi }}}}.`$
When we introduce a set of hermitian two-point functions <sup>1</sup><sup>1</sup>1These two-point functions are different from topological metrics $`\eta _\mathrm{}m=𝒪^{(\mathrm{})}𝒪^{(m)}=N\delta _{\mathrm{}+m,d}`$.
$`\overline{𝒪}^{(\mathrm{})}|𝒪^{(\mathrm{})}=e^q_{\mathrm{}}(0\mathrm{}N2),`$ (2)
the Kähler potential and Zamolodchikov metric $`g_{\stackrel{~}{\psi }\overline{\stackrel{~}{\psi }}}`$ are expressed as
$`e^{q_0}=e^K,g_{\stackrel{~}{\psi }\overline{\stackrel{~}{\psi }}}=e^{q_1q_0},`$
$`\overline{}q_0=e^{q_1q_0}.`$ (3)
The equation Eq.(3) is a part of the $`tt^{}`$-fusion equation of the Calabi-Yau $`d`$-fold with fusion couplings $`\kappa _n`$
$`\overline{}q_0+|\kappa _0|^2e^{q_1q_0}=0,`$
$`\overline{}q_{\mathrm{}}+|\kappa _{\mathrm{}}|^2e^{q_{\mathrm{}+1}q_{\mathrm{}}}|\kappa _\mathrm{}1|^2e^{q_{\mathrm{}}q_\mathrm{}1}=0(1\mathrm{}d1),`$
$`\overline{}q_d|\kappa _{d1}|^2e^{q_dq_{d1}}=0.`$ (4)
This set of equations Eq.(4) represents an $`A_d`$ type Toda system. By introducing new variables
$`\stackrel{~}{q}_0=q_0,`$
$`\stackrel{~}{q}_{\mathrm{}}=q_{\mathrm{}}+{\displaystyle \underset{n=0}{\overset{\mathrm{}1}{}}}\mathrm{log}|\kappa _n|^2(\mathrm{}1),`$
and noting relations $`\overline{}\stackrel{~}{q}_{\mathrm{}}=\overline{}q_{\mathrm{}}`$, we reexpress the Toda system into a formula
$`U_{\mathrm{}}=(\mathrm{}+1)U_0+{\displaystyle \underset{n=0}{\overset{\mathrm{}1}{}}}(\mathrm{}n)D\mathrm{log}(U_n)(1\mathrm{}d1),`$
$`U_0=Dq_0=D\mathrm{log}\left[1+{\displaystyle \underset{\mathrm{}=2}{\overset{N1}{}}}{\displaystyle \frac{I_{\mathrm{}}}{I_1}}\left|{\displaystyle \frac{\stackrel{~}{\varpi }_{\mathrm{}}}{\stackrel{~}{\varpi }_1}}\right|^2\right],`$
$`U_d=0,`$
$`U_n={\displaystyle \underset{m=0}{\overset{n}{}}}Dq_m(1nd),`$
$`D:=\overline{},`$
$`|\kappa _n|^2e^{q_{n+1}q_n}=U_n(0nd1)`$ (5)
The number of diagonal components $`I_k`$ in the Kähler potential is $`(d+1)`$ and the above Toda system with (d+1) equations gives us consistency conditions of the $`I_k`$s. Now recall that the components of the matrix $`I`$ are numerical constants. We may truncate terms of the periods $`\stackrel{~}{\varpi }_k`$s in the series expansions up to order $`𝒪(\psi ^N)`$ for the purpose of determination of $`I`$
$`\stackrel{~}{\varpi }_k(\psi )=f_k(N\psi )^k\left[1+𝒪(\psi ^N)\right](k=1,2,\mathrm{},N1),`$
$`U_0=D\mathrm{log}\left[1+{\displaystyle \underset{\mathrm{}=2}{\overset{N1}{}}}a_{\mathrm{}}(\stackrel{~}{\psi }\overline{\stackrel{~}{\psi }})^\mathrm{}1+\mathrm{}\right].`$
$`f_k:=\left[\mathrm{\Gamma }\left({\displaystyle \frac{k}{N}}\right)\right]^N{\displaystyle \frac{1}{\mathrm{\Gamma }(k)}},a_n:={\displaystyle \frac{I_n}{I_1}}\left|{\displaystyle \frac{f_n}{f_1}}\right|^2.`$ (6)
We calculate the $`U_n`$ concretely for $`n10`$, and propose a conjecture for the $`\{U_n\}`$s at the Gepner point $`\psi =0`$
$`U_0=a_2,U_d=0,`$
$`U_n=(n+1)^2{\displaystyle \frac{a_{n+2}}{a_{n+1}}}(1nd1).`$
When we use definitions of $`a_n`$ and $`f_n`$, the $`U_n`$s are expressed as
$`U_0=a_2={\displaystyle \frac{I_2}{I_1}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{2}{N}\right)}{\mathrm{\Gamma }\left(\frac{1}{N}\right)}}\right]^{2N},`$
$`U_n=(n+1)^2{\displaystyle \frac{a_{n+2}}{a_{n+1}}}={\displaystyle \frac{I_{n+2}}{I_{n+1}}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n+2}{N}\right)}{\mathrm{\Gamma }\left(\frac{n+1}{N}\right)}}\right]^{2N}(1nd1).`$
Also we can obtain expressions of hermitian two-point functions by remarking that the three-point couplings become constants $`\kappa _n=1`$ $`(n=0,1,\mathrm{},d)`$ at the $`\psi =0`$
$`e^{q_{n+1}q_n}={\displaystyle \frac{I_{n+2}}{I_{n+1}}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n+2}{N}\right)}{\mathrm{\Gamma }\left(\frac{n+1}{N}\right)}}\right]^{2N}(0nd1).`$
Equivalently, normalized two-point functions are represented as
$`{\displaystyle \frac{\overline{𝒪}^{(m)}|𝒪^{(m)}}{\overline{𝒪}^{(0)}|𝒪^{(0)}}}=e^{q_mq_0}=(1)^m{\displaystyle \frac{I_{m+1}}{I_1}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{m+1}{N}\right)}{\mathrm{\Gamma }\left(\frac{1}{N}\right)}}\right]^{2N}(0md).`$ (7)
Now let us compare these results with those of CFT calculations for minimal models associated with the $`d`$-fold $`M`$
$`e^{q_n}={\displaystyle \frac{1}{N^N}}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n+1}{N}\right)}{\mathrm{\Gamma }\left(1\frac{n+1}{N}\right)}}\right]^N(n=0,1,2,\mathrm{},N2),`$
$`e^{q_nq_0}=\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n+1}{N}\right)\mathrm{\Gamma }\left(1\frac{1}{N}\right)}{\mathrm{\Gamma }\left(\frac{1}{N}\right)\mathrm{\Gamma }\left(1\frac{n+1}{N}\right)}}\right]^N.`$ (8)
By comparing two results Eq.(7) and Eq.(8), we can obtain the components of the matrix $`I`$ up to one numerical constant $`c`$
$`I_m=c(1)^m\left(\mathrm{sin}{\displaystyle \frac{\pi m}{N}}\right)^N(m=1,2,\mathrm{},N1).`$
In this case, the $`e^K`$ in Eq.(1) is given as
$`e^K={\displaystyle \underset{m=1}{\overset{N1}{}}}c(1)^m\pi ^N\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{m}{N}\right)}{\mathrm{\Gamma }\left(1\frac{m}{N}\right)}}\right]^N\times {\displaystyle \frac{(N^2\psi \overline{\psi })^m}{[\mathrm{\Gamma }(m)]^2}}|G_m(\psi )|^2.`$
Because the Kähler potential is not a function but a section of a line bundle, there is an arbitrariness of multiplication of arbitrary (anti-)holomorphic functions. We choose the normalization factor as
$`c={\displaystyle \frac{1}{\pi ^NN^{N+2}}},`$
then the Kähler potential is written as
$`e^K=(\psi \overline{\psi }){\displaystyle \underset{m=1}{\overset{N1}{}}}(1)^{m1}\left[{\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{m}{N}\right)}{\mathrm{\Gamma }\left(1\frac{m}{N}\right)}}\right]^N{\displaystyle \frac{(N^2\psi \overline{\psi })^{m1}}{[\mathrm{\Gamma }(m)]^2}}|G_m(\psi )|^2.`$ (9)
In order to confirm the validity of this convention, we restrict ourselves to the 3-fold $`(N=5)`$ case and consider a $`3`$-point function $`\kappa `$ of the operator $`𝒪^{(1)}`$ in the B-model
$`\kappa :=\kappa _{\psi \psi \psi }={\displaystyle \frac{1}{5^3}}{\displaystyle \frac{5\psi ^2}{1\psi ^5}}.`$ (10)
(In considering the $`tt^{}`$ equation, we use normalized periods $`\omega _k`$ by taking ratios. In that case, the $`e^K`$ in Eq.(9) and the three-point function Eq.(10) are divided by a factor $`|\psi |^2`$. ) This coupling is a section of holomorphic line bundle and changes with a normalization of the Kähler potential. But there is an invariant $`3`$-point coupling
$`(g_{\psi \overline{\psi }})^{3/2}e^K|\kappa |,`$
and its value is evaluated at the $`\psi =0`$ in our normalization
$`(g_{\psi \overline{\psi }})^{3/2}e^K|\kappa |=\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{3}{5}\right)}{\mathrm{\Gamma }\left(\frac{2}{5}\right)}}\right]^{15/2}\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{5}\right)}{\mathrm{\Gamma }\left(\frac{4}{5}\right)}}\right]^{5/2}=1.5553189899632389725\mathrm{}.`$
This result coincides with that of the calculation of CFT .
## 5 Metric and Curvature
Now we have an exact formula of the $`K`$ with a moduli parameter $`\psi `$ and can evaluate the corrections in the physical observables by the marginal deformation of the CFT. For simplicity, we will consider leading corrections of the Kähler potential, metric and scalar curvature in the moduli space. When we define a function $`A_m`$
$`A_m=(N^2)^{m1}{\displaystyle \frac{1}{[\mathrm{\Gamma }(m)]^2}}\left[{\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{m}{N}\right)}{\mathrm{\Gamma }\left(1\frac{m}{N}\right)}}\right]^N(1mN1),`$
we can obtain these formulae
$`{\displaystyle \frac{A_m}{A_n}}=(N^2)^{mn}\left({\displaystyle \frac{\mathrm{\Gamma }(n)}{\mathrm{\Gamma }(m)}}\right)^2\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{m}{N}\right)\mathrm{\Gamma }\left(1\frac{n}{N}\right)}{\mathrm{\Gamma }\left(1\frac{m}{N}\right)\mathrm{\Gamma }\left(\frac{n}{N}\right)}}\right]^N(m,nN1),`$
$`e^K=\left[{\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{N}\right)}{\mathrm{\Gamma }\left(1\frac{1}{N}\right)}}\right]^N(\psi \overline{\psi })\times \left[1N^2\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{2}{N}\right)\mathrm{\Gamma }\left(1\frac{1}{N}\right)}{\mathrm{\Gamma }\left(1\frac{2}{N}\right)\mathrm{\Gamma }\left(\frac{1}{N}\right)}}\right]^N(\psi \overline{\psi })+\mathrm{}\right](N3),`$
$`g_{\psi \overline{\psi }}=N^2\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{2}{N}\right)\mathrm{\Gamma }\left(1\frac{1}{N}\right)}{\mathrm{\Gamma }\left(1\frac{2}{N}\right)\mathrm{\Gamma }\left(\frac{1}{N}\right)}}\right]^N+(\psi \overline{\psi })\left[2\left({\displaystyle \frac{A_2}{A_1}}\right)^24\left({\displaystyle \frac{A_3}{A_1}}\right)\right]+\mathrm{}(N4),`$
$`R=4+2\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{N}\right)\mathrm{\Gamma }\left(\frac{3}{N}\right)}{\mathrm{\Gamma }\left(1\frac{1}{N}\right)\mathrm{\Gamma }\left(1\frac{3}{N}\right)}}\right]^N\left[{\displaystyle \frac{\mathrm{\Gamma }\left(1\frac{2}{N}\right)}{\mathrm{\Gamma }\left(\frac{2}{N}\right)}}\right]^{2N}`$
$`+(\psi \overline{\psi })\left[24\left({\displaystyle \frac{A_3}{A_2}}\right)96\left({\displaystyle \frac{A_1}{A_2}}\right)\left({\displaystyle \frac{A_3}{A_2}}\right)^2+72\left({\displaystyle \frac{A_1}{A_2}}\right)\left({\displaystyle \frac{A_4}{A_2}}\right)\right]+\mathrm{}(N5).`$
The Ricci tensor of the moduli space is represented as $`R_{\psi \overline{\psi }}={\displaystyle \frac{1}{2}}g_{\psi \overline{\psi }}R`$. For the torus ($`N=3`$) case, its metric is a standard one of the upper-half plane and the scalar curvature is constant $`R=4`$ except for points $`\psi =e^{2\pi i\mathrm{}/3}`$ ($`\mathrm{}=0,1,2`$). Also the K3 ($`N=4`$) metric is given by the above formula and related scalar curvature is a negative constant number $`R=2`$ except for $`\psi =e^{\pi i\mathrm{}/2}`$ ($`\mathrm{}=0,1,2,3`$).
We plot this scalar curvature at the $`\psi =0`$ in Fig.1. It increases monotonically with the dimension $`d`$. But the derivative with respect to $`_\psi _{\overline{\psi }}`$ is negative for $`N11`$ cases at $`\psi =0`$ as shown in Fig.2.
Let us return to the hermitian two-point functions. They are represented as some combinations of these $`K`$, $`g_{\psi \overline{\psi }}`$ and $`R`$
$`e^{\stackrel{~}{q}_0}=e^K,e^{\stackrel{~}{q}_1\stackrel{~}{q}_0}={\displaystyle \frac{1}{N^2}}g_{\psi \overline{\psi }},e^{\stackrel{~}{q}_2\stackrel{~}{q}_1}={\displaystyle \frac{1}{N^2}}g_{\psi \overline{\psi }}\left({\displaystyle \frac{R}{2}}+2\right),`$
$`e^{\stackrel{~}{q}_3\stackrel{~}{q}_2}={\displaystyle \frac{1}{N^2}}g_{\psi \overline{\psi }}\left[3\left({\displaystyle \frac{R}{2}}+1\right)g^{\psi \overline{\psi }}_\psi \overline{}_{\overline{\psi }}\mathrm{log}\left({\displaystyle \frac{R}{2}}+2\right)\right],`$
$`e^{\stackrel{~}{q}_4\stackrel{~}{q}_3}={\displaystyle \frac{1}{N^2}}g_{\psi \overline{\psi }}[4+3R2g^{\psi \overline{\psi }}_\psi \overline{}_{\overline{\psi }}\mathrm{log}({\displaystyle \frac{R}{2}}+2)`$
$`g^{\psi \overline{\psi }}_\psi \overline{}_{\overline{\psi }}\mathrm{log}[3({\displaystyle \frac{R}{2}}+1)g^{\psi \overline{\psi }}_\psi \overline{}_{\overline{\psi }}\mathrm{log}({\displaystyle \frac{R}{2}}+2)]],`$
$`\mathrm{}.`$
Now we know the formula of the $`K`$, $`R`$, and $`g_{\psi \overline{\psi }}`$ and evaluate moduli dependences of these correlators.
In this paper we develop a method to determine the Kähler potential unambiguously by comparing the result of CFT with that of topological sigma model. We calculate the metric and curvature in the neighborhood of the Gepner point, which have dependences of moduli parameter $`\psi `$. The result represents a marginal deformation of the CFT. But the formula of the $`K`$ is exact and we can study it at all points in the moduli space by analytic continuation. The method we developed here is not restricted to the specific model and can be applied to any other Calabi-Yau spaces.
## Acknowledgement
This work is supported by the Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture 10740117. |
warning/0003/hep-th0003220.html | ar5iv | text | # Level-four approximation to the tachyon potential in superstring field theory
## 1 Introduction
It has been conjectured by Sen that at the stationary point of the tachyon potential for the D-brane-anti-D-brane pair or for the non-BPS D-brane of superstring theories, the negative energy density precisely cancels the brane tensions .
For the D-brane of bosonic string theory, this conjecture has been verified by Sen and Zwiebach starting from Witten’s open string field theory . Using a level truncation method which was proposed by Kostelecky and Samuel and including fields up to level 4, they found a contribution of $`99\%`$ of the expected value. Subsequently Taylor and Moeller continued the calculation to level 10 fields and verified the conjecture to $`99,9\%`$ .
In the superstring case the first calculation was done by Berkovits using his Wess-Zumino-Witten like proposal for the string field theory action . The pure tachyon contribution was found to amount to $`60\%`$ of the conjectured value. This computation was then continued to higher levels by the combined force of Berkovits, Sen and Zwiebach . They included fields up to level $`3/2`$ and got $`85\%`$ of the expected answer.
In this paper we perform a further check on the conjecture: we continue the calculation retaining fields up to level 2 and get $`89\%`$ of the conjectured value.
## 2 Berkovits’ superstring field theory
Using the embedding of the $`N=1`$ superstring into a critical $`N=2`$ theory found in , Berkovits proposed a superstring field theory based on a Wess-Zumino-Witten type action . With slight modifications, this action can be used to describe the NS-sector excitations of a non-BPS brane (the modification needed to include the R-sector fields is as yet unknown). In this section we briefly review the action and some of its properties. This section summarizes parts of relevant to the problem at hand with some additional comments.
A string field describing an NS excitation on a non-BPS D-brane can be represented by a vertex operator of the form
$$\widehat{\mathrm{\Phi }}=\mathrm{\Phi }(I\text{ or }\sigma _1)$$
(1)
where $`\mathrm{\Phi }`$ is an operator in the conformal field theory of the NS superstring with the superghost system bosonized as $`\beta =\xi e^\varphi `$ and $`\gamma =\eta e^\varphi `$ .<sup>1</sup><sup>1</sup>1We adopt the convention that $`e^{q\varphi }`$ is a fermion for odd values of $`q`$, i.e. it anticommutes not only with $`e^{q^{}\varphi }`$ for odd $`q^{}`$ but also with the other fermionic fields in the theory. $`\mathrm{\Phi }`$ is restricted to have ghost number 0, picture number 0 and to live in the ‘large’ Hilbert space which includes the $`\xi `$ zero mode. The string field $`\widehat{\mathrm{\Phi }}`$ should include states of both the GSO-projected and GSO-unprojected sectors. Fields in the GSO-unprojected sector are tensored with $`\sigma _1`$ and the fields in the GSO-projected sector are tensored with $`I`$.
One further defines $`\widehat{\eta }_0=\eta _0\sigma _3`$ where $`\eta _0`$ is the zero-mode of the $`\eta `$-field, and $`\widehat{Q}_B=Q_B\sigma _3`$ where $`Q_B`$ is the BRST-charge
$$Q_B=𝑑zj_B(z)=𝑑z\left\{c(T_m+T_{\xi \eta }+T_\varphi )+ccb+\eta e^\varphi G_m\eta \eta e^{2\varphi }b\right\},$$
(2)
and
$$T_{\xi \eta }=\xi \eta ,T_\varphi =\frac{1}{2}\varphi \varphi ^2\varphi ,$$
(3)
$`T_m`$ is the matter stress tensor and $`G_m`$ is the matter supercurrent. The string field action for the non-BPS D-brane then takes the following form:
$$S=\frac{1}{4g^2}(e^{\widehat{\mathrm{\Phi }}}\widehat{Q}_Be^{\widehat{\mathrm{\Phi }}})(e^{\widehat{\mathrm{\Phi }}}\widehat{\eta }_0e^{\widehat{\mathrm{\Phi }}})_0^1𝑑t(e^{t\widehat{\mathrm{\Phi }}}_te^{t\widehat{\mathrm{\Phi }}})\{(e^{t\widehat{\mathrm{\Phi }}}\widehat{Q}_Be^{t\widehat{\mathrm{\Phi }}}),(e^{t\widehat{\mathrm{\Phi }}}\widehat{\eta }_0e^{t\widehat{\mathrm{\Phi }}})\}.$$
(4)
With the double brackets we mean the following:
$$\widehat{A}_1\mathrm{}\widehat{A}_n=\left(\text{Trace of matrices}\right)f_1^{(n)}A_1(0)\mathrm{}f_n^{(n)}A_n(0).$$
(5)
The $`f_k^{(n)}`$ entering in the correlator on the right hand side denote some appropriate conformal transformations . They are defined by
$$f_k^{(n)}(z)=e^{\frac{2\pi i(k1)}{n}}\left(\frac{1+iz}{1iz}\right)^{2/n},\text{for}n1,$$
(6)
they map the unit circle to wedge-formed pieces of the complex-plane.
One can show that the action (4) is invariant under the gauge transformation
$$\delta e^{\widehat{\mathrm{\Phi }}}=\left(\widehat{Q}_B\widehat{\mathrm{\Omega }}\right)e^{\widehat{\mathrm{\Phi }}}+e^{\widehat{\mathrm{\Phi }}}\left(\widehat{\eta }_0\widehat{\mathrm{\Omega }}^{}\right),$$
(7)
where $`\widehat{\mathrm{\Omega }}`$ and $`\widehat{\mathrm{\Omega }}^{}`$ are independent gauge parameters. This gauge invariance can be fixed<sup>2</sup><sup>2</sup>2This is a reachable gauge choice (if $`L_00`$) but we have not been able to prove that it fixes the gauge freedom completely. by imposing
$$b_0|\text{State}=0\text{ and }\xi _0|\text{State}=0.$$
(8)
In the calculation of the tachyon potential, we can restrict the string field to lie in a subspace $`_1`$ formed by acting only with modes of the stress-energy tensor, the supercurrent and the ghost fields $`b`$, $`c`$, $`\eta `$, $`\xi `$, $`\varphi `$, since the other excitations can be consistently put to zero. Furthermore, when restricted to fields lying in $`_1`$ the action has a $`Z_2`$ twist invariance under which the fields in the GSO-odd sector carry charge $`()^{h+1}`$ and the fields in the GSO-even sector carry charge $`()^{h+1/2}`$, $`h`$ is the conformal weight. In the computation of the tachyon potential we can therefore further restrict ourselves to the twist even fields.<sup>3</sup><sup>3</sup>3This restriction projects out the only state with $`L_0=0`$, namely $`\xi \xi cce^{2\varphi }`$.
The non-polynomial action (4) should be formally expanded in the string field $`\widehat{\mathrm{\Phi }}`$, and each term should be accompanied by the appropriate conformal transformations. However, because we will only compute the interactions between a finite number of fields, it is easy to see that one does not need all the terms in the action. The conformal field theory correlators in the action (4) are nonvanishing only if the total $`(b,c)`$ number is $`3`$, the $`(\eta ,\xi )`$ number is 1 and the total $`\varphi `$-charge is $`2`$. In the following we will need only the terms in the action involving up to $`6`$ string fields.
Making use of the cyclicity properties derived in the appendix of , the action to this order can be written as
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2g^2}}{\displaystyle \frac{1}{2}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+{\displaystyle \frac{1}{3}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})+{\displaystyle \frac{1}{12}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }})+`$
$`+{\displaystyle \frac{1}{60}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})\left(\widehat{\mathrm{\Phi }}^3(\widehat{\eta }_0\widehat{\mathrm{\Phi }})3\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}\right)+`$
$`+{\displaystyle \frac{1}{360}}(\widehat{Q}_B\widehat{\mathrm{\Phi }})(\widehat{\mathrm{\Phi }}^4(\widehat{\eta }_0\widehat{\mathrm{\Phi }})4\widehat{\mathrm{\Phi }}^3(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}+3\widehat{\mathrm{\Phi }}^2(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\widehat{\mathrm{\Phi }}^2).`$
## 3 The fields up to level 2
Taking all this together we get the list of contributing fields up to level 2 (table 1). The level of a field is just the conformal weight shifted by 1/2. In this way the tachyon is a level 0 field. We use the notation $`|q`$ for the state corresponding with the operator $`e^{q\varphi }`$.
The level $`0`$ and level $`2`$ fields should be tensored with $`\sigma _1`$ and the level $`3/2`$ fields with $`I`$.
We list the conformal transformations of the fields in appendix A.
## 4 The tachyon potential
We have calculated the tachyon potential involving fields up to level 2, including only the terms up to level 4 (the level of a term in the potential is defined to be the sum of the levels of the fields entering into it). We have performed the actual computation in the following way. The conformal tranformations of the fields were calculated by hand. The computation of all the correlation functions between these transformed fields was done with the help of Mathematica: we have written a program to compute the necessary CFT correlation functions. We have performed an extra check by calculating some of the correlators on the upper half plane instead of the disc. Denoting
$$\widehat{\mathrm{\Phi }}=t\widehat{T}+a\widehat{A}+e\widehat{E}+f\widehat{F}+k\widehat{K}+l\widehat{L}+m\widehat{M}+n\widehat{N}+p\widehat{P}$$
(9)
we give the result with coefficients evaluated numerically up to 6 significant digits (subscripts refer to the level)
$`V(\widehat{\mathrm{\Phi }})`$ $`=`$ $`S(\widehat{\mathrm{\Phi }})`$
$`=`$ $`V_0+V_{3/2}+V_2+V_3+V_{7/2}+V_4`$
$`V(\widehat{\mathrm{\Phi }})_0`$ $`=`$ $`2\pi ^2M(0.25t^2+0.5t^4)`$
$`V(\widehat{\mathrm{\Phi }})_{3/2}`$ $`=`$ $`2\pi ^2M(at^20.25et^20.518729et^4)`$
$`V(\widehat{\mathrm{\Phi }})_2`$ $`=`$ $`2\pi ^2M\left(0.333333kt^3+1.83333lt^33.75mt^3+2.83333nt^3+0.25pt^3\right)`$
$`V(\widehat{\mathrm{\Phi }})_3`$ $`=`$ $`2\pi ^2M(2ae+5f^2+4.96405aet^20.66544e^2t^2+`$
$`+5.47589eft^2+5.82107f^2t^2+0.277778e^2t^4)`$
$`V(\widehat{\mathrm{\Phi }})_{7/2}`$ $`=`$ $`2\pi ^2M(3.03704akt7.11111alt+2.77778amt1.62963ant`$
$`1.55556apt+0.12963ekt0.296296elt+0.694444emt`$
$`1.2963ent+0.944444ept11.8519flt8.88889fmt`$
$`2.96296fpt2.87299ekt^31.94348elt^3+4.35732emt^3`$
$`4.77364ent^3+0.605194ept^3)`$
$`V(\widehat{\mathrm{\Phi }})_4`$ $`=`$ $`2\pi ^2M(3k^23kl+1.5l^2+5.625m^23n^20.75p^2+10.3958k^2t^2+`$
$`+0.791667klt^21.875kmt^2+5.54167knt^21.4375kpt^2+`$
$`+6.70833l^2t^210.3125lmt^2+11.9167lnt^20.875lpt^2+`$
$`+14.7656m^2t^215.9375mnt^21.40625mpt^2+`$
$`+5.83333n^2t^20.5npt^21.5p^2t^2).`$
Our results for $`V_0`$, $`V_{3/2}`$ and $`V_3`$ agree completely with <sup>4</sup><sup>4</sup>4Note that our field $`F`$ is defined with a different sign from the one in .
This potential has extrema at $`(\pm t_0,a_0,e_0,f_0,\pm k_0,\pm l_0,\pm m_0,\pm n_0,\pm p_0)`$ with
$$\begin{array}{ccccccccccc}\hfill t_0& =& 0.602101\hfill & & & & & & & & \\ \hfill a_0& =& 0.0521934,\hfill & & \hfill e_0& =& 0.0430366,\hfill & & \hfill f_0& =& 0.0138164,\hfill \\ \hfill k_0& =& 0.01019,\hfill & & \hfill l_0& =& 0.0450433,\hfill & & \hfill m_0& =& 0.0322127,\hfill \\ \hfill n_0& =& 0.0473113,\hfill & & \hfill p_0& =& 0.021291.\hfill & & & & \end{array}$$
At these extrema
$$V=0.891287M,$$
(10)
so we see that at the level $`4`$ we get $`89\%`$ of the conjectured value $`V=M`$ ($`M`$ is the D-brane mass). In the potential computed to level $`4`$ all the fields but $`t`$ appear only quadratically. They can be integrated out very easily to give the effective potential $`V(t)`$ see figure 1.
## 5 Conclusions and outlook
The last few months have seen an increased confidence in string field theory as a calculational tool for doing off-shell string calculations. In this letter, we have used Berkovits’ superstring field theory to calculate the tachyon potential up to level four. Our result shows a further convergence towards the total vacuum energy conjectured by Sen, albeit less rapid than the contributions of the previous levels. Therefore it would be interesting to persue the calculation to higher levels. At present, not much is known about the general convergence properties of level-truncation calculations in superstring theory. It would be nice to have a deeper understanding of this. Another interesting problem would be to study the interactions of the massless vector using Berkovits’ action and compare with the different proposals for the non-abelian Dirac-Born-Infeld action .
###### Acknowledgments.
This work was supported in part by the European Commission TMR projectERBFMRXCT96-0045. We would like to thank N. Berkovits and W. Taylor for encouragement. We have benefitted from discussions with B. Craps, F. Roose, J. Troost and W. Troost, and from correspondence with Amer Iqbal and Asad Naqvi. P.J.D.S. is aspirant FWO-Vlaanderen.
## Appendix A The conformal transformations of the fields
We now list the conformal transformations of the fields used in the calculation of the tachyon potential. To shorten the notation we denote $`w=f(z)`$.
$`fT(z)`$ $`=`$ $`(f^{}(z))^{1/2}T(w)`$
$`fA(z)`$ $`=`$ $`f^{}(z)A(w){\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}cc\xi \xi e^{2\varphi }(w)`$
$`fE(z)`$ $`=`$ $`f^{}(z)E(w){\displaystyle \frac{f^{\prime \prime }(z)}{2f^{}(z)}}`$
$`fF(z)`$ $`=`$ $`f^{}(z)F(w)`$
$`fK(z)`$ $`=`$ $`(f^{}(z))^{3/2}K(w)+2{\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}\xi c\left(e^\varphi \right)(w)+`$
$`+\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2\right](f^{}(z))^{1/2}\xi ce^\varphi (w)`$
$`fL(z)`$ $`=`$ $`(f^{}(z))^{3/2}L(w)+{\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}\xi c\varphi e^\varphi (w)+`$
$`+\left[{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2{\displaystyle \frac{2}{3}}{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}\right](f^{}(z))^{1/2}\xi ce^\varphi (w)`$
$`fM(z)`$ $`=`$ $`(f^{}(z))^{3/2}M(w)+{\displaystyle \frac{15}{12}}\left[{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}{\displaystyle \frac{3}{2}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2\right](f^{}(z))^{1/2}\xi ce^\varphi (w)`$
$`fN(z)`$ $`=`$ $`(f^{}(z))^{3/2}N(w){\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}\xi ce^\varphi (w)+`$
$`+\left[2\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}\right](f^{}(z))^{1/2}\xi ce^\varphi (w)`$
$`fP(z)`$ $`=`$ $`(f^{}(z))^{3/2}P(w)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}\xi ce^\varphi (w)+`$
$`+\left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2{\displaystyle \frac{1}{6}}{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}\right](f^{}(z))^{1/2}\xi ce^\varphi (w)`$ |
warning/0003/math0003040.html | ar5iv | text | # On two-parameter deformations of 𝑜𝑠𝑝(1|2)⁽¹⁾
ABSTRACT
An elliptic two-parameter deformation of the (universal enveloping superalgebra of) affine Lie superalgebra $`osp(1|2)^{(1)}`$ is proposed in terms of free boson realization. This deformed superalgebra is shown to fit in the framework of infinite Hopf family of superalgebras, a generalization of the infinite Hopf family of algebras proposed earlier by the authors. The trigonometric and rational degenerations are briefly discussed.
It has been becoming more and more evident that quantum affine and quantum Virasoro algebras play some essential roles in the theories of massive integrable quantum fields in (1+1)-dimensions and in 2d statistical mechanical systems off the criticality. For some systems, similar roles are played by some yet more complicated algebraic structures, e.g. some two-parameter deformations of affine Lie algebras. It has been realized that there are several kinds of two-parameter deformations of affine Lie algebras, including the standard elliptic quantum groups proposed by Felder et al. and some different variants thereof . All these two-parameter algebras fall into one of the following two classes: one is the quasi-triangular quasi-Hopf algebra –certain twists of the standard Hopf algebra structure –and the other is the so-called infinite Hopf family of algebras . To us the relation between the structures of quasi-triangular quasi-Hopf algebras and infinite Hopf family of algebras still remains a mystery because these two structures are defined respectively for algebras given in different realizations: the quasi-triangular quasi-Hopf algebra is introduced in the context of Yang-Baxter realization (or RS-realization ) and the co-structure is a deformation of the standard co-algebraic structure of quantum affine algebras, while the infinite Hopf family of algebras is introduced for algebras given in the Drinfeld new current realization only, and the co-structure is a deformation of Drinfeld’s new co-structure for quantum affine algebras .
Among the above mentioned two-parameter deformations of affine algebras, we are particularly interested in the algebras $`_{q,p}(\widehat{g})`$ studied in because, at level $`c=1`$, such algebras are intimately related to the elliptic quantum $`W`$-algebras proposed by B. Feigin and E.Frenkel : the former are basically the algebra of screening currents of the latter, with the introduction of some auxiliary semi-simple currents and perhaps some dynamical shift.
In this article, we are aiming at proposing a generalization of the algebra $`_{q,p}(\widehat{g})`$ to the case of Lie superalgebra, with $`g`$ identified as the simplest Lie superalgebra $`osp(1|2)`$. One of the reason to do such a work is due to the fact that, although the elliptic quantum $`W`$-algebras have been fairly well understood , their super-correspondences has not been properly studied, even the structure of the simplest super Virasoro algebra(if any) is still unknown. From the experiences of purely bosonic case, it may be reasonable to expect that, the elliptic quantum super Virasoro algebra, if exists, will possess an algebra of screening currents which is based on the underlying affine Lie superalgebra $`osp(1|2)^{(1)}`$ and has a similar structure as that of $`_{q,p}(\widehat{g}).`$ Therefore, the study of a two-parameter elliptic quantum superalgebra might hopefully give some hints about the unknown elliptic quantum super Virasoro algebra. On the other hand, the study of such an algebra is of its own interests also: it will provide examples of super analogues of two-parameter quantum deformed affine algebras which naturally fit in the framework of vertex operator algebras studied enthusiastically by pure mathematicians. The existence, structure and representation theories of such algebras are all worth studying.
The logic we shall be going along with is as follows. First we show that we can close a set of currents into an associative superalgebra starting from some free boson expressions. Then we show that this current superalgebra is a particular case (level $`c=1`$) of some more general, abstract superalgebra $`_{q,p}(osp(1|2)^{(1)})`$ associated with $`osp(1|2)^{(1)}`$. Next, we prove that the superalgebra $`_{q,p}(osp(1|2)^{(1)})`$ has a well-defined co-structure, which can be regarded as a super analogue of the structure of the so-called infinite Hopf family of algebras defined earlier by the authors for $`_{q,p}(\widehat{g})`$. Last, we show that the superalgebra $`_{q,p}(osp(1|2)^{(1)})`$ has some interesting scaling limit, denoted as $`𝒜_{\mathrm{},\eta }(osp(1|2)^{(1)})`$, which is the analogue of $`𝒜_{\mathrm{},\eta }(\widehat{g}),`$ the trigonometric degeneration of $`_{q,p}(\widehat{g})`$.
To begin with, we introduce the following Heisenberg algebra $``$ with generators {$`a_i,`$ $`iZ`$} and relations
$`[a_n,a_m]`$ $`=`$ $`{\displaystyle \frac{1}{n}}\left(q^nq^n\right)\left(\left(qp\right)^n\left(qp\right)^n\right)\left(p^n+p^n1\right)\delta _{n+m,0},(n0)`$
$`[P,Q]`$ $`=`$ $`1.`$
Let
$$s_n^+=\frac{a_n}{q^nq^n},s_n^{}=\frac{a_n}{\left(qp\right)^n\left(qp\right)^n}.$$
Define
$`\phi (z)={\displaystyle \underset{n0}{}}s_n^+z^n,\psi (z)={\displaystyle \underset{n0}{}}s_n^{}z^n,`$ (1)
we have
$`\phi (z)\phi (w)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{\left(\left(qp\right)^n\left(qp\right)^n\right)\left(p^n+p^n1\right)}{q^nq^n}}\left({\displaystyle \frac{w}{z}}\right)^n,`$
$`\psi (z)\psi (w)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{\left(q^nq^n\right)\left(p^n+p^n1\right)}{\left(qp\right)^n\left(qp\right)^n}}\left({\displaystyle \frac{w}{z}}\right)^n,`$
$`\phi (z)\psi (w)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\left(p^n+p^n1\right)\left({\displaystyle \frac{w}{z}}\right)^n.`$
Therefore,
$`\mathrm{exp}\left\{\phi (z)\phi (w)\right\}`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{\left(\left(qp\right)^n\left(qp\right)^n\right)\left(p^n+p^n1\right)}{q^nq^n}}\left({\displaystyle \frac{w}{z}}\right)^n\right\}`$
$`=`$ $`{\displaystyle \frac{\left(\frac{w}{z}p^2|q^2\right)_{\mathrm{}}\left(\frac{w}{z}q^2p|q^2\right)_{\mathrm{}}\left(\frac{w}{z}|q^2\right)_{\mathrm{}}}{\left(\frac{w}{z}(qp)^2|q^2\right)_{\mathrm{}}\left(\frac{w}{z}p^1|q^2\right)_{\mathrm{}}\left(\frac{w}{z}q^2|q^2\right)_{\mathrm{}}}};`$
$`\mathrm{exp}\left\{\psi (z)\psi (w)\right\}`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{\left(q^nq^n\right)\left(p^n+p^n1\right)}{\left(qp\right)^n\left(qp\right)^n}}\left({\displaystyle \frac{w}{z}}\right)^n\right\}`$
$`=`$ $`{\displaystyle \frac{\left(\frac{w}{z}p^2|\left(qp\right)^2\right)_{\mathrm{}}\left(\frac{w}{z}\left(qp\right)^2p^1|\left(qp\right)^2\right)_{\mathrm{}}\left(\frac{w}{z}|\left(qp\right)^2\right)_{\mathrm{}}}{\left(\frac{w}{z}\left(qp\right)^2p^2|\left(qp\right)^2\right)_{\mathrm{}}\left(\frac{w}{z}p|\left(qp\right)^2\right)_{\mathrm{}}\left(\frac{w}{z}\left(qp\right)^2|\left(qp\right)^2\right)_{\mathrm{}}}};`$
$`\mathrm{exp}\left\{\phi (z)\psi (w)\right\}`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\left(p^n+p^n1\right)\left({\displaystyle \frac{w}{z}}\right)^n\right\}`$
$`=`$ $`\mathrm{exp}\left\{\mathrm{log}\left(1{\displaystyle \frac{w}{z}}p\right)+\mathrm{log}\left(1{\displaystyle \frac{w}{z}}p^1\right)\mathrm{log}\left(1{\displaystyle \frac{w}{z}}\right)\right\}`$
$`=`$ $`{\displaystyle \frac{\left(1\frac{w}{z}p\right)\left(1\frac{w}{z}p^1\right)}{\left(1\frac{w}{z}\right)}},`$
where $`(z|q)_{\mathrm{}}=_{n=0}^{\mathrm{}}(1zq^n)`$. Now define
$`S^+(z)=:\mathrm{exp}\left[\phi (z)\right]:,S^{}(z)=:\mathrm{exp}[\psi (z)]:,`$ (2)
where : : means the usual normal ordering of bosonic oscillators, we have
$`S^+(z)S^+(w)`$ $`=`$ $`\mathrm{exp}\left\{\phi (z)\phi (w)\right\}:S^+(z)S^+(w):,`$
$`S^{}(z)S^{}(w)`$ $`=`$ $`\mathrm{exp}\left\{\psi (z)\psi (w)\right\}:S^{}(z)S^{}(w):,`$
$`S^+(z)S^{}(w)`$ $`=`$ $`\mathrm{exp}\left\{\phi (z)\psi (w)\right\}:S^+(z)S^{}(w):.`$
Notice that the above bosonic expressions do not contain the contribution of the zero mode and therefore do not live in the complete Fock space corresponding to the Heisenberg algebra. Adding the zero mode contributions, we now introduce
$`E(z)=e^Qz^PS^+(z),F(z)=e^Qz^PS^{}(z).`$ (3)
Then, it is easy to calculate that
$`E(z)E(w)`$ $`=`$ $`\left(zw\right){\displaystyle \frac{\left(\frac{w}{z}p^2|q^2\right)_{\mathrm{}}\left(\frac{w}{z}q^2p|q^2\right)_{\mathrm{}}}{\left(\frac{w}{z}(qp)^2|q^2\right)_{\mathrm{}}\left(\frac{w}{z}p^1|q^2\right)_{\mathrm{}}}}:E(z)E(w):,`$
$`F(z)F(w)`$ $`=`$ $`\left(zw\right){\displaystyle \frac{\left(\frac{w}{z}p^2|\left(qp\right)^2\right)_{\mathrm{}}\left(\frac{w}{z}\left(qp\right)^2p^1|\left(qp\right)^2\right)_{\mathrm{}}}{\left(\frac{w}{z}\left(qp\right)^2p^2|\left(qp\right)^2\right)_{\mathrm{}}\left(\frac{w}{z}p|\left(qp\right)^2\right)_{\mathrm{}}}}:F(z)F(w):,`$
$`E(z)F(w)`$ $`=`$ $`{\displaystyle \frac{\left(zw\right)}{z^2\left(1\frac{w}{z}p\right)\left(1\frac{w}{z}p^1\right)}}:E(z)F(w):.`$
For $`z=w`$, the OPEs $`E(z)E(w)`$ and $`F(z)F(w)`$ are both zero, showing that the currents $`E(z),F(w)`$ are essentially fermionic. Turning the above equations into commutator type, we have
$`E(z)E(w)+{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{z}{w}p^1\right)}{\theta _{q^2}\left(\frac{z}{w}p^2\right)\theta _{q^2}\left(\frac{w}{z}p^1\right)}}E(w)E(z)`$ $`=`$ $`0,`$
$`F(z)F(w)+{\displaystyle \frac{\theta _{(qp)^2}\left(\frac{w}{z}p^2\right)\theta _{(qp)^2}\left(\frac{z}{w}p\right)}{\theta _{(qp)^2}\left(\frac{z}{w}p^2\right)\theta _{(qp)^2}\left(\frac{w}{z}p\right)}}F(w)F(z)`$ $`=`$ $`0,`$
or alternatively,
$`E(z)E(w)+p{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p\right)}{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p^1\right)}}E(w)E(z)`$ $`=`$ $`0,`$
$`F(z)F(w)+p^1{\displaystyle \frac{\theta _{(qp)^2}\left(\frac{w}{z}p^2\right)\theta _{(qp)^2}\left(\frac{w}{z}p^1\right)}{\theta _{(qp)^2}\left(\frac{w}{z}p^2\right)\theta _{(qp)^2}\left(\frac{w}{z}p\right)}}F(w)F(z)`$ $`=`$ $`0,`$
where $`\theta _q(z)=(z|q)_{\mathrm{}}(z^1q|q)_{\mathrm{}}(q|q)_{\mathrm{}}`$ is essentially the usual Jacobi $`\theta `$-function<sup>1</sup><sup>1</sup>1 To be precise, we have $`\theta _1(u,\tau )=iq^{1/8}z^{1/2}\theta _q(z)`$, where $`z=e^{2i\pi u},q=e^{2\pi i\tau }`$.. Meanwhile, we also have
$`\{E(z),F(w)\}`$ $`=`$ $`{\displaystyle \frac{zw}{(pp^1)zw}}\left\{\delta \left({\displaystyle \frac{z}{wp}}\right)\delta \left({\displaystyle \frac{w}{zp}}\right)\right\}:E(z)F(w):`$
$`=`$ $`{\displaystyle \frac{1}{(pp^1)}}\left\{{\displaystyle \frac{p1}{wp}}\delta \left({\displaystyle \frac{z}{wp}}\right)+{\displaystyle \frac{p1}{zp}}\delta \left({\displaystyle \frac{w}{zp}}\right)\right\}:E(z)F(w):`$
$`=`$ $`{\displaystyle \frac{1}{p^{1/2}+p^{1/2}}}\left\{(wp^{1/2})^1\delta \left({\displaystyle \frac{z}{wp}}\right)+(zp^{1/2})^1\delta \left({\displaystyle \frac{w}{zp}}\right)\right\}:E(z)F(w):.`$
Therefore, defining
$$H^\pm (z)=z^1:E(zp^{\pm 1/2})F(zp^{1/2}):,$$
we get
$$\{E(z),F(w)\}=\frac{1}{p^{1/2}+p^{1/2}}\left\{\delta \left(\frac{z}{wp}\right)H^+(wp^{1/2})+\delta \left(\frac{w}{zp}\right)H^{}(zp^{1/2})\right\}$$
and
$`H^+(z)E(w)`$ $`=`$ $`p{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{q^2}\left(\frac{w}{z}pp^{1/2}\right)}{\theta _{q^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{q^2}\left(\frac{w}{z}p^1p^{1/2}\right)}}E(w)H^+(z),`$
$`H^{}(z)E(w)`$ $`=`$ $`p{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{q^2}\left(\frac{w}{z}pp^{1/2}\right)}{\theta _{q^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{q^2}\left(\frac{w}{z}p^1p^{1/2}\right)}}E(w)H^+(z),`$
$`H^+(z)F(w)`$ $`=`$ $`p^1{\displaystyle \frac{\theta _{(qp)^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{(qp)^2}\left(\frac{w}{z}p^1p^{1/2}\right)}{\theta _{(qp)^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{(qp)^2}\left(\frac{w}{z}pp^{1/2}\right)}}F(w)H^+(z),`$
$`H^{}(z)F(w)`$ $`=`$ $`p^1{\displaystyle \frac{\theta _{(qp)^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{(qp)^2}\left(\frac{w}{z}p^1p^{1/2}\right)}{\theta _{(qp)^2}\left(\frac{w}{z}p^2p^{1/2}\right)\theta _{(qp)^2}\left(\frac{w}{z}pp^{1/2}\right)}}F(w)H^{}(z),`$
$`H^\pm (z)H^\pm (w)`$ $`=`$ $`{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p\right)}{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p^1\right)}}{\displaystyle \frac{\theta _{(qp)^2}\left(\frac{w}{z}p^2\right)\theta _{(qp)^2}\left(\frac{z}{w}p\right)}{\theta _{(qp)^2}\left(\frac{z}{w}p^2\right)\theta _{(qp)^2}\left(\frac{w}{z}p\right)}}H^\pm (w)H^\pm (z),`$
$`H^+(z)H^{}(w)`$ $`=`$ $`{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2p^1\right)\theta _{q^2}\left(\frac{w}{z}pp^1\right)}{\theta _{q^2}\left(\frac{w}{z}p^2p^1\right)\theta _{q^2}\left(\frac{w}{z}p^1p^1\right)}}{\displaystyle \frac{\theta _{(qp)^2}\left(\frac{w}{z}p^2p\right)\theta _{(qp)^2}\left(\frac{z}{w}pp\right)}{\theta _{(qp)^2}\left(\frac{z}{w}p^2p\right)\theta _{(qp)^2}\left(\frac{w}{z}pp\right)}}H^{}(w)H^+(z).`$
Notice that unlike $`E(z)`$ and $`F(z)`$, the currents $`H^\pm (z)`$ are bosonic. For later reference, we introduce the Grassmann parity operator $`\pi `$ such that
$`\pi [E(z)]`$ $`=`$ $`\pi [F(z)]=1,`$
$`\pi [H^\pm (z)]`$ $`=`$ $`0.`$
Then we see that the currents $`H^\pm (z),`$ $`E(z)`$ and $`F(z)`$ close into a superalgebra with $`Z_2`$ gradation provided by the parity operator $`\pi `$.
Definition 1: *The superalgebra* $`_{q,p}(osp(1|2)^{(1)})`$ *is a* $`Z_2`$*-graded associative algebra generated by the unit* $`1`$*, coefficients of the formal power series* $`H^\pm (z)`$ *(invertible),* $`E(z)`$ *and* $`F(z)`$ *in* $`z`$ *and the central element* $`c`$ *with relations*
$`H^+(z)E(w)`$ $`=`$ $`p{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{q^2}\left(\frac{w}{z}pp^{c/2}\right)}{\theta _{q^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{q^2}\left(\frac{w}{z}p^1p^{c/2}\right)}}E(w)H^+(z),`$
$`H^{}(z)E(w)`$ $`=`$ $`p{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{q^2}\left(\frac{w}{z}pp^{c/2}\right)}{\theta _{q^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{q^2}\left(\frac{w}{z}p^1p^{c/2}\right)}}E(w)H^+(z),`$
$`H^+(z)F(w)`$ $`=`$ $`p^1{\displaystyle \frac{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^1p^{c/2}\right)}{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}pp^{c/2}\right)}}F(w)H^+(z),`$
$`H^{}(z)F(w)`$ $`=`$ $`p^1{\displaystyle \frac{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^1p^{c/2}\right)}{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2p^{c/2}\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}pp^{c/2}\right)}}F(w)H^{}(z),`$
$`H^\pm (z)H^\pm (w)`$ $`=`$ $`{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p\right)}{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p^1\right)}}{\displaystyle \frac{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2\right)\theta _{\stackrel{~}{q}^2}\left(\frac{z}{w}p\right)}{\theta _{\stackrel{~}{q}^2}\left(\frac{z}{w}p^2\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p\right)}}H^\pm (w)H^\pm (z),`$
$`H^+(z)H^{}(w)`$ $`=`$ $`{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2p^c\right)\theta _{q^2}\left(\frac{w}{z}pp^c\right)}{\theta _{q^2}\left(\frac{w}{z}p^2p^c\right)\theta _{q^2}\left(\frac{w}{z}p^1p^c\right)}}{\displaystyle \frac{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2p^c\right)\theta _{\stackrel{~}{q}^2}\left(\frac{z}{w}pp^c\right)}{\theta _{\stackrel{~}{q}^2}\left(\frac{z}{w}p^2p^c\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}pp^c\right)}}H^{}(w)H^+(z),`$
$`E(z)E(w)`$ $`=`$ $`p{\displaystyle \frac{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p\right)}{\theta _{q^2}\left(\frac{w}{z}p^2\right)\theta _{q^2}\left(\frac{w}{z}p^1\right)}}E(w)E(z),`$
$`F(z)F(w)`$ $`=`$ $`p^1{\displaystyle \frac{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^1\right)}{\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p^2\right)\theta _{\stackrel{~}{q}^2}\left(\frac{w}{z}p\right)}}F(w)F(z),`$
$`\{E(z),F(w)\}`$ $`=`$ $`{\displaystyle \frac{1}{p^{1/2}+p^{1/2}}}\left\{\delta \left({\displaystyle \frac{z}{wp^c}}\right)H^+(wp^{c/2})+\delta \left({\displaystyle \frac{w}{zp^c}}\right)H^{}(zp^{c/2})\right\},`$
*where* $`\stackrel{~}{q}=qp^c`$. $`\mathrm{}`$
Obviously, we have
Proposition 1: *Equations (1, 2, 3) give a free boson realization for the superalgebra* $`_{q,p}(osp(1|2)^{(1)})`$ *at level* $`c=1`$*.* $`\mathrm{}`$
To understand the properties of the superalgebra $`_{q,p}(osp(1|2)^{(1)})`$ at levels other than $`c=1`$, one natural way is to consider its co-structure. Since the superalgebra $`_{q,p}(osp(1|2)^{(1)})`$ has some similar properties as those of the $`_{q,p}(\widehat{g})`$, e.g. the Cartan involution is broken due to the different periods of the structure functions for the $`EE`$ and $`FF`$ relations, we hope that the structure of infinite Hopf family of algebras for the latter also holds for $`_{q,p}(osp(1|2)^{(1)})`$. It is indeed so, however, the definition of infinite Hopf family of algebras has to be modified into an infinite Hopf family of superalgebras, as is expected naturally. This generalized co-structure for $`_{q,p}(osp(1|2)^{(1)})`$ will be given in the following proposition.
Let us first prepare some notations. Let $`𝒥`$ be an additive semigroup which may be identified with the set of non-negative integer numbers. $`c_n𝒥`$ are elements of $`𝒥`$. Let $`q^{(0)}=q`$ and define $`q^{(n+1)}=q^{(n)}p^{c_n}`$ iteratively. We set $`𝒜_n=_{q^{(n)},p}(osp(1|2)^{(1)})`$ whose generator are denoted $`H^\pm (z;q^{(n)}),`$ $`E(z;q^{(n)}),`$ $`F(z;q^{(n)})`$ and $`c_n`$ respectively. The generating relations for each $`𝒜_n`$ are nothing but those of $`_{q,p}(osp(1|2)^{(1)})`$ with parameters $`q,\stackrel{~}{q}`$ replaced by $`q^{(n)},q^{(n+1)}`$ respectively.
Let $`\{v_i^{(n)},i=1,\mathrm{},dim(𝒜_n)\}`$ be a basis of $`𝒜_n`$. The maps
$`\tau _n^\pm :𝒜_n`$ $``$ $`𝒜_{n\pm 1}`$
$`v_i^{(n)}`$ $``$ $`v_i^{(n\pm 1)}`$
are morphisms from $`𝒜_n`$ to $`𝒜_{n\pm 1}`$. For any two integers $`n,m`$ with $`n<m`$, we can specify a pair of morphisms
$`Mor(𝒜_m,𝒜_n)\tau ^{(m,n)}\tau _{m1}^+\mathrm{}\tau _{n+1}^+\tau _n^+:𝒜_n𝒜_m,`$
$`Mor(𝒜_n,𝒜_m)\tau ^{(n,m)}\tau _{n+1}^{}\mathrm{}\tau _{m1}^{}\tau _m^{}:𝒜_m𝒜_n`$
with $`\tau ^{(m,n)}\tau ^{(n,m)}=id_m,\tau ^{(n,m)}\tau ^{(m,n)}=id_n`$. Clearly the morphisms $`\tau ^{(m,n)},n,mZ`$ satisfy the associativity condition $`\tau ^{(m,p)}\tau ^{(p,n)}=\tau ^{(m,n)}`$ and thus make the family of superalgebras $`\{𝒜_n,nZ\}`$ into a category.
The following definition is a straightforward generalization of the structure of infinite Hopf family of algebras originally presented in :
Definition 2: *The category of superalgebras* $`\{𝒜_n,\{\tau ^{(n,m)}\},n,mZ\}`$ *is called an infinite Hopf family of superalgebras if on each object* $`𝒜_n`$ *of the category one can define the morphisms* $`\mathrm{\Delta }_n^+:𝒜_n𝒜_n𝒜_{n+1}`$*,* $`\mathrm{\Delta }_n^{}:𝒜_n𝒜_{n1}𝒜_n`$*,* $`ϵ_n:𝒜_nC`$ *and antimorphisms* $`S_n^\pm :𝒜_n𝒜_{n\pm 1}`$ *such that the following axioms hold,*
* $`(ϵ_nid_{n+1})\mathrm{\Delta }_n^+=\tau _n^+,(id_{n1}ϵ_n)\mathrm{\Delta }_n^{}=\tau _n^{}`$ *(a1)*
* $`m_{n+1}(S_n^+id_{n+1})\mathrm{\Delta }_n^+=ϵ_{n+1}\tau _n^+,m_{n1}(id_{n1}S_n^{})\mathrm{\Delta }_n^{}=ϵ_{n1}\tau _n^{}`$ *(a2)*
* $`(\mathrm{\Delta }_n^{}id_{n+1})\mathrm{\Delta }_n^+=(id_{n1}\mathrm{\Delta }_n^+)\mathrm{\Delta }_n^{}`$ *(a3)*
*in which* $`m_n`$ *is the (super)multiplication for* $`𝒜_n`$*.* $`\mathrm{}`$
Notice that throughout this article, the symbol $``$ denotes a graded direct product, or *direct super-product*, obeying, e.g. for elements $`A,B,C,D`$ with definite Grassmann parity,
$`(AB)(CD)=(1)^{\pi (B)\pi (C)}ACBD.`$
Proposition 2: *The family of superalgebras* $`\{𝒜_n,`$$`nZ\}`$ *form an Infinite Hopf family of algebras with comultiplications* $`\mathrm{\Delta }_n^\pm ,`$ *counits* $`ϵ_n`$ *and antipodes* $`S_n^\pm `$ *given as follows,*
* *the comultiplications* $`\mathrm{\Delta }_n^\pm `$*:*
$`\mathrm{\Delta }_n^+c_n`$ $`=`$ $`c_n+c_{n+1},`$
$`\mathrm{\Delta }_n^+H^+(z;q^{(n)})`$ $`=`$ $`H^+(zp^{c_{n+1}/2};q^{(n)})H^+(zp^{c_n/2};q^{(n+1)}),`$
$`\mathrm{\Delta }_n^+H^{}(z;q^{(n)})`$ $`=`$ $`H^{}(zp^{c_{n+1}/2};q^{(n)})H^{}(zp^{c_n/2};q^{(n+1)}),`$
$`\mathrm{\Delta }_n^+E(z;q^{(n)})`$ $`=`$ $`E(z;q^{(n)})1H^{}(zp^{c_n/2};q^{(n)})E(zp^{c_n};q^{(n+1)}),`$
$`\mathrm{\Delta }_n^+F(z;q^{(n)})`$ $`=`$ $`1F(z;q^{(n+1)})+F(zp^{c_{n+1}};q^{(n)})H^+(zp^{c_{n+1}/2};q^{(n+1)}),`$
$`\mathrm{\Delta }_n^{}c_n`$ $`=`$ $`c_{n1}+c_n,`$
$`\mathrm{\Delta }_n^{}H^+(z;q^{(n)})`$ $`=`$ $`H^+(zp^{c_n/2};q^{(n1)})H^+(zp^{c_{n1}/2};q^{(n)}),`$
$`\mathrm{\Delta }_n^{}H^{}(z;q^{(n)})`$ $`=`$ $`H^{}(zp^{c_n/2};q^{(n1)})H^{}(zp^{c_{n1}/2};q^{(n)}),`$
$`\mathrm{\Delta }_n^{}E(z;q^{(n)})`$ $`=`$ $`E(z;q^{(n1)})1H^{}(zp^{c_{n1}/2};q^{(n1)})E(zp^{c_{n1}};q^{(n)}),`$
$`\mathrm{\Delta }_n^{}F(z;q^{(n)})`$ $`=`$ $`1F(z;q^{(n)})+F(zp^{c_n};q^{(n1)})H^+(zp^{c_n/2};q^{(n)});`$
* *the counits* $`ϵ_n`$*:*
$`ϵ_n(c_n)`$ $`=`$ $`0,`$
$`ϵ_n(1_n)`$ $`=`$ $`1,`$
$`ϵ_n(H^\pm (z;q^{(n)}))`$ $`=`$ $`1,`$
$`ϵ_n(E(z;q^{(n)}))`$ $`=`$ $`0,`$
$`ϵ_n(F(z;q^{(n)}))`$ $`=`$ $`0;`$
* *the antipodes* $`S_n^\pm `$*:*
$`S_n^\pm c_n`$ $`=`$ $`c_{n\pm 1},`$
$`S_n^\pm H^+(z;q^{(n)})`$ $`=`$ $`[H^+(z;q^{(n\pm 1)})]^1,`$
$`S_n^\pm H^{}(z;q^{(n)})`$ $`=`$ $`[H^{}(z;q^{(n\pm 1)})]^1,`$
$`S_n^\pm E(z;q^{(n)})`$ $`=`$ $`H^{}(zp^{c_{n\pm 1}/2};q^{(n\pm 1)})^1E(zp^{c_{n\pm 1}};q^{(n\pm 1)}),`$
$`S_n^\pm F(z;q^{(n)})`$ $`=`$ $`F(zp^{c_{n\pm 1}};q^{(n\pm 1)})H^+(zp^{c_{n\pm 1}/2};q^{(n\pm 1)})^1.`$
$`\mathrm{}`$
Let us stress that, among the defining relations of the superalgebra $`_{q,p}(osp(1|2)^{(1)})`$, the unusual signature in between the two $`\delta `$-function terms in the relation containing the anti-commutator of $`E(z)`$ and $`F(w)`$ is superficial: we can always replace $`H^{}(z)`$ by $`H^{}(z)`$ and e.g. $`F(z)`$ by $`F(z)(p^{1/2}+p^{1/2})/(pp^1)`$ and everything looks standard as in the usual q-affine algebra case.
It is remarkable that the comultiplication $`\mathrm{\Delta }_n^+`$ can be applied iteratively onto $`𝒜_n,`$ so that beginning from the $`c=1`$ realization one can obtain a realization of higher $`cZ_+.`$
If we re-parameterize the parameters $`q,`$ $`p`$ and $`z`$ as
$$q=e^{ϵ/\eta },p=e^ϵ\mathrm{},z=e^{iϵu}$$
and taking the scaling limit $`ϵ0,`$ the superalgebra $`_{q,p}(osp(1|2)^{(1)})`$ degenerates into
$`H^\pm (u)E(v)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi \eta (uv+2\mathrm{}\pm \mathrm{}c/2)\mathrm{sin}2\pi \eta (uv\mathrm{}\pm \mathrm{}c/2)}{\mathrm{sin}\pi \eta (uv2\mathrm{}\pm \mathrm{}c/2)\mathrm{sin}2\pi \eta (uv+\mathrm{}\pm \mathrm{}c/2)}}E(v)H^+(u),`$
$`H^\pm (u)F(v)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi \eta ^{}(uv+2\mathrm{}\mathrm{}c/2)\mathrm{sin}2\pi \eta ^{}(uv\mathrm{}\mathrm{}c/2)}{\mathrm{sin}2\pi \eta ^{}(uv2\mathrm{}\mathrm{}c/2)\mathrm{sin}2\pi \eta ^{}(uv+\mathrm{}\mathrm{}c/2)}}F(v)H^\pm (u),`$
$`H^\pm (u)H^\pm (v)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi \eta (uv+2\mathrm{})\mathrm{sin}2\pi \eta (uv\mathrm{})}{\mathrm{sin}2\pi \eta (uv2\mathrm{})\mathrm{sin}2\pi \eta (uv+\mathrm{})}}`$
$`\times {\displaystyle \frac{\mathrm{sin}2\pi \eta ^{}(uv2\mathrm{})\mathrm{sin}2\pi \eta ^{}(uv+\mathrm{})}{\mathrm{sin}2\pi \eta ^{}(uv+2\mathrm{})\mathrm{sin}2\pi \eta ^{}(uv\mathrm{})}}H^\pm (v)H^\pm (u),`$
$`H^+(u)H^{}(v)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi \eta (uv+2\mathrm{}+\mathrm{}c)\mathrm{sin}2\pi \eta (uv\mathrm{}+\mathrm{}c)}{\mathrm{sin}2\pi \eta (uv2\mathrm{}+\mathrm{}c)\mathrm{sin}2\pi \eta (uv+\mathrm{}+\mathrm{}c)}}`$
$`\times {\displaystyle \frac{\mathrm{sin}2\pi \eta ^{}(uv2\mathrm{}\mathrm{}c)\mathrm{sin}2\pi \eta ^{}(uv+\mathrm{}\mathrm{}c)}{\mathrm{sin}2\pi \eta ^{}(uv+2\mathrm{}\mathrm{}c)\mathrm{sin}2\pi \eta ^{}(uv\mathrm{}\mathrm{}c)}}H^{}(v)H^+(u),`$
$`E(z)E(w)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi \eta (uv+2\mathrm{})\mathrm{sin}2\pi \eta (uv\mathrm{})}{\mathrm{sin}2\pi \eta (uv2\mathrm{})\mathrm{sin}2\pi \eta (uv+\mathrm{})}}E(w)E(z),`$
$`F(z)F(w)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi \eta ^{}(uv+2\mathrm{})\mathrm{sin}2\pi \eta ^{}(uv\mathrm{})}{\mathrm{sin}2\pi \eta ^{}(uv2\mathrm{})\mathrm{sin}2\pi \eta ^{}(uv+\mathrm{})}}F(w)F(z),`$
$`\{E(u),F(v)\}`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}\left\{\delta \left(uv\mathrm{}c\right)H^+(v+\mathrm{}c/2)+\delta \left(uv+\mathrm{}c\right)H^{}(u+\mathrm{}c/2)\right\},`$
where
$$\frac{1}{\eta ^{}}\frac{1}{\eta }=\mathrm{}c.$$
This superalgebra is clearly an $`osp(1|2)^{(1)}`$ analogue of the earlier studied algebras $`𝒜_{\mathrm{},\eta }(\widehat{g})`$ and hence we call it $`𝒜_{\mathrm{},\eta }(osp(1|2)^{(1)}).`$ In the particular case of $`\eta 0`$ this superalgebra further degenerates into the super Yangian double $`DY_{\mathrm{}}(osp(1|2)^{(1)})`$ – the relations of which (first introduced in ) is just those of $`𝒜_{\mathrm{},\eta }(osp(1|2)^{(1)})`$ but with all the $`\mathrm{sin}2\pi \eta `$ and $`\mathrm{sin}2\pi \eta ^{}`$ removed – however with a sign difference appeared in the last relation. The reason for this sign difference has already been mentioned earlier in the context. This last degeneration fully clarifies the connection of our superalgebras $`_{q,p}(osp(1|2)^{(1)})`$ and $`𝒜_{\mathrm{},\eta }(osp(1|2)^{(1)})`$ with the underlying affine superalgebra $`osp(1|2)^{(1)}`$.
In closing, let us point out some related unsolved problems. The existence of two parameter deformation of affine Lie (super)algebras with the structure of infinite Hopf family of (super)algebras seems to be a universal phenomenon, which means that there should be such a (super)algebra associated with each underlying affine Lie (super)algebra. However, what we have known about these (super)algebras is only a tiny top of an iceberg. We know only a little about the structure theory and the representation theory. The definition of these algebras themselves were only known for untwisted affine Lie algebras associated with simply-laced Lie algebras and the present article add to this picture the simplest affine Lie superalgebra $`osp(1|2)^{(1)}`$. It seems that there remains a lot of pure algebraic works to do toward these (super)algebras.
On the other hand, physicists and/or applied mathematicians may be particularly interested in the application aspect of these (super)algebras. ¿From this point of view, we would like to mention the following problems which we would like to see a solution:
* Realizations of these (super)algebras other than the current realization.
It is well known that for $`q`$-affine algebras and Yangian doubles , there are mainly three different realizations which are proved to be connected to each other: the current realization, Drinfeld realization in terms of Laurent components of the currents and the Yang-Baxter realization . The last one is very important when physics applications are considered because it relates the structure of quantum symmetry algebra and the physical two-body $`S`$-matrix. So far we only know that, among the two parameter deformed affine Lie (super)algebras with the structure of infinite Hipf family of (super)algebras, only $`𝒜_{\mathrm{},\eta }(sl(2)^{(1)})`$ have a Yang-Baxter realization which contains a“dynamical operator” (spectral-shifting operator or weight vector of the underlying Lie algebra) . Attempts in obtaining Yang-Baxter realizations for all other current algebras of this kind have not lead to any success.
* The potential applications of these (super)algebras in physics problems. Besides being related to the algebra of screening currents of the quantum Virasoro and $`W`$-algebras, it would be interesting to see whether there is any physical model which bears any of these algebras as the underlying quantum symmetry. However, besides the case of $`𝒜_{\mathrm{},\eta }(sl(2)^{(1)})`$, nobody has ever been able to say a word on this possibility.
* Reconstruction of deformed $`W`$-algebras from the two parameter deformed affine algebras. To us, this seems the most plausible route to seek for physics applications, because these algebras are closely related to the algebra of screening currents of the deformed $`W`$-algebras, and it is indeed possible to reconstruct the $`W`$-algebras out of the screening currents. This problem is particularly interesting when the superalgebra associated with $`osp(1|2)^{(1)}`$ is considered because we then will be able to gain some knowledge about the deformed super Virasoro algebra – an object expected both for mathematical completeness and for physics applications!
Acknowledgement: L.Zhao would like to thank Niall MacKay for hospitality at Dept. Appl. Math., Sheffield University and at Dept. Math., Univ. of York during the preparation of this manuscript. The content of this article has been presented at Sheffield University in an informal seminar. This work is supported in part by the National Natural Science Foundation of China. |
warning/0003/cond-mat0003343.html | ar5iv | text | # Field induced ordering in highly frustrated antiferromagnets
\[
## Abstract
We predict that an external field can induce a spin order in highly frustrated classical Heisenberg magnets. We find analytically stabilization of collinear states by thermal fluctuations at a one-third of the saturation field for kagome and garnet lattices and at a half of the saturation field for pyrochlore and frustrated square lattices. This effect is studied numerically for the frustrated square-lattice antiferromagnet by Monte Carlo simulations for classical spins and by exact diagonalization for $`S=1/2`$. The field induced collinear states have a spin gap and produce magnetization plateaus.
\]
Frustrated magnets, classical and quantum, exhibit spectacular and often unexpected behaviors . A few of them do not order at any $`T>0`$. Extensive residual entropy in the classical ground state in such cases prevents the realization of the usual order by disorder scenario . An external magnetic field changes the degeneracy and topology of the ground state manifold stabilizing, eventually, the (nondegenerate) saturated state at $`H>H_{\mathrm{sat}}`$. If the order by disorder effect occurs in a finite field and suppresses residual entropy, then, such an effect can be used for practical applications. During a demagnetization process a spin system has to regain its entropy and, therefore, the whole crystal will cool down. Examples of geometrically frustrated AFMs on kagome , pyrochlore , and garnet lattices include magnetic compounds with rather small exchange constants opening the way for experimental tests of their finite field behavior.
In this Letter, we predict that fluctuations stabilize collinear spin configurations at rational values of the field $`H/H_{\mathrm{sat}}=1/2`$ or $`1/3`$. We present an analytical proof of the field induced ordering driven by thermal fluctuations and suggest a similar role for quantum fluctuations on the basis of spin-wave and numerical results. In fact, an ordered spin phase in a finite field has already been observed in gadolinium garnet Gd<sub>3</sub>Ga<sub>5</sub>O<sub>12</sub>, though dipolar anisotropy plays a crucial role in this material .
To be specific we use as an example the frustrated square-lattice antiferromagnet (FSAFM) in an external field:
$$\widehat{}=J\underset{\mathrm{n}.\mathrm{n}.}{}𝐒_i𝐒_j+J^{}\underset{\mathrm{n}.\mathrm{n}.\mathrm{n}.}{}𝐒_i𝐒_jH\underset{i}{}S_i^z.$$
(1)
First, we present analytical arguments, making them as general as possible, in order to include the other highly frustrated AFMs. Second, we consider numerical results for classical and quantum FSAFMs in a magnetic field, which confirm our predictions.
Let us briefly discuss magnetically ordered phases of the FSAFM. For small diagonal exchange $`J^{}<0.5J`$, classical spins form the Néel state in zero field. At $`J^{}>0.5J`$ the classical ground state consists of two $`\sqrt{2}\times \sqrt{2}`$ interpenetrating antiferromagnets, which are locked by fluctuations in a striped AFM state described by a single wave-vector $`(\pi ,0)`$ or $`(0,\pi )`$ . An applied magnetic field cants the spins and creates an easy plane for the AFM sublattices. Thus, both the Néel and striped states have a transverse spin order in a finite field. Thermal and quantum fluctuations destroy the transverse order in the vicinity of the highly frustrated point $`J^{}=0.5J`$. As an illustration we present in Fig. 1 the phase diagram for the spin-1/2 model at $`T=0`$ obtained in the linear spin-wave theory. This phase diagram agrees with the previous zero-field studies , which suggest a quantum disordered ground state for $`S=\frac{1}{2}`$ near $`J^{}=0.5J`$. Since the disordered singlet phase has a spin gap, it becomes unstable above a finite critical field.
We now focus on the classical critical point $`J^{}=0.5J`$, where the Hamiltonian (1) can be written up to a constant term as a sum over edge-sharing plaquettes:
$$\widehat{}=\frac{1}{2^n}\underset{\alpha }{}\left\{J\left|𝐋_\alpha \right|^2HL_\alpha ^z\right\},$$
(2)
where $`𝐋_\alpha =_{i\alpha }𝐒_i`$ is the total spin in plaquette $`\alpha `$ and $`n=2`$. The block form of the spin Hamiltonian (2) is common to all highly frustrated spin models. For AFMs on kagome and garnet lattices the blocks are triangles and for a pyrochlore AFM they are tetrahedra with the corner-sharing arrangements and $`n=1`$ in all three cases.
A spin configuration minimizes the energy (2) at $`H=0`$ provided $`𝐋_\alpha =0`$ for each plaquette. This constraint can be satisfied for many classical states. We estimate the number of continuous degrees of freedom in the ground state of the $`N`$-site FSAFM as $`N^{1/2}`$. For pyrochlore and kagome AFMs the degeneracy is larger and the dimensionality of the ground state manifold scales with $`N`$ . Thus, zero-field properties are not universal for disordered frustrated AFMs. We now show that a universal behavior does appear in a magnetic field.
In a finite field the classical energy (2) is minimized for spin configurations with the plaquette magnetization $`L_\alpha ^z=H/(2J)`$. There are many degenerate classical states which satisfy this constraint below the saturation field $`H_{\mathrm{sat}}=8JS`$ (pyrochlore and FSAFM) or $`6JS`$ (kagome and garnet). Generally, all these states are noncollinear as, e.g., canted Néel and striped phases for FSAFM. Collinear spin arrangements appear only at special rational values of the applied field: at $`H_c=\frac{1}{2}H_{\mathrm{sat}}`$ for 4-spin blocks (FSAFM, pyrochlore), the up-up-up-down (uuud) structure, see Fig. 2, and at $`H_c=\frac{1}{3}H_{\mathrm{sat}}`$ for 3-spin blocks (kagome, garnet), the up-up-down structure. The distribution of down-spins on the lattice in a collinear state is not unique: it only obeys the one down-spin per plaquette condition.
Zero-energy modes correspond to continuous distortions of a given classical ground state, which do not violate the magnetization constraint. To construct a zero mode for the collinear states described above one has to draw an (open) line through the lattice points with adjacent sites on a line occupied by antiparallel spins with no plaquettes crossed more than once. A zero mode consists of a simultaneous rotation of all spins along the line by an angle $`\theta `$ (Fig. 2). The complete set of zero modes is constructed when a set of such lines goes through all up-spins. Zero modes connect a uuud state to noncollinear ground state configurations. Hence, the collinear states are singular points on the ground state manifold. There are no other collinear states in the local neighborhood of a selected collinear state.
Moessner and Chalker have recently discussed the order by disorder phenomenon in zero field for classical frustrated models in terms of a local topology of the ground state manifold. They reached definite conclusions only for an $`XY`$ magnet on the pyrochlore lattice and for a Heisenberg kagome lattice AFM. We show that a similar analysis predicts stabilization of collinear states near special rational values of the applied field. Let $`𝐱`$ denote the coordinates on the ground state manifold and $`𝐲`$ be the transverse directions, which span the rest of the configuration space. Low-energy excited states are described by the quadratic Hamiltonian $`_2=_lϵ_l(𝐱)y_l^2`$ resulting in a probability distribution
$$Z(𝐱)=𝑑y_le^{\beta _2}\underset{l}{}[T/ϵ_l(𝐱)]^{1/2}$$
(3)
over the ground state manifold at low temperatures. At a special point $`𝐱_0`$ some of the stiffnesses $`ϵ_l(𝐱_0)`$ may vanish making $`Z(𝐱_0)`$ divergent. The corresponding coordinates $`y_l`$ describe soft modes. The appearance of the order by disorder depends on a number of soft modes, which exist for a given classical ground state $`𝐱_0`$.
A special feature of the collinear states with finite magnetization is that each zero mode generates exactly one soft mode. They are constructed in the following way. First, all spins along a zero mode line are rotated by $`\theta `$ about the axis perpendicular to the field and, second, all down spins on the same line are rotated by an angle $`\phi `$ about the field direction, see Fig. 2. The total spin of each plaquette in a deformed state is $`𝐋_\alpha S(0,\theta \phi ,2)`$ for small $`\theta `$ and $`\phi `$ and contributes to an energy increase $`\theta ^2\phi ^2`$. Thus, we identify $`x_l`$ with $`\theta `$ and $`y_l`$ with $`\phi `$.
To allow the order by disorder selection, not only has the probability density $`Z(𝐱_0)`$ to diverge, but the statistical weight $`Z(𝐱)𝑑𝐱`$ must be concentrated entirely near the collinear spin states. We check this by integrating $`Z(𝐱)`$ in the $`D`$-dimensional neighborhood of a collinear state parameterized by zero modes. Since each zero mode has one soft mode with $`ϵ_s(𝐱)x^2`$, the integral
$$Z(𝐱)𝑑𝐱x^Dd^Dx$$
(4)
diverges independently of the actual value of $`D`$. There is a vanishingly small probability of finding a spin system in one of the noncollinear states surrounding a given collinear configuration. Hence, the order by disorder selection occurs: thermal fluctuations stabilize a discrete set of collinear states. This result relies only on the special block structure (2) of the spin Hamiltonian and the symmetry properties of the collinear states. Therefore, the order by disorder in an external field is a universal effect and occurs for all frustrated Heisenberg AFMs: apart from the FSAFM the uuud states at $`H=\frac{1}{2}H_{\mathrm{sat}}`$ are stabilized for a pyrochlore AFM and the uud states appear for Heisenberg AFMs on kagome and garnet lattices at $`H=\frac{1}{3}H_{\mathrm{sat}}`$. A subsequent selection between collinear states with different patterns of down-spins is made to maximize $`D`$ in Eq. (4) and corresponds in the case of FSAFM to the most symmetric $`𝐪=0`$ uuud state. The above proof of the field induced ordering is based on the special property of the collinear states with finite magnetization. Namely, an equal number of zero and soft modes. For a Heisenberg magnet on the pyrochlore lattice in zero field the same equality holds only approximately: in the leading order in the number of spins $`N`$. Therefore, no analytical conclusion has been reached in a zero field case, while numerical simulations indicated no ordering .
For quantum models in the limit $`S1`$ we check the relative stability of different states by comparing their zero-point oscillation energies: $`\frac{1}{2}_𝐤\omega _𝐤`$. The magnon spectrum of the $`𝐪=0`$ uuud state has four modes:
$$\omega _{1,2}=JS\sqrt{\eta _𝐤},\omega _{3,4}=2JS\pm JS\sqrt{4\eta _𝐤}$$
(5)
with $`\eta _𝐤=(1\mathrm{cos}k_x)(1\mathrm{cos}k_y)`$. Zero-point contributions to the energies of the $`𝐪=0`$ uuud state, the Néel state, and the striped AFM at $`J^{}=0.5J`$ and $`H=\frac{1}{2}H_{\mathrm{sat}}`$ are $`0.57JS`$, $`0.69JS`$, and $`0.81JS`$, respectively. (They have the same classical energy.) Thus, this comparison suggests that quantum fluctuations also select the collinear states in a magnetic field due to their large number of soft modes. Collinear states preserve the $`O(2)`$-rotational symmetry. Hence, their renormalized magnon spectrum becomes gapped and plateaus arise on the magnetization curve. The spin-density wave of the $`𝐪=0`$ uuud state in FSAFM is a superposition of spin harmonics with wave-vectors $`(0,0)`$, $`(\pi ,\pi )`$, $`(\pi ,0)`$ and $`(0,\pi )`$.
To test these predictions we have performed Monte Carlo (MC) simulations for the model Eq. (1) with unit vector spins. Fig. 3, top panel shows the magnetization curve at $`T=0.1J`$ obtained with $`2\times 10^5`$ MC steps per spin for each point. The field derivative of the magnetization has a dip around $`H=4J=\frac{1}{2}H_{\mathrm{sat}}`$, which indicates the presence of a new plateau phase. A jump in the magnetization and a hysteresis of about $`\mathrm{\Delta }H0.3J`$ suggests a first-order transition to a high-field phase. To determine the nature of the spin state at the plateau we calculated different components of the static structure factor $`S^{\alpha \beta }(q)=\frac{1}{N^2}_{r,x}\mathrm{e}^{i𝐪𝐱}S_r^\alpha S_{r+x}^\beta `$. The Néel and the collinear states have nonzero transverse components $`S^{xx}=S^{yy}`$ at $`𝐪=(\pi ,\pi )`$ and $`𝐪=(\pi ,0)`$ or $`(0,\pi )`$, respectively. The uuud state has nonzero longitudinal components $`S^{zz}(q)`$ at all the above vectors simultaneously. The static structure factor presented in Fig. 3, bottom panel was obtained by averaging 50 ‘instant shots’ separated by $`10^3`$ MC steps. In the region of weak (strong) diagonal exchange $`J^{}`$ the data clearly support the Néel (striped) type of spin correlations. Nonzero harmonics in the longitudinal structure factor both at $`𝐪=(\pi ,0)`$ and $`𝐪=(\pi ,\pi )`$ exist only for $`0.494<J^{}/J<0.508`$. We have also checked that neither a lower field of $`H=2J`$, nor a higher field of $`H=6J`$ induces a nonzero value of $`S^{zz}(q)`$ at these points in the Brillouin zone. Thus, these results unambiguously identify the spin configuration on the plateau as the $`𝐪=0`$ uuud state.
The uuud phase breaks the translational symmetry in such a way that the wave-vectors belonging to different irreducible representations of the space group are mixed: $`(\pi ,\pi )`$ vs. $`(\pi ,0)`$ and $`(0,\pi )`$. Hence, the spin structure of the uuud phase is a mixture of different order parameters. In such a case the phase transition to a disordered spin liquid state at the high-field end of the plateau, Fig. 3, must go either via a first order transition or in several steps with an intermediate supersolid state.
We have also studied the quantum spin-1/2 model (1) at $`T=0`$ by Lanczos diagonalizations of finite clusters. Fig. 4 presents the magnetization $`m=M/S`$ vs. field at $`J^{}/J=0.6`$. There are no magnetization plateaus in the thermodynamic limit for $`m>1/2`$. To determine whether the plateau at $`m=\frac{1}{2}`$ remains after finite size scaling we show its width in the inset for three cluster sizes as a function of $`J^{}/J`$. If a plateau disappears in the thermodynamic limit, its width should decrease as $`\mathrm{\Delta }m=1/N`$. Therefore, the ratio of the plateau-widths of the $`4\times 4`$ and $`6\times 6`$ clusters should be $`16/36`$ or less if the magnetization curve has a nonzero slope for $`N\mathrm{}`$. The ratio $`16/36`$ is exceeded for $`m=1/2`$ only in the interval $`0.49J^{}/J0.66`$, where we would expect to have a plateau on an infinite lattice.
Quantum fluctuations have a strong effect on the stability of the $`m=\frac{1}{2}`$ plateau. Both, the width of the plateau for the spin-1/2 system and the range where it appears are larger than for the classical model. In addition, the parameter range for the plateau is shifted asymmetrically around the classical critical point $`J^{}=0.5J`$. The plateau is most pronounced for $`J^{}0.6J`$. This is the value we used for the presented magnetization curve. The bold line in Fig. 4 was obtained by connecting the mid-points of the steps for the largest available systems size, except for the plateau at $`m=\frac{1}{2}`$, for which the corners of the $`6\times 6`$ cluster data were used. The value of the spin gap (i.e. the boundary of the $`m=0`$ plateau) was taken from Kotov et al. Fig. 4 also shows the peaks of the static structure factor vs. $`J^{}/J`$ for the $`m=\frac{1}{2}`$ plateau. The peaks in $`S^{zz}(q)`$, which indicate the presence of the uuud state, exist for $`0.51J^{}/J0.67`$.
The discussed mechanism for the magnetization plateaus in highly frustrated AFMs has to be contrasted with that of weakly coupled 1D spin systems (see, e.g., ). There, magnetization plateaus correspond to disordered states with a gap, that are stable under a weak higher-dimensional coupling, while intermediate gapless regions are immediately ordered at sufficiently low temperatures once a higher-dimensional coupling is switched on. Here, we found that an external field induces a long-range collinear ordering with an excitation gap on the plateau, whereas intermediate gapless regions remain disordered. Further exploration of such possibilities should also be important for the general problem of classifying magnetization plateaus in two dimensions .
We thank J.T. Chalker, T.M. Rice, and O.P. Sushkov for helpful discussions. A.H. is grateful to the Institut für Theoretische Physik of the ETH Zürich for hospitality during the course of this project, to the Alexander von Humboldt-foundation for financial support and to the Max-Planck-Institut für Mathematik, Bonn for allocation of CPU time. The work of M.E.Z. was supported by the Swiss National Fund. |
warning/0003/hep-ph0003093.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Investigation of the multiplicity distributions was popular since seventies . The modern hadron theory based on the local QCD Lagrangians and the experimental consequences was given in the review papers . The very high multiplicity (VHM) processes, as the attempt to get beyond this standard multiperipheral hadron physics, was offered in . It considered as a possible physical program for LHC experiments.
It was shown in seventieth that the multiperipheral kinematics dominates inclusive cross sections. Moreover, the created particles spectra do not depend on $`s`$ at high energies in the multiperipheral region:
$$f(s,p_c)=2E_c\frac{d\sigma }{d^3p_c}=\frac{dt_1dt_2s_1s_2\varphi _1(t_1)\varphi _2(t_2)}{(2\pi )^2s(t_1m^2)^2(t_2m^2)^2},s_1s_2=sE_c^2,E_c^2=m_c^2+\stackrel{}{p}_c^2.$$
Here $`s_1=(p_a+p_c)^2,s_2=(p_b+p_c)^2`$, and $`\varphi _i(t_i)`$ are the impact factors of hadrons. So the particle $`c`$ forgot the details of its creation. It was found experimentally that the ratio
$$\frac{f(\pi ^+p\pi _{}+\mathrm{})}{\sigma (\pi ^+p)}=\frac{f(K^+p\pi _{}+\mathrm{})}{\sigma (K_+p)}=\frac{f(pp\pi _{}+\mathrm{})}{\sigma (pp)}$$
(1.1)
is universal . This take place due to the two Pomeron multiperipheral exchange providing the nonvanishing contribution in the $`s`$ asymptotics to the cross section. It was implied that the Pomeron intercept is exactly equal to one. Just this kinematics leads to $`c_m=\gamma _m(c_1)^m`$ (the correlators $`c_m`$ are introduced in (1.8)), i.e. to the KNO-scaling .
We begin with general analysis to formulate an aim of this paper. Let $`\sigma _n(s)`$ be the cross section of $`n`$ particles (hadrons) creation at the total CM energy $`\sqrt{s}`$. We introduce the generating function:
$$T(s,z)=\underset{n=1}{\overset{n_{max}}{}}z^n\sigma _n(s),s=(p_1+p_2)^2>>m^2,n_{max}=\sqrt{s}/m.$$
(1.2)
So, the total cross section and the averaged multiplicity will be:
$$\sigma _{tot}=T(s,1)=\sigma _n,\sigma _{tot}\overline{n}=n\sigma _n=\frac{d}{dz}T(s,z)|_{z=1}.$$
(1.3)
At the same time
$$\sigma _n=\frac{1}{2\pi i}\frac{dz}{z^{n+1}}T(s,z)=\frac{1}{2\pi i}\frac{dz}{z}e^{(n\mathrm{ln}z+\mathrm{ln}T(s,z))}.$$
(1.4)
The essential values of $`z`$ in this integral are defined by the equation (of state):
$$n=z\frac{}{z}\mathrm{ln}T(z,s).$$
(1.5)
Considering the tail, i.e. $`n>>\overline{n}`$, let us assume that one can find such values of $`n<<n_s`$ at high energies $`\sqrt{s}>>m`$ that we can neglect in (1.2) dependence on the upper boundary $`n_{max}`$. This formal trick allows to consider $`T(z,s)`$ as the nontrivial function of $`z`$. Then the asymptotics over n ($`n<<n_s`$ is assumed) is governed by the mostleft situated singularity $`z_s`$ of $`T(z,s)`$:
$$\sigma _n(s)e^{n\mathrm{ln}z_c(n,s)},$$
(1.6)
where $`z_c(n,s)`$ is the smallest solution of eq.(1.5). It is important that
$$z_s(n,s)z_c\mathrm{at}n\mathrm{}.$$
(1.7)
One can put this method of asymptotic estimation in the basis of VHM processes phenomenology.
We may distinguish following possibilities at $`n\mathrm{}`$:
1)$`z_s=1`$: $`\sigma _n>O(e^n)`$;
2)$`z_s=\mathrm{}`$: $`\sigma _n<O(e^n)`$;
3)$`i<z_s<\mathrm{}`$: $`\sigma _n=O(e^n)`$.
The asymptotics 1) assumes the condensation phenomena . The asymptotics 2) belong to the multiperipheral processes kinematics: created particles form jets moving in the CM frame with different velocities along the incoming particles directions, i.e. with restricted transverse momentum. The third type asymptotics is predicted by stationary Markovian processes with the QCD jets kinematics of the high transverse momentum particles creation. It is evident that the $`n`$ asymptotics should be governed by largest among 1) - 3). Just under this idea we hope that we get beyond the multiperipheral kinematics in the VHM region.
Additional information is coded in the expansion over binomial moments $`c_m(s)`$:
$$\mathrm{ln}T(s,z)=\frac{(z1)^m}{m!}c_m(s).$$
(1.8)
So, for instance, if all $`c_m=0,m>1`$, then we have the Poisson distribution:
$$\sigma _n=\sigma _{tot}\frac{(\overline{n})^n}{n!}e^{(\overline{n})}.$$
But if $`c_m=\gamma _m(c_1)^m`$, where $`\gamma _m`$ is the some constant, then the so called KNO scaling take place:
$$\overline{n}\sigma _n=\sigma _{tot}\mathrm{\Psi }(n/\overline{n}).$$
The multiple production processes are typical inelastic reactions of the initial kinetic energy dissipation into the particles mass. Consequently, the mean multiplicity $`\overline{n}(s)`$ is the measure of entropy $`𝒮`$ production at given energy. Experimentally $`\overline{n}(s)\mathrm{ln}^2s<<n_{max}`$. This testify to the incomplete energy dissipation in the mostly probable channels of hadrons production.
This phenomena is explained naturally by presence of the space-time local non-Abelian symmetry constraints. So, it is known that there is not thermalization phenomena in the completely integrable systems. But, at all evidence, the quantum Yang-Mills theory is not completely integrable, i.e. it admits the dissipation, but nevertheless the symmetry constraints play essential role.
It is natural to assume that $`𝒮`$ exceed its maximum if $`n>>\overline{n}(s)`$. So, our essentially inelastic process is happened so rapidly that the non-Abelian symmetry constraints becomes frozen. On other hand, maximum of entropy $`𝒮`$ means that the final state of the dissipation process is equilibrium.
Last one means relaxation of energy correlations, i.e. absence of the macroscopical energy flows in the system, and the Gauss energy spectrum of created particles. We would like to say that in such a state one get to the VHM.
The aim of this article is to build the suitable mechanism of maximal initial energy dissipation, where the correlations are relaxed and the energy spectra are Gaussian.. The classification 1)-3) is mostly general and we will put it in the basis of consideration. So, our main purpose is to show as the multiperipheral kinematics transform in the VHM domain.
## 2 Pomeron,DIS and Double-Logarithmic kinematics
Let us consider process of type $`22+n`$ in different kinematics
$$A(P_1)+B(P_2)A^{}(p_1^{})+B^{}(p_2^{})+h_1(k_1)+h_2(k_2)+\mathrm{}+h_n(k_n),s=(p_1+p_2)^2>>m^2.$$
We will distinguish the peripheral,deep inelastic and large-angles kinematical regions with different physical content. First we build two light-like 4-momenta from the momenta of initial particles $`p_{1,2}=P_{1,2}P_{2,1}m_{2,1}^2/s`$ and present the 4-momenta of final particles in form:
$$p_1^{}=\alpha _1^{}p_2+\beta _1^{}p_1+p_1^{};a_{}p_{1,2}=0;p_2^{}=\alpha _2^{}p_2+\beta _2^{}p_1+p_2^{};k_i=\alpha _ip_2+\beta _ip_1+k_i^{}.$$
(2.1)
Sudakov’s parameters $`\alpha ,\beta `$ are not independent.The mass shell conditions and the conservation low give the relations:
$`s\alpha _1^{}\beta _1^{}=m_1^2+\stackrel{}{p}_1^{}_{}{}^{}2=E_1^2,`$ (2.2)
$`s\alpha _2^{}\beta _2^{}=E_2^2,s\alpha _i\beta _i=E_i^2,`$
$`\stackrel{}{a}^2=a_{}^2>0;\alpha _1^{}+\alpha _2^{}+{\displaystyle \alpha _i}=1;`$
$`\beta _1^{}+\beta _2^{}+{\displaystyle \beta _i}=1.`$
Consider first the peripheral kinematics. For it is characteristic the weak dependence of differential cross sections on the center of mass (CM) total energy $`2E=\sqrt{s}`$,the strict ordering of parameters $`\alpha ,\beta `$
$$1\beta _1^{}>>\beta _1>>\mathrm{}>>\beta _n>>\beta _2^{}\frac{m^2}{s};\frac{m^2}{s}<<\alpha _1^{}<<\alpha _1<<\mathrm{}<<\alpha _n<<\alpha _2^{}1$$
and the restrictiveness on the transvers momenta $`|\stackrel{}{k}_i|m`$. It corresponds to small emission angles moving along 3-momentum $`\stackrel{}{P}_1`$
$$\theta _i=\frac{|\stackrel{}{k}_i|}{E\beta _i}<<1,|\beta _i|>>|\alpha _i|,$$
and the similar expression for particles moving in opposite direction $`|\beta _i|<<|\alpha _i|`$. The central region $`|\alpha _i||\beta _i|E_i/E<<1`$ corresponds to particles of low energies moving at large angles(in cm frame). The differential cross section have the form:
$`d\sigma _{22+n}`$ $`=`$ $`{\displaystyle \frac{(2\alpha _s)^{2+n}}{16\pi ^{2n}}}C_V^n{\displaystyle \frac{d^2q_1}{q_1^2+m^2}}{\displaystyle \frac{d^2q_2}{(q_1q_2)^2+\lambda ^2}}\mathrm{}`$
$`\times `$ $`{\displaystyle \frac{d^2q_{n+1}}{(q_nq_{n+1})^2+\lambda ^2}}{\displaystyle \frac{1}{q_{n+1}^2+\lambda ^2}}{\displaystyle \frac{d\alpha _1}{\alpha _1}}\theta (\alpha _2\alpha _1)\mathrm{}`$
$`\times `$ $`{\displaystyle \frac{d\alpha _n}{\alpha _n}}{\displaystyle \underset{i=1}{\overset{n+1}{}}}({\displaystyle \frac{s_i}{m^2}})^{2\alpha (q_i^2)}={\displaystyle \frac{1}{q_{n+1}^2+\lambda ^2}}dZ_n,`$
where $`C_V=3`$ and we imply $`q_i`$ the two-dimensional euclidean vectors,the 4-momentum of the $`i`$-th emitted particle (gluon)is
$$k_i=(\alpha _i\alpha _{i+1})p_2+(\beta _i\beta _{i+1})p_1+(q_iq_{i+1})_{}=\alpha _{i+1}p_2+\beta _ip_1+(q_iq_{i+1})_{};$$
(2.4)
$`s_i`$ are the partial squares of invariant mass of nearest emitted particles:
$$s_1=(p_1^{}+k_1)^2=s|\alpha _2|,s_{n+1}=(k_n+p_2^{})^2=\frac{E_n^2}{\alpha _n},s_i=E_{,i1}^2\frac{\alpha _{i+1}}{\alpha _{i1}},s_1s_2\mathrm{}s_n=sE_1^2\mathrm{}E_n^2,$$
(2.5)
and the trajectory of reggeized gluon is
$$\alpha (q^2)=\frac{q^2\alpha _s}{2\pi ^2}\frac{d^2k}{(k^2+\lambda ^2)((qk)^2+\lambda ^2}.$$
Here $`\lambda `$ is gluon mass.It was shown that the infrared singularities absent in the limit $`\lambda 0`$.In this point we will suggest the gluon to be massive: $`\lambda M`$ and will decay to the jet of hadrons (pions) with the probability to create $`n`$ particles $`dW_n(M)=\frac{c}{\overline{n}}exp(n/\overline{n})dn,\overline{n}=\mathrm{ln}(M^2/m_\pi ^2)`$.
The Monte Carlo simulation shows a tendency to minimization of the number of rungs at large $`n`$.
For the pure deep inelastic case, when one of the initial hadrons is scattered at the angle $`\theta `$ have the energy $`E^{}`$ in the cms of beams whereas the another is scattered at small angle and the large transfer momentum $`Q=4EE^{}\mathrm{sin}^2(\theta /2)>>m^2,`$ is distributed to the some number of the emitted particles due to evolution mechanism we have ($`\theta `$ is small):
$`d\sigma _n^{DIS}={\displaystyle \frac{4\alpha ^2E^{}_{}{}^{}2}{Q^4M}}dD_ndE^{}d\mathrm{cos}\theta ,`$
$`dD_n=({\displaystyle \frac{\alpha _s}{4\pi }})^n{\displaystyle _{m^2}^{Q^2}}{\displaystyle \frac{dk_n^2}{k_n^2}}{\displaystyle _{m^2}^{k_n^2}}{\displaystyle \frac{dk_{n1}^2}{k_{n1}^2}}\mathrm{}{\displaystyle _{m^2}^{k_2}}{\displaystyle \frac{dk_1^2}{k_1^2}}{\displaystyle _x^1}𝑑\beta _n\mathrm{\Theta }_n^{(1)}{\displaystyle _{\beta _n}^1}𝑑\beta _{n1}\mathrm{\Theta }_{n1}^{(1)}\mathrm{}`$
$`\times {\displaystyle _{\beta _2}^1}d\beta _1\mathrm{\Theta }_1^{(1)}P({\displaystyle \frac{\beta _n}{\beta _{n1}}})\mathrm{}P(\beta _1),P(z)=2{\displaystyle \frac{1+z^2}{1z}},`$ (2.6)
where the limits of integrals show the intervals of variation and the integrand is the differential cross section. Again the rapidities $`\beta _i`$ are arranged and the transverse momenta square are rigorously arranged. Here $`\mathrm{\Theta }^{(i)}=\theta (\theta _{i+1}\theta _i)`$ the condition which forbids the destructive interference of jets and jets emission angles are $`\theta _i=|\stackrel{}{k}_i|/(E\mathrm{max}|\alpha _i|,|\beta _i|)`$ regarding the beam axe direction.
Compared with peripheral production regime DIS one gives the contribution which fall down with increasing $`Q^2`$. Using the analogy to statistical nonequilibrium processes we may consider DIS regime as an nonequilibrium process of diminishing of large virtualities to small ones by means of evolution. The peripheral regime may be associated with equilibrium process when the random fluctuations on the transvers momenta values become essential.
Let consider now the regime of hard particles production at large angles. It is known as a double-logarithm regime. Any exclusive process is suppressed in this regime by Sudakov form factor. For inclusive set-up of experiments when arbitrary number of photons (gluons) may be emitted the Sudakov’s form factor suppression disappears and we obtain the cross section of the form :
$$d\sigma (s)=\frac{1}{s}F(\alpha \mathrm{ln}^2\frac{s}{m^2}).$$
(2.7)
For instance the cross section of annihilation of electron-positron to muon pair accompanied by emission of arbitrary number of photons is
$$\sigma (s)_{e\overline{e}\mu \overline{\mu }+\mathrm{}}=\frac{4\pi \alpha ^2}{3s}chx,x^2=\frac{2\alpha }{\pi }\mathrm{ln}^2(\frac{s}{m_\mu m_e}).$$
The characteristic squares of transvers momenta of created particles in this regime are big and of order of $`s`$.
The typical process -annihilation of electron -positron pair to $`n`$ hard photons accompanied by emission of any number of soft and virtual photons have a form :
$`d\sigma _n={\displaystyle \frac{2\pi \alpha ^2}{s}}dF_n^{DL},`$ (2.8)
$`dF_n^{DL}=({\displaystyle \frac{\alpha }{2\pi }})^n{\displaystyle \underset{i=1}{\overset{i=n}{}}}dx_idy_i\theta (y_iy_{i1})\theta (x_iy_i)\theta (x_i)\theta (y_i)\theta (\rho x_n)\theta (\rho x_n),`$
$`x_i=\mathrm{ln}{\displaystyle \frac{\stackrel{}{q}_i^2}{m^2}},y_i=\mathrm{ln}{\displaystyle \frac{1}{\beta _i}},\rho =\mathrm{ln}{\displaystyle \frac{s}{m^2}}.`$
## 3 Various scenario
Keeping in mind the kinematical restriction $`\mathrm{\Pi }s_js`$ we may build the combinations of regimes considered above. Le construct the relevant cross sections. It is convenient to separate them to the classes
a) Pomeron regime (P);
b) Evolution regime (DIS);
c) Double logarithmic regime (DL);
d) DIS+P regime;
e) P+DL+P regime.
The description of every regime may be performed in terms of effective ladder-type Feynman diagrams (The set of relevant FD depends on the gauge chosen and include much more number of them).
The Pomeron is treated as a (infinite) set of particles emitted close to the CM beams direction (within the small angles of order $`\theta _i2m_h/\sqrt{s}<<1`$). We expect that these type of particles will not be detected by the detectors since they are move into the beams pipe. The collider experiment detectors locate at finite angles $`\theta _D1`$ and will measure the products only of particle $`c`$ decay.
What will happened when instead of one particle (see (1.1)) a set of particles with invariant mass square $`s_t`$ is created at large angles? Then the cross section will have a form:
$`d\sigma _n={\displaystyle \frac{\alpha _s^2}{s_t}}NdF_n^{DL}(\alpha _s\mathrm{ln}^2({\displaystyle \frac{s_t}{s_0}})),N=({\displaystyle \frac{s}{s_t}})^\mathrm{\Delta },`$
$`\mathrm{\Delta }=\alpha _P1={\displaystyle \frac{12\mathrm{ln}2}{\pi }}\alpha _s0.55,\alpha _s=0.2`$ (3.1)
Radiative corrections to the intercept was calculated in recent time. The resulting value is $`\mathrm{\Delta }0.2`$.
The way to obtain detected large multiplicity is to organize DIS-like experiments, expecting the large-angle scattered hadrons in the detectors. Large transfers momenta will be decreasing by ordinary evolution mechanism to the value of order $`m_\pi `$ and then the Pomeron mechanism of peripheral scattering of the created hadrons from the pionization region will start.
This phenomena is quite close to flea-dog model of Euhrenfest (see ): in the nonequilibrium process (DIS regime) the fluctuations are suppressed and they take place (Pomeron regime) when the equilibrium take place.
What the characteristic multiplicities expected from Pomeron mechanism with the intercept exceeding unity, $`\mathrm{\Delta }0.2`$? It is the quantity of order $`(s/m_\pi ^2)^\mathrm{\Delta }200`$ for $`\sqrt{s}=14TeV`$. This rather rough estimation is in agreement with the phenomenological analysis of A.Kaidalov , based on multi-pomeron exchange in the scattering channel.
Construct now the cross sections of combined processes. When the one of the initial particles $`h_1`$ is scattered on small but sufficient enough angle to fit the detectors and other is scattered almost forward the combination of DIS and Pomeron regimes take place:
$$d\sigma _{n,m}=d\sigma _n^{DIS}dZ_m,|q_n|^2m^2$$
(3.2)
provided that the virtuality of the last step of evolution regime of order of hadron mass. For the kinematical case of almost forward scattering of both initial hadrons the situation may be realized with large angles hadron production from the central region (see (3.1)):
$$d\sigma _{n,m,k}=dZ_nd\sigma _m^{DL}dZ_k.$$
(3.3)
## 4 Discussion
Every regime $`a)e)`$ considered above provide creation of $`(n+2)`$ gluons and theirs subsequent decay in the universal way on jets.
Let us discuss more carefully the large-angle moving jets creation mechanism due to annihilation of initial partons into the hadrons system. The intermediate state of single heavy photons decay into jet was analyzed by A.Polyakov . It was shown that the scaling regime works, providing the behavior of the $`n`$ jets creation cross section and the mean multiplicity as:
$$\sigma _n(s)\frac{1}{s}(\frac{m^2}{s})^\delta F(n(\frac{m^2}{s})^\delta ),\overline{n}(\frac{s}{m^2})^\delta ,0<\delta <\frac{1}{2}.$$
(4.4)
Another mechanism of multiple production takes into account the channels of incident partons annihilation onto the arbitrary large number of real and virtual gluons (photons). On this way some intermediate regime between the single logarithm (the renormalization group approach) and the double logarithmical, $`\alpha \mathrm{ln}^\rho (s/m^2)1,1<\rho <2`$, is realized.
All the considered regimes of many particles state production describes the hard stage of process. The hadronization stage will impose its features which can not be expressed in terms of pQCD as well as it concern the confinement region. Here the identity of produced particles must be taken into account . The effective Lagrangian approach also may be applied here . We hope to consider last questions in another work.
Acknowledgments
We are grateful to V.G.Kadyshevski for interest to discussed in the paper questions. |
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